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warning/0004/cond-mat0004343.html	ar5iv	text	# Dynamical response of a pinned two-dimensional Wigner crystal ## I Introduction The problem of pinning of an elastic solid by an external random potential has been a long-standing problem of condensed matter physics, first addressed by Larkin in 1970. Physical systems where this problem arises include charge-density waves, vortex lattices in superconductors, interfaces, magnetic bubbles, liquid crystals in aerogels, and many others. An interesting example of the charge-density wave is a Wigner crystal (WC). Under terrestrial conditions a stable WC has been realized in its two-dimensional (2D) form, on cryogenic surfaces and in semiconductor heterostructures with low carrier density and/or in strong magnetic fields. In this work we examine the finite-frequency response of a collectively pinned 2D WC. The collective pinning is the regime where individual defects of the substrate are too weak to significantly deform the crystalline lattice, so that the WC is well ordered (“unpinned”) at small length scales, but the cumulative effect of the disorder eventually dominates the elasticity at a pinning length, which is much larger than the lattice constant. The linear-response conductivity $`\sigma _{\alpha \beta }(q,\omega )`$ of the WC in an external magnetic field is a tensor. We focus our study on the real part of its diagonal component $`\sigma _{xx}(q=0,\omega )`$, the quantity that determines the power absorption in the presence of a uniform electric field. If the frequency $`\omega `$ is not too small, the absorption is dominated by the collective excitations. In the absence of the magnetic field and random potential, a 2D WC possesses two branches of such excitations: the longitudinal (L) and the transverse (T) phonons. The long-wavelength L-phonon is also referred to as plasmon. In a finite magnetic field L- and T-phonons hybridize into the magnetophonons (gapless lower hybrid mode) and the magnetoplasmons (gapped higher hybrid mode). It is noteworthy that without the pinning potential, the frequency and the oscillator strength of the magnetophonons vanish at $`q=0`$ by virtue of Kohn’s theorem. The absorption of a spacially uniform ac field is possible only by exciting the magnetoplasmon mode, at the cyclotron frequency $`\omega _c=eB/m_ec`$ set by the magnetic field $`B`$. The pinning shifts the original $`q=0`$ magnetophonon mode to a nonzero disorder-dependent frequency $`\omega _p`$, broadens it, and imparts it with a finite oscillator strength. The resulting absorption line, which can be called the pinning mode, was first discussed by Fukuyama and Lee (FL). In those early works the linewidth $`\mathrm{\Delta }\omega _p`$ of the pinning mode was predicted to be of the order of $`\omega _p`$ itself. Strictly speaking, FL considered not a WC but a more conventional charge-density wave where the charge modulation has only one harmonic, i.e., is cosine-like. Experiments on such materials, e.g., $`\mathrm{K}_{0.3}\mathrm{Mo0}_3`$ (“blue bronze”), $`\mathrm{TaS}_3`$, etc., performed in zero magnetic field indeed revealed broad absorption lines at disorder-dependent frequencies. A later work of Normand et al. devoted specifically to the WC in a strong $`B`$ revised some results of FL but left unchanged the prediction that $`\mathrm{\Delta }\omega _p\omega _p`$. At the time, experiments seemed to confirm that. It came as a surprise when most recent measurements in strong magnetic fields demonstrated that $`\mathrm{\Delta }\omega _p`$ can be more than order of magnitude smaller than $`\omega _p`$. Such unexpected findings revived interest in this long-standing problem. An important step towards resolution of the puzzle has been made by Fertig, who showed that long-range Coulomb interaction plays an important role in reducing the inhomogeneous broadening of the pinning mode. However, Fertig used an oversimplified model (commensurate pinning) and did not calculate $`\mathrm{\Delta }\omega _p`$ directly. His results for another quantity may be interpreted as an indication that $`\mathrm{\Delta }\omega _p/\omega _p`$ is exponentially small for weak pinning. More realistic model was studied by Chitra et al. who considered pinning by a short-range Gaussian random potential of a general type. We comment on their results later in this section. In the present paper we study essentially the same model as FL and Chitra et al. except we treat the electrons classically. This limits the applicability of our results to the case $`\xi >l_B`$, where $`\xi `$ is the correlation length of the pinning potential and $`l_B=\sqrt{\mathrm{}c/eB}`$ is the magnetic length. The classical approximation enables us to focus on the interaction of the WC with disorder, which is really the essence of the pinning mode phenomenon. We critically re-examine the previous work on the subject, identify the physical mechanism responsible for the line narrowing, and spell out the conditions under which it occurs: * the lattice dynamics is inertial (not overdamped) * magnetic field is sufficiently strong, such that $`\omega _c`$ is larger than $`\omega _p`$ * the pinning is sufficiently weak so that $`\omega _p`$ is much smaller than the magnetophonon bandwidth * the compression modulus $`\lambda `$ of the WC evaluated at the pinning length is larger than its shear modulus $`\mu `$ (satisfied for long-range interaction). If any of the conditions above is violated, then the line is conventionally broad, $`\mathrm{\Delta }\omega _p\omega _p`$. Although we were unable to find the exact expression for $`\mathrm{\Delta }\omega _p`$, we show that there is an asymptotically exact relation between $`\mathrm{\Delta }\omega _p`$ at intermediate $`B`$ and the low-frequency tail of the phonon spectral function in zero $`B`$ (we remind the reader here that we study classical electrons, which form the WC at any $`B`$). Establishing such a connection and elucidating the physical mechanism of the line narrowing are our main achievements. The phonon spectral function at low frequencies is expected to have a power law form $$\mathrm{Im}D_{\alpha \alpha }(q=0,\omega )\omega ^{2s+1},$$ (1) with a nontrivial exponent $`s`$. This additional input and further considerations enable us to predict that as $`\omega _c`$ increases and becomes larger than $`\omega _p`$, $`\mathrm{\Delta }\omega _p/\omega _p`$ first decreases in a power-law fashion $$\mathrm{\Delta }\omega _p/\omega _p(\omega _p/\omega _c)^s$$ (2) but then eventually saturates at a value (see Fig. 1) $$\mathrm{\Delta }\omega _p/\omega _p(\mu /\lambda )^s$$ (3) (here $`\lambda `$ is meant to be evaluated at the pinning length). The recent experiments have probed the high-field regime where Eq. (3) is supposed to apply. In the aforementioned work of Chitra et al. a set of equations was derived, which, when analyzed further, also yield Eq. (3), with $`s=\frac{1}{2}`$. This specific value of $`s`$ can be traced down to the fact that Chitra et al. unknowingly rederive a common form of the self-consistent Born approximation (SCBA), where Eq. (1) with $`s=\frac{1}{2}`$ is satisfied. Since the SCBA is uncontrolled at small $`\omega `$, these results are unreliable. Moreover, appealing to the known properties of Lifshitz tails in other disordered systems, we argue that the SCBA strongly overestimates the low-frequency spectral function. We propose that the correct approach should be based on considering certain low-propability disorder configurations (“optimal fluctuations”), which leads to $$s=\frac{3}{2}.$$ (4) The paper is organized as follows. In Sec. II we discuss the ground state of the pinned WC, placing emphasis on the parameters that provide the input for the dynamical problem of interest. In Sec. III we formulate the model and the general framework for study of the finite-frequency response. In Sec. IV we derive Eq. (2). In Sec. V we study the soft phonon modes at $`B=0`$ and obtain the above formula for $`s`$. The comparison with the experiments is given in Sec. VII after a brief discussion of quantum and finite-temperature effects in Sec. VI. ## II Static properties The distortions of the lattice in the ground state of a collectively pinned Wigner crystal accumulate gradually over length scales much larger than the lattice constant $`a`$. Such distortions can be described in terms of a smooth displacement field $`𝐮^{(0)}(𝐫)`$. The Hamiltonian of the system $`H=H_{el}+H_p`$ is the sum of the elastic term, $`H_{el}={\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐪}{}}u_\alpha ^{(0)}(𝐪)H_{\alpha \beta }^0(𝐪)u_\beta ^{(0)}(𝐪),`$ (5) $`H_{\alpha \beta }^0=\left(\delta _{\alpha \beta }{\displaystyle \frac{q_\alpha q_\beta }{q^2}}\right)H_T^0(q)+{\displaystyle \frac{q_\alpha q_\beta }{q^2}}H_L^0(q),`$ (6) $`H_T^0(q)=\mu q^2,H_L^0(q)=\lambda (q)q^2,`$ (7) and the pinning term, $`H_p=n_e{\displaystyle \underset{𝐫}{}}U(𝐫)\left\{{\displaystyle \underset{𝐊}{}}\mathrm{cos}𝐊[𝐫𝐮^{(0)}(𝐫)]\mathbf{}𝐮^{(0)}\right\},`$ (8) $`U(𝐫)U(0)=C(r).`$ (9) The notations here are as follows. The $`𝐪`$-integral is taken over the Brillouin zone with measure $`d^2q/(2\pi )^2`$; the $`𝐫`$-integrals is taken over the area of the system with measure $`d^2r`$. Parameters $`\mu `$ and $`\lambda `$ represent the shear and bulk elastic moduli of the WC, respectively. In strong magnetic fields, i.e., at low filling factors they approach their classical values, $`\mu 0.245e^2n_e^{3/2}/\kappa `$ and $`\lambda (q)=5\mu +Y/\|q\|`$. Here $`n_e=2/\sqrt{3}a^2`$ is the average electron concentration, $`\kappa `$ is the dielectric constant of the medium, and $`Y=2\pi e^2n_e^2/\kappa `$. The infrared divergence of $`\lambda (q)`$ originates from the long-range Coulomb interaction. $`𝐊`$’s are the reciprocal lattice vectors. Finally, the correlator $`C(r)`$ is assumed to rapidly decay at $`r>\xi `$ where $`\xi <a`$ is the correlation length of the short-range Gaussian random potential $`U(𝐫)`$. The static properties of the system defined by Eqs. (59) can be investigated using classical statistical mechanics methods. Problems of this type have been studied extensively in the past, see Ref. for review and references, and also some recent contributions, Refs. , . The 2D WC generally conforms to the $`d=2`$, $`m=2`$ class of elastic media ($`d`$ is the total number of spacial dimensions and $`m`$ is the number of components of $`𝐮`$). However, some caution is required in transferring the known results to the WC case because the L- and T-components are not equivalent. The large disparity $`\lambda \mu `$ of the elastic moduli causes strongly suppression of the longitudinal distortions compared to the transverse ones. Therefore, in certain formulas one has to use $`m=1`$. In this respect, the WC is similar to the Abrikosov vortex lattice where the analogous inequality $`c_{11}c_{66}`$ exists. Let us define the two-point correlation function (roughness) $`W(r)=[𝐮^{(0)}(𝐫)𝐮^{(0)}(0)]^2^{1/2}`$. Many static properties of the system can be deduced from the behavior of $`W(r)`$. Three regimes can be distinguished: random-force, random-manifold, and the asympotic (Fig. 2). The random-force or Larkin regime appears for small $`r`$ where $`W(r)<\xi `$. The last inequality legitimizes expanding the exponential in Eq. (8) and keeping only the linear in $`u^{(0)}`$ terms. The Hamiltonian becomes quadratic, and its ground state is easily found to be $$u_\alpha ^{(0)}(𝐪)=in_e\underset{𝐊}{}[𝐇^0]_{\alpha \beta }^1(𝐪)K_\beta \stackrel{~}{U}(𝐪+𝐊).$$ (10) Here and below tildes denote Fourier transforms. Averaging over the disorder, we arrive at $`W^2(𝐫)`$ $``$ $`V^{}(0){\displaystyle \underset{𝐪}{}}(1e^{i\mathrm{𝐪𝐫}})\left[(H_T^0)^2+(H_L^0)^2\right]`$ (11) $``$ $`r^2\mathrm{ln}(R_c/r),`$ (12) $`V(z)`$ $``$ $`{\displaystyle \frac{n_e^2}{2}}{\displaystyle \underset{𝐊}{}}\stackrel{~}{C}(K)e^{zK^2/2}.`$ (13) To obtain Eq. (12) we cut off the log-divergent integral in Eq. (11) at $`q=q_c`$ where $`q_c=\pi /R_c`$ and $`R_c`$ is the Larkin length defined by the relation $`W(R_c)=\xi `$. For a weak random potential $`R_c`$ is much larger than $`a`$ and is given by $$R_c\mu a\xi ^2[C(0)]^{1/2}.$$ (14) The same estimate for $`R_c`$ can be obtained in a simpler way, by equating the elastic energy density $`\mu \xi ^2/R_c^2`$ to the pinning energy density $`[C(0)]^{1/2}\sqrt{N}/R_c^2`$. In the last expression $`N=n_eR_c^2`$ is the number of electrons within the area $`R_c\times R_c`$ and the square root takes into account that individual pinning energies have random signs. The elastic-manifold regime is realized at intermediate distances, $`R_c<r<R_a`$ where $`\xi <W(r)<a`$. Here the roughness grows slower than in the Larkin regime, $`W(r)r^\zeta `$, $`\zeta <1`$. For the 2D Wigner crystal the Flory estimate $`\zeta =\frac{1}{3}`$ (as appropriate for the $`d=2`$, $`m=2`$ case) should be accurate but perhaps not exact. The slower growth rate of $`W(r)`$ is due to a less efficient mechanism by which the lattice adjusts to the random potential on length scales exceeding $`R_c`$. The steepest descent towards the nearest local minima of the potential landscape in the Larkin regime gets replaced by selection among a set of many possible yet roughly equivalent nearby minima. The third and the last regime begins when $`W(r)`$ reaches $`a`$ at the length scale $`R_a=R_c(a/\xi )^{1/\zeta }`$. The accumulation of the pinning energy density at such large length scales becomes additionally suppressed by lattice periodicity effects. Indeed, a uniform shift $`uu+a`$ reduces to relabeling the lattice sites with no physical significance. As a result, the roughness grows only logarithmically with distance, $`W(r)a\mathrm{ln}(r/R_a)`$. The asymptotic $`\mathrm{ln}r`$-roughness corresponds to a “quasi-short-range” order: $`W(r)`$ increases so slowly that it is barely enough to remove the algebraic divergencies (Bragg peaks) of the static structure factor $$F(𝐪)=n_e^2\underset{𝐫}{}\mathrm{exp}\left[i(𝐪𝐊)𝐫\frac{1}{2}K^2W^2(r)\right]$$ (15) at the reciprocal lattice vectors $`𝐪=𝐊`$. Remarkably, it also creates enough elastic stress to generate unbound dislocations, on the length scale $`\xi _DR_a\mathrm{exp}[C_0\mathrm{ln}^{1/2}(R_a/a)]`$. At this scale $`𝐮`$ becomes multi-valued and Eq. (5) must be modified. Although very interesting, the random-manifold and the asymptotic regimes do not play much role in the further discussion. We mentioned them here for the sake of a more comprehensive introduction and also to draw attention to the difference between two characteristic length scales, $`R_c`$ and $`R_a`$. $`R_a`$ is the length that seems to be important only for rather subtle properties mentioned above: the destruction of the Bragg peaks and the appearance of the topological defects. In contrast, $`R_c`$ determines most of the physically important response characteristics such as the threshold electric field in the dc transport and the frequency $`\omega _p`$ of the pinning mode. This role of $`R_c`$ originates from the fact that it is the length scale where the dominant contribution $`\mu \xi ^2/R_c^2`$ to the pinning energy density is accumulated. In this sense $`R_c`$ plays the role of the pinning length mentioned in Sec. I. Let us now elaborate on the difference between the T- and L-correlation functions mentioned above. To do so we decompose the elastic displacement into its T- and L-components, $$𝐮(𝐪)=\widehat{𝐪}u_L(𝐪)+[\widehat{𝐳}\times \widehat{𝐪}]u_T(𝐪),$$ (16) where $`\widehat{𝐪}=𝐪/\|𝐪\|`$, and define the two components of roughness $`W_\alpha =u_\alpha ^{(0)}(𝐫)u_\alpha ^{(0)}(0)^{1/2}`$, $`\alpha =L,T`$. First of all, let us examine the short-distance behavior using Eq. (10). For $`W_L(r)`$ we obtain the expression similar to Eq. (11): $$W_L^2(r)V^{}(0)\underset{𝐪}{}\frac{1e^{i\mathrm{𝐪𝐫}}}{[H_L^0(q)]^2}\frac{V^{}(0)}{2\pi Y^2}\mathrm{ln}\left(\frac{r}{a}\right),$$ (17) and so $`W_L(r)`$ is much smaller than $`W_T(r)W(r)`$ given by Eq. (12). Formula (17) is certainly valid for $`r<R_c`$. For larger $`r`$ we can no longer treat the elastic distortion $`\|𝐮^{(0)}(𝐫)𝐮^{(0)}(0)\|`$ as small. Nonetheless, we expect that the region of validity of Eq. (17) extends somewhat beyond $`r=R_c`$. Indeed, instead of doing the perturbation theory around the ideal lattice state, we can do it around the ground state $`𝐮_{\mathrm{}}^{(0)}`$ of the incompressible crystal, $`\lambda =\mathrm{}`$. In the second approach we have to account for a finite pinning energy density, but as long as the longitudinal elastic stiffness $`Yq`$ where $`q\pi /r`$ exceeds the characteristic curvature $`\mu q_c^2`$ of the pinning energy density landscape, i.e., for $`rR_L=(Y/\mu )R_c^2R_c^2/a`$, this should be unimportant. Substituting $`r=R_L`$ into Eq. (17) we find that $`W_L(R_L)\xi `$, contrary to some previous suggestions that $`W_L(R_L)W_T(R_L)`$ \[recall that $`W_T(R_L)W(R_L)W(R_c)=\xi `$\]. The large-distance ($`r>R_L`$) behavior of $`W_L(r)`$ is unaccessible within the perturbation theory. However, we should get qualitatively correct results by using the Gaussian variational replica method (GVRM). Assuming for simplicity that $`R_LR_a`$ and suitably modifying GVRM for the present problem, we find $$W_L^2(r)V^{}(0)\underset{𝐪}{}\frac{1\mathrm{cos}(\mathrm{𝐪𝐫})}{H_L^0(q)[H_L^0(q)+\mathrm{\Pi }_a]},$$ (18) where $`\mathrm{\Pi }_a=\mu R_a^2`$. It is easy to see that $`W_L(r)`$ tends to a finite limiting value as $`r\mathrm{}`$. So, the long-range positional order is absent in the T- but is preserved in the L-part of the elastic displacement field. To conclude the review of the statics, let us consider a matrix of the second derivatives of $`H_p`$, $$𝐒(𝐫)=\frac{^2H_p}{𝐮^{(0)}𝐮^{(0)}}=n_e\underset{𝐊}{}\mathrm{𝐊𝐊}U(𝐫)\mathrm{cos}𝐊[𝐫𝐮^{(0)}(𝐫)].$$ (19) Matrix $`𝐒`$ plays an important role in the dynamics and we need to understand its properties. In the case of a short-range disorder we study here, the correlations in $`U(𝐫)`$ decay more rapidly with $`r`$ than those of $`𝐮^{(0)}(𝐫)`$. This enables us to find the two-point correlator of $`S_{\alpha \beta }`$ with little effort. It suffices to use the approximation $`𝐮^{(0)}(𝐫)=0`$ in Eq. (19), which leads to $`\stackrel{~}{S}_{\alpha \mu }(𝐤_1)\stackrel{~}{S}_{\nu \beta }(𝐤_2)\stackrel{~}{S}_{\alpha \mu }\stackrel{~}{S}_{\nu \beta }{\displaystyle \frac{V^{\prime \prime }(0)}{2}}`$ (20) $`\times (\delta _{\alpha \mu }\delta _{\nu \beta }+\delta _{\alpha \beta }\delta _{\mu \nu }+\delta _{\alpha \nu }\delta _{\mu \beta })\delta _{𝐤_1𝐤_2}L^2`$ (21) for the WC with the hexagonal lattice and in the long-wavelength limit $`k_1,k_21/a`$. Here $`L^2`$ is the area of the system. Another important property of $`S_{\alpha \beta }`$ is a nonzero mean, $$S_{\alpha \beta }(𝐫)=\delta _{\alpha \beta }S_0.$$ (22) Parameter $`S_0`$ is positive, in accordance with the pinning phenomenon: the crystal in its ground state is distorted in such a way that the electrons are preferentially situated near the minima of the random potential $`U(𝐫)`$ where its curvature is positive. $`S_0`$ can be estimated to be $$S_0\frac{V^{\prime \prime }(0)}{2\pi \mu }\mathrm{ln}\left(\frac{R_c}{a}\right)$$ (23) using the same procedure as for deriving Eq. (10). Assuming that $`\|𝐮^{(0)}(𝐫)𝐮^{(0)}(0)\|`$ is “small,” we expand the exponential in Eq. (19) to obtain $$\stackrel{~}{S}_{\alpha \beta }(𝐪)\frac{in_e}{L^2}\underset{𝐊}{}K_\alpha K_\beta K_\gamma \underset{𝐪_1}{}u_\gamma ^{(0)}(𝐪_1)\stackrel{~}{U}(𝐪𝐪_1𝐊).$$ (24) Combining Eqs. (10) and (24), averaging over the disorder, and using the inequality $`H_L^0(q)H_T^0(q)`$, we get $$S_0V^{\prime \prime }(0)\underset{𝐪}{}[H_T^0(q)]^1,$$ (25) which leads to Eq. (23). Here we again had to cut off the infrared logarithmic divergence by hand, at $`q=q_c=\pi /R_c`$. More sophisticated approximation schemes such as the GVRM would implement this cutoff more gracefully but as long as we are not interested in the numerical factor in the argument of the log, the result is the same. At this point we conclude the discussion of the ground state properties of the pinned WC, as we are now ready to address its dynamical response. ## III Finite-frequency response: perturbation theory and qualitative considerations To study the dynamics of the WC we introduce a time-dependent displacement field $`𝐮`$, which is the deviation of the total displacement from its ground-state value $`𝐮^{(0)}(𝐫)`$. We will restrict ourselves to the harmonic approximation where the action is quadratic in $`𝐮`$. The response of a harmonic oscillator is the same in quantum and classical mechanics, so we can use the convenient imaginary-time formalism without compromising our original intention to treat the system classically. The resultant action contains two terms: one which describes a uniformly pinned WC, $$A_0=\frac{1}{2}\frac{1}{\mathrm{}\beta }\underset{\omega _n}{}\underset{𝐪}{}𝐮^{}(𝐪,i\omega _n)[𝐃^0]^1𝐮(𝐪,i\omega _n),$$ (26) and the other which takes into account fluctuations in local pinning strength, $$A_1=\frac{1}{2}\frac{1}{\mathrm{}\beta }\underset{\omega _n}{}\underset{𝐪_1}{}\underset{𝐪_2}{}𝐮^{}(𝐪_1,i\omega _n)\delta \stackrel{~}{𝐒}(𝐪_1𝐪_2)𝐮(𝐪_2,i\omega _n).$$ (27) Here $`\beta =1/(k_BT)`$ is the inverse temperature, $`\omega _n=2\pi n/(\mathrm{}\beta )`$ are the bosonic Matsubara frequencies, $`\delta \stackrel{~}{𝐒}\stackrel{~}{𝐒}S_0𝐈`$, $`𝐈`$ is the identity matrix, $`𝐃^0=𝐃^0(𝐪,i\omega _n)`$ is the phonon propagator of a uniformly pinned WC (cf. Refs. and ) $$[𝐃^0]^1=𝐑_𝐪^{}\left[\begin{array}{cc}H_T^0+S_0+\rho \omega _n^2& \rho \omega _n\omega _c\\ \rho \omega _n\omega _c& H_L^0+S_0+\rho \omega _n^2\end{array}\right]𝐑_𝐪,$$ (28) $`\rho =m_en_e`$ is the average mass density, and $`𝐑_𝐪`$ is the $`O(2)`$ rotation by angle $`\mathrm{arg}(q_x+iq_y)`$. We also define the disorder-averaged propagator $`𝐃(𝐪,i\omega _n)`$, $$𝐃(𝐪,i\omega _n)=\frac{𝒟𝐮𝒟𝐮^{}𝐮(𝐪,i\omega _n)𝐮^{}(𝐪,i\omega _n)e^A}{(\mathrm{}\beta L^2)𝒟𝐮𝒟𝐮^{}e^A},$$ (29) where $`AA_0+A_1`$ and $`L^2`$ is again the area of the system. The quantity we set out to calculate is the ac conductivity $`𝝈(q,\omega )`$, which in this model is given by $$𝝈(q,\omega )=ie^2n_e^2\omega 𝐃(q,i\omega _n)\|_{i\omega _n\omega +i0}.$$ (30) The conductivity can also be expressed in terms of the phonon self-energy $$𝚷(𝐪,\omega )S_0𝐈+\left(𝐃^1[𝐃^0]^1\right)\|_{i\omega _n\omega +i0}.$$ (31) We are interested primarily in the case $`q=0`$ where the most general form of $`𝚷(𝐪,\omega )`$ consistent with rotational symmetry is $$\mathrm{\Pi }_{\alpha \beta }(0,\omega )=\delta _{\alpha \beta }\mathrm{\Pi }(\omega )i\epsilon _{\alpha \beta }\rho \omega \omega _cf_{xy}(\omega ),$$ (32) $`f_{xy}(\omega )`$ being the relative correction $`\mathrm{\Delta }\rho _{xy}/\rho _{xy}^0`$ to the bare Hall resistivity $`\rho _{xy}^0(\omega )=B/(n_eec)`$. Approximate calculations presented in this paper give $`f_{xy}=0`$, which leads us to believe that $`f_{xy}(\omega )`$ must be small for weak pinning. We choose to neglect such fine details and to assume that $`f_{xy}`$ vanishes. In this case $`𝚷(0,\omega )`$ is a scalar and $`\sigma _{xx}(\omega )`$ is given by $$\mathrm{Re}\sigma _{xx}(\omega )=e^2n_e^2\omega \mathrm{Im}\frac{\mathrm{\Pi }(0,ϵ)ϵ}{[\mathrm{\Pi }(0,ϵ)ϵ]^2ϵϵ_c},$$ (33) where we introduced convenient “energy” variables $$ϵ\rho \omega ^2,ϵ_c\rho \omega _c^2.$$ (34) From Eq. (33) one can see that in strong magnetic fields, $`\omega _c\omega _p`$, the power absorption in the uniform electric field takes place mainly within the frequency interval of width $$\mathrm{\Delta }\omega _p=\frac{\mathrm{Im}\mathrm{\Pi }(0,ϵ_p)}{\rho \omega _c}$$ (35) centered at $`\omega =\omega _p`$, where $$\omega _p=\sqrt{\frac{ϵ_p}{\rho }}=\frac{\mathrm{Re}\mathrm{\Pi }(0,ϵ_p)}{\rho \omega _c}.$$ (36) The last equation is the implicit definition of $`ϵ_p`$. In the following we will work predominantly with the “energy” variable $`ϵ`$ rather than the frequency $`\omega `$. At small $`ϵ`$ the calculation of $`𝚷`$ is a difficult problem but at large energies, $`ϵϵ_p`$, the first Born approximation suffices. The only parameters needed to implement it are the mean and the variance of the matrix elements of $`𝐒`$ given by Eqs. (21) and (22). For $`q1/a`$ we obtain that $`\mathrm{\Pi }_{\alpha \beta }(𝐪,ϵ)=\delta _{\alpha \beta }\mathrm{\Pi }(ϵ)`$, where $$\mathrm{\Pi }(ϵ)=S_0+V^{\prime \prime }(0)\underset{𝐤}{}\mathrm{tr}𝐃^0(𝐤,ϵ).$$ (37) $`\mathrm{\Pi }`$ has both imaginary and real parts, corresponding to the broadening and the frequency shift of the pinning mode, Eqs. (35) and (36), respectively. The imaginary part comes solely from the pole(s) of $`𝐃^0(𝐤,ϵ)`$, i.e., the solutions of $`[H_T^0(𝐤)+S_0ϵ][H_L^0(𝐤)+S_0ϵ]ϵϵ_c=0`$. A more accurate expression for $`\mathrm{\Pi }`$ can be obtained using the self-consistent Born approximation (SCBA), $$\mathrm{\Pi }(ϵ)=S_0+V^{\prime \prime }(0)\underset{𝐤}{}\mathrm{tr}𝐃(𝐤,ϵ),$$ (38) as long as the full propagator $`𝐃(𝐤,ϵ)`$ has a pole at $`\|k\|q_c=\pi /R_c`$. Under this condition the diagrams with intersecting lines not included into the SCBA are suppressed, much like they are suppressed in a dirty metal with $`k_Fl^1`$, $`k_F`$ being the Fermi momentum and $`l`$ being the mean free path. The indicated condition is satisfied when $$ϵ\mathrm{min}\{\mathrm{\Pi }_0,\frac{\lambda }{\mu }\frac{\mathrm{\Pi }_0^2}{ϵ_c}\}.$$ (39) Here and below $`\lambda `$ is meant to be evaluated at $`q=q_c`$ unless it is indicated otherwise. Parameter $`\mathrm{\Pi }_0`$, which can be estimated to be $$\mathrm{\Pi }_0\mu q_c^2$$ (40) represents the real part of the self-energy at the lowest $`ϵ`$ allowed by inequality (39). $`\mathrm{\Pi }_0`$ is smaller than $`S_0`$ by the logarithmic factor due to the partial cancellation between the first and the second terms in Eq. (38). This unfortunate cancellation leaves us only with the order of magnitude estimate (40). When inequality (39) is violated, the SCBA becomes an uncontrolled approximation. The only SCBA result, which can presumably be trusted at such $`ϵ`$, is a slow dependence of $`\mathrm{Re}\mathrm{\Pi }(ϵ)`$ on energy. This is because the real part of the integral in Eq. (38) is dominated by large $`k`$’s where the lowest-order perturbation theory is valid. To estimate $`ϵ_p`$ we can take $`\mathrm{Re}\mathrm{\Pi }(ϵ)\mathrm{\Pi }_0`$, which yields $$ϵ_p\mathrm{min}\{\mathrm{\Pi }_0,\frac{\mathrm{\Pi }_0^2}{ϵ_c}\}.$$ (41) represented graphically in Fig. 3 together with the SCBA domain. As one can see from this graph, for $`ϵ_c<\mathrm{\Pi }_0`$ (weak magnetic fields) there is no parametric separation between the SCBA domain and the curve $`ϵ=ϵ_p`$; therefore, the SCBA is qualitatively correct even along this curve where it predicts $`\mathrm{Im}\mathrm{\Pi }(ϵ_p)\mathrm{\Pi }_0`$. Hence, $`\mathrm{\Delta }\omega _p\omega _p`$. The resultant broad absorption maximum is sketched in Fig. 4. If $`\lambda `$ and $`\mu `$ are comparable, the same argument works also in strong magnetic fields ($`ϵ_c\mathrm{\Pi }_0`$) where $`\sigma _{xx}(\omega )`$ looks qualitatively similar, except the position of the maximum depends on the magnetic field: $`\omega _p=\omega _{p0}^2/\omega _c`$, where $`\omega _{p0}`$ is the pinning frequency at $`B=0`$. On the other hand, if $`\lambda \mu `$, as for the WC, the pinning mode is situated in the interior of a parametrically wide range of $`ϵ`$ where the SCBA fails. Thus, the calculation of $`\mathrm{\Delta }\omega _p`$ requires other methods. Later we will show that $`\mathrm{\Delta }\omega _p`$ in strong magnetic fields is related to the low-frequency tail of $`\sigma _{xx}(\omega )`$ in the absence of the magnetic field. Although calculating the functional form of such a tail is another nontrivial problem (after all, it is beyond the SCBA!), this intriguing relation is sufficient to establish that the absorption line narrows down, in agreement with the experiments. We unfold our argument gradually over the remaining sections. In this section we start implementing this task by clarifying why $`\sigma _{xx}(\omega )`$ is virtually independent of the longitudinal stiffness $`\lambda `$ provided $`\lambda \mu `$ and $`B=0`$. Since action $`A`$ is quadratic in $`𝐮_L`$, the L-phonon degrees of freedom can be easily integrated out, leading to the effective Hamiltonian for the T-phonons, $$𝐇_T=𝐇_T^0+𝐒_T𝐒_{X+}^{}(𝐇_L^0+𝐒_Lϵ𝐈)^1𝐒_{X+},$$ (42) where $`𝐇_T^0`$, $`𝐇_L^0`$, and $`𝐒_i`$ should be understood as operators. The first two are diagonal in the basis of plane waves, with matrix elements $`H_T^0(q)`$ and $`H_L^0(q)`$, respectively; $`𝐒_i`$’s have both diagonal and off-diagonal matrix elements: $`S_{L\mathrm{𝐪𝐪}^{}}`$ $`=`$ $`\stackrel{~}{S}_{\alpha \beta }(𝐪𝐪^{})\widehat{q}_\alpha \widehat{q}_\beta ^{},`$ (44) $`S_{T\mathrm{𝐪𝐪}^{}}`$ $`=`$ $`(\delta _{\alpha \beta }\mathrm{tr}\stackrel{~}{𝐒}\stackrel{~}{S}_{\alpha \beta })\widehat{q}_\alpha \widehat{q}_\beta ^{},`$ (45) $`S_{X\mathrm{𝐪𝐪}^{}}`$ $`=`$ $`\stackrel{~}{S}_{\alpha \beta }(𝐪𝐪^{})ϵ_{\beta \gamma }\widehat{q}_\alpha \widehat{q}_\gamma ^{}.`$ (46) Finally, $`𝐒_{X+}𝐒_X+\sqrt{ϵϵ_c}𝐈`$. Integrating out $`𝐮_L`$ is an exact algebraic transformation, which preserves the spectrum of the collective modes: the resolvent $`𝐃_T(ϵ)`$ of the operator $`𝐇_Tϵ𝐈`$ has poles at the same $`ϵ`$ as the full propagator. In fact, $`𝐃_T`$, which can be written in the form $$𝐃_T=([𝐃_T^{}]^1𝐒^{})^1,$$ (47) where $`𝐃_T^{}=(𝐇_T^0+𝐒_Tϵ𝐈)^1`$ (48) $`𝐒^{}=𝐒_{X+}^{}(𝐇_L^0+𝐒_Lϵ𝐈)^1𝐒_{X+},`$ (49) is nothing else than the T-T component of the full phonon propagator before the disorder averaging. The self-energy of the averaged propagator at $`q=0`$ is a scalar (see above); therefore, the T-T component determines the entire propagator for this particular $`q`$. Let us now show that in zero magnetic field the mixing between the T- and L-phonons represented by $`𝐒^{}`$ has little effect on $`𝐃_T`$. Indeed, in the diagrammatic expansion of $`𝐃_T(ϵ)`$ in powers of $`\delta 𝐒_T`$ and $`𝐒^{}`$, the typical momentum transfer $`k`$ for $`ϵ\mathrm{\Pi }_0`$ is of the order of $`q_c`$. A single occurence of $`𝐒^{}`$ contributes $`k^2\|V^{\prime \prime }(0)\|/\lambda k^2(\mu /\lambda )\mathrm{\Pi }_0`$ to the self-energy, compared to $`\mathrm{\Pi }_0`$ from $`\delta 𝐒_T`$. Diagrams with multiple occurences of $`𝐒^{}`$ are suppressed by even higher powers of the small parameter $`\mu /\lambda `$. Thus, in the parameter range relevant for the observation of the pinning mode $$𝐃_T(𝐪,ϵ)𝐃_T^{}(𝐪,ϵ)\frac{1}{H_T^0(q)+\mathrm{\Pi }^{}(𝐪,ϵ)ϵ}.$$ (50) Note that $`𝐃_T^{}`$ is determined only by the shear modulus $`\mu `$ and the disorder in the T-T channel $`\delta 𝐒_T`$, while bulk modulus $`\lambda `$ drops out. It follows from this discussion that for $`\lambda \mu `$ we can calculate the $`q=0`$ response pretending that the L-degree of freedom does not exist. Let us indeed imagine that the L-phonons are forbidden and try to investigate the nature of the T-phonon eigenstates that would compose the pinning mode, i.e., the states that would respond to the uniform electric field (still at $`B=0`$). To do so we need to analyze the solutions of the eigenvalue problem for the single-particle Hamiltonian $`𝐇^{}=𝐇_T^0+𝐒_T`$. In the uniformly pinned WC, $`𝐒_T=S_0𝐈`$ and the eigenstates are just the plane waves labelled by momenta $`𝐤`$ of the Brillouin zone. In the actual random system the eigenstates are wavepackets of the plane waves with characteristic spread of momenta of the order of $`L_T^1`$, where $`L_Tv_T/\|\mathrm{Im}\mathrm{\Pi }\|`$ is the T-phonon mean free path, $`v_T=\sqrt{ϵ\mu }`$ playing the role T-phonon “velocity.” Since $`\|\mathrm{Im}\mathrm{\Pi }(ϵ_p)\|ϵ_p\mu q_c^2`$, at $`ϵϵ_p`$ we have $`L_TR_c`$. Hence, the eigenstates that respond appreciably to a $`q=0`$ external drive are wavepackets built from $`0<k<q_c`$ plane waves. For such wavepackets the average $`k`$ is of the order of the inverse mean free path and the Ioffe-Regel criterion maintains that these states are localized. In other words, $`R_c`$ is not only the mean free path but also the localization length of the states within the pinning mode. Let us define the $`B=0`$ T-phonon density of states, $$\nu (ϵ)=\frac{1}{L^2}\underset{i}{}\delta (ϵϵ_i)=\frac{1}{\pi }\underset{𝐤}{}𝐃_T^{}(𝐤,ϵ).$$ (51) At $`ϵϵ_p`$ it tends to the bare value $`\nu _0=(4\pi \mu )^1`$, and at $`ϵϵ_p`$ it is only slightly smaller. It is easy to see then that within an area $`R_c\times R_c`$ and energy interval $`0<ϵ<ϵ_p`$, there is typically only one (localized) state. The picture that emerges is illustrated in Fig. 5: we have a collection of localized states of roughly the same size (i.e., the inverse participation ratio) and roughly the same distance from each other in real space, both of the order of $`R_c`$. The broad distribution of shapes and sizes characteristic of localized states in disordered systems yields a large inhomogeneous broadening $`\mathrm{\Delta }\omega _p\omega _p`$ of the absorption line, as we found earlier based on a different line of reasoning. In early works on pinning a heuristic imagery of “domains” was sometimes invoked to describe the ground state structure. However, it was never clear where to draw the boundaries between different domains. This difficulty is partially resolved in our picture of phonon localization where domains can be defined as areas where individual localized phonon wavefunctions are appreaciable. However, modern understanding of the subject reviewed in the previous section no longer appeals to any “domains,” and so this interpretation may be regarded as a historic sentiment. One more historic comment is in order. Previous work on phonon localization found that the phonon eigenstates remain extended even in the $`\omega 0`$ limit, seemingly in contradiction to our results. In fact, there is no contradiction because the model we study here and the conventional formulations, where disorder originates from the defects of the crystalline lattice, e.g., substitutions by impurity atoms, are quite different. While in those conventional models the defects move with the lattice and their influence rapidly diminishes at small $`𝐤`$’s, the pinning potential for the WC is static, so it is not affected by the lattice motion. It couples to the phonons via the matrix $`𝐒`$ whose fluctuations have approximately the same rms magnitude at all $`𝐤`$’s within the Brillouin zone, see Eq. (21). This is the reason why in our case the disorder effects are much stronger, causing the localization of the phonon eigenstates. Actually, our model has more in common with those of noninteracting electrons in dirty metals and semiconductors. (At the one-particle level the distinction between the bosonic statistics of phonons and the fermionic statistics of electrons is irrelevant). Concluding this section, we would like to reiterate one of its main results, that the presence (or absence) of the stiff L-phonon degree of freedom affects the $`B=0`$, $`q=0`$ response only weakly. Naturally, it makes the fate of the L-phonons seem quite mysterious. The finite-$`B`$ dynamical response has also been barely touched upon. These gaps in understanding will be filled in the next section. ## IV Scattering of L-phonons and pinning mode linewidth in a finite magnetic field In the previous section we discussed the scattering and localization of the T-component of WC lattice vibrations. Let us now turn to the L-component. This time we integrate out $`𝐮_T`$ to obtain the effective Hamiltonian for the L-phonons, the corresponding Green’s function, and the self-energy: $`𝐇_L=𝐇_L^0+𝐒_L𝐒_{X+}(𝐇_T^0+𝐒_Tϵ𝐈)^1𝐒_{X+}^{},`$ (52) $`𝐃_L=(𝐇_Lϵ𝐈)^1,`$ (53) $`𝚷_L(ϵ)=𝐃_L^1(𝐇_L^0ϵ𝐈).`$ (54) Since the self-energy $`𝚷(𝐪,ϵ)`$ of the original system (before integrating out $`𝐮_T`$) is a scalar at $`𝐪=0`$, $`\mathrm{\Pi }_L`$ must obey the relation $$\mathrm{\Pi }_L(0,ϵ)=\mathrm{\Pi }(0,ϵ)\frac{ϵϵ_c}{\mathrm{\Pi }(0,ϵ)ϵ}.$$ (55) Our first task is to reproduce the results of Sec. III by demonstrating that in sufficiently weak magnetic fields $`\mathrm{\Pi }_L`$ indeed has the above form with $`\mathrm{\Pi }\mathrm{\Pi }^{}`$. There are two types of processes that contribute to $`\mathrm{\Pi }_L`$: the intraband scattering (L-L channel) and the scattering through an intermediate T-state (L-T-L channel). They originate from the second and the third terms of $`𝐇_L`$, respectively. The intraband scattering will be treated in more detail in Appendix A where using a combination of the renormalization group and instanton methods we show that it is exponentially small. For our current purposes a weaker result is sufficient, that the intraband contribution $`\mathrm{\Pi }_{LL}`$ to $`\mathrm{\Pi }_L`$ is much smaller than $`\mathrm{\Pi }_0`$. This can be established without detailed calculations. Indeed, the effective random potential $`𝐒_L`$ in the L-L channel has basically the same mean and rms fluctiations as $`𝐒_T`$; therefore, $`𝐒_L`$ and $`𝐒_T`$ can be regarded as disorder of approximately the same strength. Then the alluded inequality $`\|\mathrm{\Pi }_{LL}\|\|\mathrm{\Pi }\|\mathrm{\Pi }_0`$ follows simply from the fact that L-phonons have much steeper bare dispersion relation $`ϵ(q)=\lambda (q)q^2`$ than T-phonons and are scattered much weaker by disorder of the same strength. It is legitimate to neglect the L-L scattering altogether by replacing $`𝐒_L`$ with its average value $`S_0𝐈`$, and to concentrate exclusively on the L-T-L processes. In this case, much like for a dirty semiconductors with two bands of carriers, the diagrammatic expansion of $`\mathrm{\Pi }_L`$ can be formulated as the sum of all one-particle irreducible (1PI) graphs involving $`𝐒_{X+}`$ and $`𝐃_T`$, see Eqs. (4749) and Fig. 6a–b. The crucial point substantiated below is that $`𝐒^{}`$ can be neglected, i.e., that $`𝐃_T`$ can be replaced by $`𝐃_T^{}`$, provided $`\lambda \mu `$ and the magnetic field is not too strong, $$ϵ_c(\lambda /\mu )\mathrm{\Pi }_0.$$ (56) For such $`ϵ_c`$ only two graphs contribute to $`\mathrm{\Pi }_L`$ (Fig. 6c), which is a significant simplification. However, we still have to specify how to average over the disorder. Naive averaging of $`𝐒_X`$ and $`𝐃_T^{}`$ separately from each other yields $$\mathrm{\Pi }_L(𝐪,ϵ)=S_0ϵϵ_c𝐃_T^{}(𝐪,ϵ)+V^{\prime \prime }(0)\underset{𝐤}{}𝐃_T^{}(𝐤,ϵ)$$ (57) but it is certainly not accurate because $`𝐒_X`$ and $`\delta 𝐒_T`$ are not independent. Indeed, their two-point correlator $`S_{X𝐪_1𝐪}S_{T𝐪_2𝐪_3}`$ $``$ $`(\widehat{𝐪}_1\widehat{𝐪}_2)(\widehat{𝐪}_3\widehat{𝐪})+(\widehat{𝐪}_1\widehat{𝐪}_3)(\widehat{𝐪}_2\widehat{𝐪})`$ (58) $`+`$ $`(\widehat{𝐪}_1\widehat{𝐪})(\widehat{𝐪}_3\widehat{𝐪}_2)`$ (59) is in general nonvanishing. In the special case of $`𝐪=0`$ a further progress is possible. Comparing the above equation with $`S_{T𝐪_1𝐪}S_{T𝐪_2𝐪_3}`$ $``$ $`(\widehat{𝐪}_1\widehat{𝐪}_2)[\widehat{𝐪}\times \widehat{𝐪}_3]\widehat{𝐳}+(\widehat{𝐪}_1\widehat{𝐪}_3)[\widehat{𝐪}\times \widehat{𝐪}_2]\widehat{𝐳}`$ (60) $`+`$ $`(\widehat{𝐪}_1\widehat{𝐪})[\widehat{𝐪}_3\times \widehat{𝐪}_2]\widehat{𝐳},`$ (61) we observe that the former is transformed into the latter by the replacement $`𝐪[\widehat{𝐳}\times 𝐪]`$. If the incoming momentum $`𝐪`$ is zero, rotating it by $`\pi /2`$ has no effect. Thus, for $`𝐪=0`$ we can replace $`𝐒_X`$ in the second diagram of Fig. 6c by $`\delta 𝐒_T`$, after which it becomes identical to the totality of diagrams in Fig. 6a for $`\mathrm{\Pi }^{}`$; therefore, $$\mathrm{\Pi }_L(0,ϵ)=\mathrm{\Pi }^{}(0,ϵ)\frac{ϵϵ_c}{\mathrm{\Pi }^{}(0,ϵ)ϵ}.$$ (62) As explained in Sec. III, $`\mathrm{\Pi }^{}\mathrm{\Pi }`$; hence, we will succeed in reproducing Eq. (55) as soon as we show that inequality (56) is indeed the relevant condition for dropping $`𝐒^{}`$ in the diagrammatic series for $`\mathrm{\Pi }_L`$. To do so we develop one step further the argument of Sec. III. One-particle irreducible diagrams containing $`𝐒^{}`$ have at least one L-phonon line with typical momentum transfer $`kq_c`$, see Fig. 6d. If $`ϵ_c=0`$ such diagrams are suppressed by at least a factor of $`\mu /\lambda `$, as we found previously in Sec. III. However, if $`ϵ_c0`$, these diagrams also generate terms of the order of $$\sqrt{ϵϵ_c}\frac{1}{H_L^0(k)}\sqrt{ϵϵ_c}\frac{ϵϵ_c}{\lambda q_c^2},$$ (63) which act as an additional self-energy correction to $`𝐃_T^{}`$ in Fig. 6c. They can be approximately accounted for if we replace $`ϵ`$ in the argument of $`\mathrm{\Pi }^{}`$ in Eq. (62) by a renormalized value $`\stackrel{~}{ϵ}`$, $`\mathrm{\Pi }_L(0,ϵ)`$ $`=`$ $`\mathrm{\Pi }^{}(0,\stackrel{~}{ϵ}){\displaystyle \frac{ϵϵ_c}{\mathrm{\Pi }^{}(0,\stackrel{~}{ϵ})ϵ}},`$ (64) $`\stackrel{~}{ϵ}`$ $`=`$ $`ϵ+C_1{\displaystyle \frac{ϵϵ_c}{\lambda q_c^2}},C_11.`$ (65) The difference between $`ϵ`$ and $`\stackrel{~}{ϵ}`$ can be safely ignored as long as the inequality (56) is satisfied, in which case Eq. (62) is asymptotically exact. In stronger fields Eq. (64) applies, which gives only the order of magnitude of $`\mathrm{\Pi }_L`$ because of the uncertainty in the parameter $`C_1`$. The meaning of Eqs. (62) and (64) can be elucidated returning to the the original formulation where both T- and L-degrees of freedom are present: $`\mathrm{\Pi }(0,ϵ)`$ $`=`$ $`\mathrm{\Pi }^{}(0,ϵ),ϵ_c(\lambda /\mu )\mathrm{\Pi }_0,`$ (67) $`=`$ $`\mathrm{\Pi }^{}(0,\stackrel{~}{ϵ}),ϵ_c(\lambda /\mu )\mathrm{\Pi }_0.`$ (68) We see that $`\mathrm{\Pi }(0,ϵ)`$ remains magnetic field independent up to rather high $`B`$. As explained above, this occurs because the enhancement of the L-T mixing due to the Lorentz force is strongly impeded by the disparity of the two elastic moduli. Let us now evaluate the consequences of the obtained equations (67) and (68). In the intermediate field range, $`\mathrm{\Pi }_0ϵ_c(\lambda /\mu )\mathrm{\Pi }_0`$, where Eq. (67) applies, the consequences are two-fold: the suppression of $`ϵ_p`$, which is well known, and the suppression of the pinning mode linewidth, which is nontrivial. Indeed, since $`ϵ_p`$ becomes much less than $`\mathrm{\Pi }_0`$ \[Eq. (41)\], it slips into the low-energy tail of the zero-field phonon spectrum where $`\mathrm{Im}\mathrm{\Pi }\mathrm{\Pi }_0`$; hence, $`\mathrm{\Delta }\omega _p\omega _p`$. The connection between the linewidth of the absorption line in strong magnetic fields and the low-frequency phonon modes in zero magnetic field, which we just established, is the keystone of the present paper. The properties of the zero-field soft modes will be discussed in more detail in Sec. V where we argue that they give rise to the power-law dependence of $`\mathrm{Im}\mathrm{\Pi }`$ on $`ϵ`$, $$\mathrm{Im}\mathrm{\Pi }(ϵ)C_2\mathrm{\Pi }_0(ϵ/\mathrm{\Pi }_0)^s,ϵ\mathrm{\Pi }_0,$$ (69) with exponent $`s=3/2`$ and numerical prefactor $`C_21`$. Combining Eqs. (3536) and (IV69), we find for the intermediate-$`B`$ regime: $`\sigma _{xx}(\omega )`$ $`=`$ $`i{\displaystyle \frac{e^2n_e\omega }{m\omega _{p0}^2}}{\displaystyle \frac{1if_1(\omega )}{[1if_1(\omega )]^2(\omega \omega _c/\omega _{p0}^2)^2}},`$ (70) $`f_1(\omega )`$ $``$ $`\{\begin{array}{cc}C_2(\omega /\omega _{p0})^{2s},& \omega \omega _{p0},\hfill \\ \mathrm{const},& \omega _{p0}\omega \omega _c,\hfill \end{array}`$ (74) $`\omega _{p0}\omega _c\omega _{p0}\sqrt{\lambda /\mu }.`$ In this regime the absorption line narrows down according to Eq. (2), which we reproduce here for convenience: $$\mathrm{\Delta }\omega _p/\omega _p(\omega _p/\omega _c)^s.$$ (75) The high-field regime is described by Eq. (68). In this case the relative linewidth is field-independent and is of the order of $$\mathrm{\Delta }\omega _p/\omega _p(\mu /\lambda )^s$$ (76) as given by Eq. (3) and illustrated in Fig. 1. Although Eq. (68) was derived with much less rigor than Eq. (67), the saturation of $`\mathrm{\Delta }\omega _p/\omega _p`$ in strong fields is certainly to be expected. Indeed, in strong fields the dynamics is dominated by the Lorentz force; therefore, $`ϵ`$ must enter in the combination $`ϵϵ_c`$, not by itself. But then a phonon eigenstate, which for a given (large) $`ϵ_c=ϵ_c^{}=\rho (\omega _c^{})^2`$ has an energy $`ϵ_i`$, is also an eigenstate of the system for larger $`ϵ_c`$, with the eigenvalue $`ϵ_i(ϵ_c^{}/ϵ_c)`$. Hence, as the magnetic field increases, all relevant eigenfrequencies scale inversely proportional to $`\omega _c`$, while the conductivity varies in the self-similar way, $$\sigma _{xx}(\omega )=\frac{\omega _c^{}}{\omega _c}\sigma _{xx}^{}\left(\frac{\omega \omega _c}{\omega _c^{}}\right),$$ (77) and posseses constant $`\mathrm{\Delta }\omega _p/\omega _p`$. The explicit expression for $`\sigma _{xx}(\omega )`$ is similar to Eq. (74), $`\sigma _{xx}(\omega )`$ $`=`$ $`i{\displaystyle \frac{e^2n_e\omega }{m\omega _{p0}^2}}{\displaystyle \frac{1if_2(\omega )}{[1if_2(\omega )]^2(\omega \omega _c/\omega _{p0}^2)^2}},`$ (78) $`f_2(\omega )`$ $``$ $`\{\begin{array}{cc}(\omega /\mathrm{\Omega })^{2s},& \omega \mathrm{\Omega },\hfill \\ \mathrm{const},& \mathrm{\Omega }\omega \omega _c,\hfill \end{array}`$ (81) $`\mathrm{\Omega }`$ $`=`$ $`C_3{\displaystyle \frac{\omega _{p0}^2}{\omega _c}}\sqrt{\lambda /\mu },\omega _c\omega _{p0}\sqrt{\lambda /\mu },`$ (82) where $`C_31`$. Formulas (74) and (81) are our final results. They describe the entire lineshape of the pinning mode both in intermediate and strong magnetic fields. However, their derivation was presented in diagrammatic rather than physical terms. Next, we will give an alternative derivation, which elucidates the physics of the line narrowing and also helps to clarify the structure of the phonon eigenstates that compose the pinning mode. We will start with the qualitative picture of localized T-phonons developed in Sec. III and add a new ingredient, the stiff L-degree of freedom. It is clear that an admixture of the L-component produced by a joint action of the disorder and the Lorentz force makes phonon eigenstates much more extended in real space. Unlike the softer T-phonons, the stiffer L-ones cannot be confined in small areas of size $`R_c`$. In order to understand the large scale structure of the eigenstates, we can coarse-grain the system by integrating out the degrees of freedom on the spacial scales between $`a`$ and $`R_c`$. An insight on the form of the effective Hamiltonian $`𝐇_L^{eff}`$ after the coarse-graining is given by the spectral decomposition of the matrix element $$𝐫\left\|\frac{1}{𝐇_T^0+𝐒_Tϵ𝐈}\right\|𝐫^{}=\underset{i}{}\frac{u_{Ti}(𝐫)u_{Ti}^{}(𝐫^{})}{ϵ_iϵ}.$$ (83) Here $`ϵ_i`$ and $`u_{Ti}(𝐫)`$ are the eigenvalues and the eigenfunctions of localized T-phonons. Since each of $`u_{Ti}`$ is localized within an area of size $`R_c`$, we expect that after the coarse-graining the numerators $`u_{Ti}(𝐫)u_{Ti}^{}(𝐫^{})`$ transform into local operators $`R_c^2\delta (𝐫𝐫^{})\delta (𝐫𝐫_i)`$, where $`𝐫_i`$ is the center-of-gravity of the $`i`$th mode. As for the denominators, the discussion in Sec. III shows that in weak and intermediate magnetic fields \[defined by the inequality (56)\], $`ϵ_i`$’s should remain unaffected. Finally, naive averaging of $`𝐒_L`$ in Eq. (7) yields $`S_0𝐈`$ but the interaction with high-$`k`$ T-modes, which can be calculated within the SCBA-type perturbation theory, renormalizes $`S_0`$ to $`\mathrm{\Pi }_0`$. The resultant coarse-grained Hamiltonian is $$𝐇_L^{eff}=𝐇_L^0+\mathrm{\Pi }_0𝐈\underset{i}{}\frac{c_i\mathrm{\Pi }_0^2+ϵϵ_c}{ϵ_iϵ}R_c^2\delta (𝐫𝐫_i),$$ (84) where $`c_i1`$ are random and $`ϵ_i(0,2\mathrm{\Pi }_0)`$. Without losing essential physics, we can assume that $`𝐫_i`$’s form a regular square lattice with lattice constant $`R_c`$. What are the properties of the obtained lattice model? One of them is the Lorentzian tails of the distribution function $`P(U)`$ of the on-site disorder terms $`U_i(c_i\mathrm{\Pi }_0^2+ϵϵ_c)/(ϵϵ_i)`$. For example, if $`ϵϵ_p`$, then $$P(U)\frac{\mathrm{\Pi }_0^2R_c^2\nu (ϵ)}{U^2},\|U\|\mathrm{},$$ (85) where $`\nu (ϵ)`$ is the T-phonon density of states introduced in Sec. III. Clearly, the variance $`U^2`$ of the on-site disorder is unbounded. One can therefore expect much stronger disorder effects than in the case of a Gaussian distribution with the same typical values of $`U_i`$, i.e., same $`\mathrm{ln}\|U\|`$. Let $`L_L`$ be the mean-free path of L-component of the low-energy phonons. As explained in Sec. III, Ioffe-Regel criterion suggests that $`L_L`$ is simultaneously their localization length. For weak Gaussian disorder and unscreened Coulomb interactions $`L_L`$ is exponentially large, see Appendix A. Now we will show that in the presence of the long tails (85), $`L_L`$ is only as a power-law function of the disorder strength. The crucial point is that the scattering is dominated by a few largest $`U_i`$’s. As soon as we realize this, we can replace $`P(U)`$ by the Cauchy distribution $$P(U)=\frac{1}{\pi }\frac{\mathrm{\Gamma }}{(UU_0)^2+\mathrm{\Gamma }^2},\mathrm{\Gamma }(ϵ)\mathrm{\Pi }_0^2R_c^2\nu (ϵ)$$ (86) on the grounds that it has the same tails. At this point we make a reasonable assumption that correlations among $`U_i`$’s at different sites can be neglected, and arrive at the famous Lloyd model. The self-energy can now be calculated exactly, $$\mathrm{\Pi }_L(ϵ)=\mathrm{\Pi }_0U_0(ϵ)i\mathrm{\Gamma }(ϵ),$$ (87) while the mean-free path is the solution of the equation $`H_L^0(L_L^1)=\|\mathrm{Im}\mathrm{\Pi }_L\|`$, which gives $`L_LY/\mathrm{\Gamma }`$. To compare with our earlier results we need to know $`\nu `$ and $`U_0`$. The integral in Eq. (51) is determined by $`\|𝐤\|q_c`$; therefore, $$\nu (ϵ)q_c^2\frac{\mathrm{Im}\mathrm{\Pi }^{}(ϵ)}{\mathrm{\Pi }_0^2}\mu ^1\left(\frac{ϵ}{\mathrm{\Pi }_0}\right)^s,ϵ\mathrm{\Pi }_0.$$ (88) The estimate of $`U_0`$ is $`U_0ϵϵ_c/\mathrm{\Pi }_0`$. Combining these together, we find that Eq. (87) is in agreement with Eq. (62), which validates our mapping onto the Lloyd model. To describe the dominant scattering mechanism in the Lloyd model notice that $`\mathrm{Im}\mathrm{\Pi }_L`$ comes from the inhomogeneous broadening of the phonon energies, i.e., the energy shifts $`\delta ϵ`$ caused by on-site disorder $`U_i`$. For each localized L-phonon the largest energy shift comes from a site with largest $`U_i`$ within the localization area $`L_L\times L_L`$, which is typically $`U_i(L_L/R_c)^2\mathrm{\Gamma }`$ for the Cauchy distribution (86). Since the amplitude of the L-phonon wavefunction at the position of the oscillator is $`u_L1/L_L`$, the resultant energy shift is $`\delta ϵU_iu_L^2R_c^2\mathrm{\Gamma }`$. The energy shifts from other sites are subdominant; thus, $`\mathrm{Im}\mathrm{\Pi }_L\mathrm{\Gamma }`$, in agreement with the exact result, Eq. (87). In the original formulation large $`U_i`$ come from the localized T-oscillators whose energy $`ϵ_i`$ is almost in resonance with $`ϵ`$, i.e., $`\|ϵ_iϵ\|1/\nu (ϵ)L_L^2`$. A few of such resonant scatterers will localize the L-component on the scale of the mean-free path $`L_L`$. The proposed picture resembles the situation in a crystal pinned by strong dilute impurities, the role of impurities played by the resonant sites. The structure of a typical phonon is illustrated in Fig. 7. It involves of the order of $`(L_L/R_c)^2`$ T-modes of Fig. 5 oscillating with the same frequency imposed by the mutual coupling mediated by the L-phonon. The oscillation amplitude is roughly the same for all the T-modes except for a few resonant ones, where it is much larger. These resonant sites “drain” the energy of the L-degree of freedom thereby localizing it. Let us now describe the evolution of the phonon eigenstates as the magnetic field increases. Since the pinning mode frequency goes down, lower and lower frequency resonant sites are required to effectively scatter and localize the mixed L-T phonon vibrations. Such sites are therefore the soft modes, the oscillators of unusually low frequency, which shape up the tails of the zero-field phonon spectrum (Lifshitz tails). The soft modes appear due to very rare, peculiar disorder configurations. The lower the frequency, the mode dilute in real space they are. Correspondingly, the mean-free path (localization length) of the magnetophonons becomes larger and larger. The narrowing of the absorption line is then similar to the motional narrowing phenomenon. Eventually, in very strong magnetic fields, the enhancement of the L-T mixing by the Lorentz force start to diminish the frequencies of the localized T-phonons. At this point slipping of the pinning frequency $`\omega _p`$ deeper into the soft-mode tail stops. The wavefunctions of the phonons cease to change, only their eigenfrequencies continue decreasing in inverse proportion to the magnetic field. As a result, $`\mathrm{\Delta }\omega _p/\omega _p`$ remains constant. The physical structure of the pinned phonon modes being clarified, the only important questions that remain to be considered are the nature of the soft modes and the derivation of Eq. (69). This will be the subject of the next section. ## V Soft modes It is clear that the soft phonon modes must come from rare places where pinning is unusually weak. The actual calculation of the density of states $`\nu (ϵ)`$ in the low-$`ϵ`$ tail is however nontrivial. The method most suitable for the task seems to be the method of optimal fluctuation. It has been successfully employed for calculation of the Lifshitz tails in other disordered systems, where it proves to be asymptotically exact. The method is based on the following idea. The energy $`ϵ`$ of a given eigenstate is determined in a complicated way by the distribution of the random potential within the entire area supporting the eigenstate and so there are many different random potential configurations, which give the same eigenenergy. If $`ϵ`$ is small, any such configuration is untypical, i.e., a fluctuation of some sort. The method is based on the assumption that certain fluctuations have much higher probability than the rest and dominate the quantity of interest, $`\nu (ϵ)`$ in this case. The objective is then to design such an optimal fluctuation and evaluate its probability. Let us mention one type of fluctuation, which is not optimal. This is a configuration where the random potential is strongly suppressed in a large area $`L\times L`$, where $`L\sqrt{\mu /ϵ}`$. (It may be useful to recall at this point that we are investigating the soft modes of a system where L-phonons are integrated out). The probability of such a fluctuation can be estimated by mentally dividing the area into uncorrelated blocks of size $`R_c\times R_c`$ and multiplying together the probabilities $`ϵ/\mathrm{\Pi }_0`$ that the random potential is suppressed within each block. The total probability is exponentially small $`\mathrm{exp}[(\mathrm{\Pi }_0/ϵ)\mathrm{ln}(\mathrm{\Pi }_0/ϵ)]`$. This is a general situation for fluctuations spread over a large area in real space. Thus, the optimal fluctuation must be of the smallest possible size. It is easy to see that this size is $`R_c`$. Indeed, to get a small phonon eigenenergy $`ϵ`$ in a small volume, the positive “kinetic energy” $`\mu L^2`$ must be accurately compensated by the negative “potential energy” $`\delta 𝐒_T_L`$. The fluctuations of the latter averaged over the area $`L\times L`$ are of the order of $`\mathrm{\Pi }_0(R_c/L)`$. Thus, for $`LR_c`$ the cancellation of $`\mu L^2`$ can be done by a typical fluctuation of $`\delta 𝐒_T`$ and the probability of such an event would not have any exponential suppression. Instead, we expect the power-law dependence of $`\nu (ϵ)`$ on $`ϵ`$, Eq. (88). In order to possess a soft mode, the system must have a soft direction, i.e., a certain collective coordinate $`X`$ such that the total energy of the system $`E`$ as a function of $`X`$ has a very shallow minimum. Let $`X=0`$ be the ground state. In its vicinity $`E`$ should be Taylor-expandable, $$E(X)=\frac{\alpha }{2}X^2+\beta X^3+\gamma X^4+\mathrm{}$$ (89) The coefficient $`\alpha `$ has the meaning of an effective spring constant of the local oscillator, while the phonon energy $`ϵ`$ is determined by the ratio of the spring constant and the effective mass. The latter is proportional to the area involved in the oscillations. As we argued above, this area is of the order of $`R_c^2`$ for all $`ϵ<ϵ_p`$. Therefore, $`\alpha `$ should scale linearly with $`ϵ`$ in the limit $`ϵ0`$, $$\alpha =C_4ϵ.$$ (90) If we find a way to estimate the probability of such untypically small $`\alpha `$, we will succeed in calculating the exponent $`s`$ of the power-law function $`\nu (ϵ)`$. This kind of calculation is hardly possible without specifying $`X`$. Actually, ascribing the precise meaning to $`X`$ amounts more or less to designing the optimal fluctuation. At the moment we do not have a good understanding how to do it in the 2D case. However, in one dimension the general structure of the calculation is clear. We are guided by the earlier work on the subject by Feigelman et al. and by Aleiner and Ruzin. In the case of a weakly pinned 1D elastic chain, the role of the desired collective coordinate $`X`$ can be played by the elastic displacement field $`u`$ at an arbitrary point in the middle of the chain. This point divides the system into two parts, which communicate only through the single variable $`u`$. For any fixed $`u`$ serving as a boundary condition, we can find separately the ground state energies $`E_<(u)`$ and $`E_>(u)`$ of the left and the right halves of the chain. The ground state of the whole system is determined by minimizing the sum $`E(u)=E_<(u)+E_>(u)`$ with respect to $`u`$. Both $`E_<(u)`$ and $`E_>(u)`$ are periodic function of $`u`$ with the period equal to the lattice constant $`a`$. On their period they have several (typically, of the order of $`a/\xi `$) minima and maxima. The beautiful idea of Aleiner and Ruzin (totally overlooked in Ref. ) was that the most efficient way to generate a soft mode is via a frustration, when a maximum of $`E_<(u)`$ occurs near a minimun of $`E_>(u)`$. In this case the second derivative of $`E`$, i.e., $`\alpha =E^{\prime \prime }`$, can be very small even though each of the second derivatives of $`E_<(u)`$ and $`E_>(u)`$ are typical. The next step is to show that the probability density distribution $`P(\alpha )`$ of $`\alpha `$ at local minima of $`E`$ vanishes linearly with $`\alpha `$ in the limit $`\alpha 0`$. Indeed, $`P(\alpha )`$ $`=`$ $`\delta (E^{\prime \prime }\alpha )\|_{E^{}=0}`$ (91) $`=`$ $`{\displaystyle \frac{1}{a}}{\displaystyle \underset{0}{\overset{a}{}}}𝑑u\delta (E^{\prime \prime }\alpha )\delta (E^{})\left\|{\displaystyle \frac{E^{}}{u}}\right\|`$ (92) $``$ $`\|\alpha \|\delta (E^{\prime \prime })\delta (E^{})`$ (93) $`=`$ $`\|\alpha \|\delta (E_<^{\prime \prime }+E_>^{\prime \prime })\delta (E_<^{}+E_>^{}).`$ (94) The last equation explicitly demonstrates the aforementioned cancellations between the derivatives of $`E_<`$ and $`E_>`$. The value of the delta-function product average is determined by properties of typical configurations. Hence, the low probability of having unusually shallow local minima is entirely due to the first factor on the left-hand side of Eq. (94) and so $`P(\alpha )\|\alpha \|`$. In particular, $`Pϵ`$ in the case of interest (90). But this is not yet the end of the story. Following Aleiner and Ruzin we argue that the probability of having a shallow global minimum is additionally suppressed. Indeed, if the coefficient $`\beta `$ in front of the the cubic term in Eq. (89) is too large, $`E(X)`$ would have a second minimum which is deeper than that at the ground state $`X=0`$. To avoid that $`\|\beta \|`$ must not exceed $`\gamma \sqrt{\alpha }\sqrt{ϵ}`$. The typical dependence of the resultant $`E(X)`$ is illustrated in Fig. 8. The total probability density of the optimal fluctuation is therefore proportional to $`ϵ\sqrt{ϵ}=ϵ^s`$ with $`s=3/2`$, as we claimed. It is very likely that $`s=3/2`$ is in fact the general result independent of the number of dimensions because the optimal fluctuation would always have only one soft direction, i.e., the system is essentially one-dimensional. Finally, let us compare our results with those in the literature. One group of works, by Feigelman et al. has been already mentioned before. We borrowed from them the idea of splitting the system into two statistically independent parts. The major mistake of these authors is overlooking the possibility of frustrations in the system. In other words, they did not realize that it is not necessary to require that both $`E_<`$ and $`E_>`$ have a small second derivative with respect to $`u`$ when only their sum $`E_<+E_>`$ needs to be so. Another point of disagreement between us and them is the size of the localized phonons. Feigelman et al. take for granted that it is of the order of $`\sqrt{\mu /ϵ}`$, while we gave an argument that it should be $`R_c`$, i.e., much smaller. In another large group of papers, Refs. , $`\nu _\omega (\omega )\omega ^2`$ for the density of states, and a similar dependence, $`\mathrm{Re}\sigma _{xx}(\omega )\omega ^2`$, for the conductivity were calculated. To facilitate the comparison we should point out that the conventionally defined phonon density of states $`\nu _\omega (\omega )`$ is related to our $`\nu (ϵ)`$ by $$\nu _\omega (\omega )=2\nu (\rho \omega ^2)\rho \omega .$$ (95) Hence, our result is $`\mathrm{Re}\sigma _{xx}(\omega )\nu _\omega (\omega )\omega ^4`$, same as in Refs. (see also Ref. ). All the papers in the second group are explicitly or implicitly based on the SCBA (see Appendix B for more details). The second power of $`\omega `$ can be traced to the unphysical square-root singularity $`\nu (ϵ)\|ϵϵ_{\mathrm{th}}\|^{1/2}`$ near a band edge $`ϵ_{\mathrm{th}}`$ ($`ϵ_{\mathrm{th}}=0`$ in our case), which is a well-known basic flaw of the SCBA. Note however that even within the SCBA $`\nu (ϵ)`$ conforms to the generic form (88), except it corresponds to $`s=1/2`$ instead of what we believe is the correct result, $`s=3/2`$. ## VI Quantum and thermal effects So far we have neglected any effects of quantum nature or due to a finite temperature. Some of them will be addressed in this section but their systematic, in-depth treatment is deferred for future work. Let us begin with noting that the electrons of the pinned WC constantly fluctuate around their equilibrium positions. The typical size of such fluctuations is easy to find. For example, in strong magnetic fields and at low temperatures ($`k_BT\mathrm{}\omega _c`$) where all the electrons are confined to the lowest Landau level, we have $`𝐮^2`$ $`=`$ $`l_B^2+{\displaystyle \frac{\mathrm{}}{\rho \omega _c}}{\displaystyle \underset{𝐪}{}}\mathrm{coth}{\displaystyle \frac{\mathrm{}\mathrm{\Omega }(𝐪)}{2k_BT}}`$ (97) $`\times {\displaystyle \frac{1}{\rho \omega _c\mathrm{\Omega }(𝐪)}}\left[\mathrm{\Pi }_0+{\displaystyle \frac{\mu (𝐪)+\lambda (𝐪)}{2}}q^2\right],`$ $`\mathrm{\Omega }(𝐪)`$ $`=`$ $`(\rho \omega _c)^1\sqrt{[\mathrm{\Pi }_0+\mu (𝐪)q^2)][\mathrm{\Pi }_0+\lambda (𝐪)q^2]}.`$ (98) This formula can be used away from the thermal or quantum melting transitions, where the phenomenological criterion $`𝐮^2^{1/2}<\epsilon a`$ is satisfied. Here $`\epsilon 0.2`$ is the Lindenmann parameter. Incidentally, at $`T=\mathrm{\Pi }_0=0`$, Eq. (97) gives the variance of the electron fluctuations in the correlated WC of Lam and Girvin. Let us focus on the range of temperatures $`k_BT\mathrm{}\omega _B`$ where $`\omega _B4\pi \mu /m_e\omega _c`$ is the magnetophonon bandwidth (usually, $`\mathrm{}\omega _B/k_B2\mathrm{K}`$). At such temperatures $`𝐮^2^{1/2}`$ is totally dominated by quantum fluctuations and is of the order of $`l_B`$. We expect that as long as $`𝐮^2^{1/2}`$ is smaller than the correlation length of the pinning potential $`\xi `$, which is the smallest length scale in the classical theory, all the results obtained in the previous sections acquire at most minor corrections. Recently Fertig and Chitra et al. studied the opposite limit, $`\xi l_B`$, and found a novel dependence of the pinning frequency $`\omega _p`$ on the magnetic field, although they disagree with each other on the functional form of such a dependence. To sort things out we offer a different perspective on this question. Recall that the classical theory predicts the $`1/B`$-behavior, $$\omega _p\frac{C(0)c}{e\mu \xi ^4B},$$ (99) which follows from Eqs. (14), (36), and (40). We propose that a reasonably accurate estimate for $`\omega _p`$ in the regime $`\xi l_B`$ can be obtained within a quasiclassical approximation. The electrons are visualized as compact wavepackets whose centers of gravity perform classical motion, while the quantum effects are presumed to be contained in the form-factor of the wavepakets. This type of quantum effects amount to replacing the bare pinning potential $`U(𝐫)`$ by its convolution with the form-factor $`F_e(𝐫)`$ of the wavepackets. For small fluctuations $`l_Ba`$ the appropriate form-factor is Gaussian, $`F_e(𝐫)=(\pi 𝐮^2)^1\mathrm{exp}(r^2/𝐮^2)`$. Upon the convolution, we obtain an effective random potential, with the correlation length $`l_B`$ and variance $`C(0)\xi ^2/l_B^2`$. In view of such modifications, Eq. (99) transforms into $$\omega _p\frac{e^2}{\mu \mathrm{}^3c^2}C(0)\xi ^2B^2,\xi <l_B.$$ (100) The quasiclassical approximation is certainly not exact; nonetheless, when used away from the quantum Hall fractions, at sufficiently low $`T`$, and for robust quantities such as $`\omega _p`$, it should be qualitatively correct. The linear in $`B`$-dependence of $`\omega _p`$ derived by Fertig \[see Eq. (44) in Ref. \] is at odds with the $`B^2`$-dependence found above \[Eq. (100)\] and in Ref. . Although Fertig consider a Poissonian disorder (potential wells or “pits” of size $`s<l_B`$, energy $`\mathrm{\Delta }V`$, and the areal density $`n_i<1/\pi l_B^2`$) rather than the Gaussian one, the real source of the discrepancy is his ascertion that the elastic distortion at the Larkin length scale is of the order of $`n_i^{1/2}`$. Instead, $`l_B`$ should be used. Upon this correction, the formula for $`\omega _p`$ acquires the above form (100) with $`C(0)\xi ^2`$ replaced by $`n_is^4\mathrm{\Delta }V^2`$. The finite amplitude of the quantum fluctuations has another important repercussion: nonvanishing high-order in $`𝐮`$ terms in the pinning energy. As a result, the absorption line acquires extra broadening even at $`T=0`$. It originates from finite widths $`\mathrm{\Gamma }_n`$ of the energy levels $`n=1,2,\mathrm{}`$ of the localized magnetophonon oscillators (the ground state $`n=0`$ remains unbroadened). At zero $`T`$ only the transitions between $`n=0`$ and $`n=1`$ levels contribute to the absorption. The level width $`\mathrm{\Gamma }_1`$ of the $`n=1`$ state is determined by the decay of the given magnetophonon into two other magnetophonon excitations of lower frequencies, $`\mathrm{}\omega _1\mathrm{}\omega _2+\mathrm{}\omega _3`$ (see Fig. 9). Such a process is always possible in a truly random system, the necessary three-phonon matrix elements, $`M_{\alpha \beta \gamma }(𝐫)=n_e{\displaystyle \underset{𝐊}{}}K_\alpha K_\beta K_\gamma U(𝐫)\mathrm{sin}𝐊[𝐫𝐮^{(0)}(𝐫)],`$ being provided by the cubic anharmonicities. The corresponding contribution $`\mathrm{Im}\mathrm{\Pi }_3`$ to the imaginary part of the self-energy can be obtained by evaluating the diagram in Fig. 9, $`\mathrm{Im}\mathrm{\Pi }_3(𝐪,\omega )`$ $`=`$ $`{\displaystyle \frac{2\mathrm{}}{L^2}}{\displaystyle \underset{𝐤}{}}{\displaystyle \underset{𝐤^{}}{}}\stackrel{~}{M}_{\alpha \mu \sigma }\stackrel{~}{M}_{\beta \nu \tau }{\displaystyle \frac{d\omega ^{}}{2\pi }}`$ (101) $`\times `$ $`[n_B(\omega ^{})+n_B(\omega \omega ^{})+1]`$ (102) $`\times `$ $`\mathrm{Im}D_{\mu \nu }(𝐤,\omega ^{})\mathrm{Im}D_{\sigma \tau }(𝐤^{},\omega \omega ^{}),`$ (103) $`n_B(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{exp}(\mathrm{}\omega /k_BT)1}}.`$ (104) At $`k_BT\mathrm{}\omega _p`$ only the frequency interval $`0<\omega ^{}<\omega `$ contributes to the integral, which leads us to the estimate $$\mathrm{Im}\mathrm{\Pi }_3(\omega _p)\mathrm{Im}\mathrm{\Pi }(\omega _p/2)\frac{l_B^2}{\xi ^2}\frac{a^2}{R_c^2}\frac{\|\mathrm{Im}\mathrm{\Pi }(\omega _p/2)\|}{\mathrm{\Pi }_0}.$$ (105) This formula is derived assuming that $`\xi l_B`$; otherwise, in the spirit of the quasiclassical approximation, we have to replace $`\xi `$ by $`l_B`$. Even in that case $`\mathrm{Im}\mathrm{\Pi }_3(\omega _p)`$ is smaller than the previously found $`\mathrm{Im}\mathrm{\Pi }(\omega _p)`$ by a large factor $`(R_c/a)^2(\lambda /\mu )^s`$. At higher temperatures, the anharmonic effects are enhanced by a factor of $`k_BT/\mathrm{}\omega _p`$ originating from the Bose distribution functions $`n_B`$ in Eq. (103). They should start to affect the linewidth when $`k_BT(R_c/a)^{7/2}\mathrm{}\omega _p`$ (cf. Ref. ). A word of caution is in order. The diagram in Fig. 9 does not contain vertex corrections, which contain information about, e.g., the soft modes. The soft modes are crucial for the dynamical response in strong magnetic fields and at the same time are potential sources of much larger anharmonicities. Thus, the response of a weakly pinned quantum WC may prove to be more nontrivial. Concluding this section, we would like to re-emphasize that for the problem studied in the main part of the paper (zero-temperature linear-response of a classical WC), the anharmonicities are not important. ## VII Comparison with experiment In this section we will attempt to analyze the most recent experimental data in the light of our understanding of the pinning mode. Let us first discuss the position of the pinning resonance. From Eqs. (99) and (100) we see that depending on the type of disorder, $`\omega _p`$ can either decrease as $`1/B`$ or increase as $`B^2`$ when the magnetic field increases. Weak magnetic field dependence of $`\omega _p`$ found in the experiments leads us to conclude that $`\xi `$ remains of the order of $`l_B`$ in the limited range $`10`$$`15\mathrm{T}`$ of the magnetic fields where the pinning mode was detected, i.e., $`\xi l_B75\mathrm{\AA }`$. The source of such a short-range random potential is likely to be the roughness of the heterostructure interface as suggested by Fertig. The experimental value of $`\omega _p8\times 10^9s^1`$ for $`n_e=5.4\times 10^{11}\mathrm{cm}^2`$ can then be reproduced with reasonable values of fitting parameters (see Sec. VI) $`s=30\mathrm{\AA }`$, $`\mathrm{\Delta }V=4.3\mathrm{K}`$, and $`n_i=4\times 10^{11}\mathrm{cm}^2`$. Unfortunately, the disorder parameters are poorly known, and so we have to use $`\omega _p`$ to estimate the disorder characteristics, not the other way around. For example, we can solve Eq. (100) for the rms amplitude $`\sqrt{C(0)}`$ of the pinning potential, which gives a value of the order $`0.2\mathrm{K}`$. Let us now discuss the density dependence of $`\omega _p`$. Explicitly, it enters Eq. (100) only through the shear modulus $`\mu (n_e)`$. The latter should behave as $`\mu n_e^{3/2}`$ (see Sec. II) away from the quantum Hall fractions, e.g., the $`\frac{1}{3}`$-filling. If $`C(0)`$ is $`n_e`$-independent, we therefore expect $`\omega _pn_e^{3/2}`$. Such a dependence was indeed observed by Li et al. for the concentration of holes ($`p`$-type samples were used) in the range $`3`$$`5\times 10^{10}\mathrm{cm}^2`$, which corresponds to the filling factor range $`0.1`$$`0.16`$. At lower concentrations the pinning frequency continued growing as $`n_e`$ decreased but less rapidly. This comes presumably from the fact that the pinning is no longer weak at such low $`n_e`$: one can verify that $`R_c`$ is approximately $`6a`$ at the highest densities, but approaches $`a`$ (the lattice constant) at the lowest densities used in Ref. . (In Ref. the density dependence of $`\omega _p`$ was not investigated). The relationships among the empirical values of $`\omega _p`$ and $`\mathrm{\Delta }\omega _p`$ and the experimental parameters ($`n_e`$ and $`B`$) provide another means to test the agreement between the theory and the experiment. To this end we rewrite Eq. (3) (with $`s=3/2`$) in the following way: $$Q\frac{\omega _p}{\mathrm{\Delta }\omega _p}\left[\left(26\frac{ec}{\kappa }\frac{n_e^{3/2}}{Bf_{pk}}\right)^{1/2}5\right]^{3/2},$$ (106) where $`f_{pk}=\omega _p/(2\pi )`$ is the pinning frequency in cycles per second. The largest quality factor $`Q8`$ reported by Li et al. was achieved for $`n_e=5.4\times 10^{10}\mathrm{cm}^2`$, $`B=13\mathrm{T}`$, where $`f_{pk}`$ was measured to be $`1.4\mathrm{GHz}`$. In Ref. $`Q45`$ was found for comparable $`n_e`$, $`B`$, and $`f_{pk}`$. Our theoretical estimate from Eq. (106) is $`Q260`$ and far exceeds both. We should point out, however, that specific details of the experimental setup become important for observability of such a high $`Q`$. Henceforth we focus on the work of Li et al. since most of their data is available in the published record. Because of the finite distance ($`30\mu \mathrm{m}`$) between the plates of the coplanar microwave waveguide (see Fig. 10), these were not the ideal $`𝐪=0`$ measurements. We estimate that the spread of $`𝐪`$’s imposes un upper bound of about $`30`$ for the largest observable $`Q`$. The remaining discrepancy seems to be rooted in finite-temperature and to a lesser degree finite incident power effects ignored in the derivation of Eq. (106). In the experiments, the linewidth was decreasing roughly linearly as the temperature varied from $`T=200\mathrm{mK}`$ to $`T=50\mathrm{mK}`$ with other parameters held fixed. Below $`50\mathrm{mK}`$ the linewidth appeared to reach saturation, but only at the nominal level of incident microwave power ($`80\mathrm{pW}`$). When lower power levels were used, the line continued sharpening up. At the lowest experimental temperature of $`25\mathrm{mK}`$, the linewidth did not show any signs of saturation (this time as a function of power) even upon ten-fold input power reduction, which was near the sensitivity limit. Several sources of the thermal broadening can be envisioned: (a) thermally excited single-particle excitations of the WC, (b) phonons of the host semiconductor, (c) the anharmonisms of the collective modes of the pinned WC, and perhaps some others. The classical activation energy of vacancies and interstitials is of the order of $`0.5\mathrm{K}`$ at the densities studied, and mechanism (a) should be frozen out at $`25\mathrm{mK}`$. Furthermore, for the WC of low density the bandwidth of such excitations is very narrow and so they are easily localized by the random potential. The weak electron-phonon coupling in GaAs and small phonon phase space at the frequencies involved ($`\omega _p`$) are likely to render the mechanism (b) inefficient. The preliminary estimate of the broadening due to mechanism (c) was given in the previous section. It is too small to explain the experimental observations. We speculate that stronger anharmonic effects and a better agreement with the experiments may be obtained if one properly accounts for untypically large anharmonicities at the locations of the soft modes, which as we showed in this paper, play a very important role in the response. It is also quite possible that the actual pinning potential is more complicated than the one studied here. In reality, it can be due to a combination of the interface roughness and dilute residual ions in the vicinity of the WC plane. In this case $`\omega _p`$ may be determined by the former, while $`Q`$ could be limited by the latter. Such more complicated models as well as the dynamic response of a pinned quantum WC are interesting subjects awaiting further investigation. It would be also interesting to investigate if the discussed phenomena appear in the conventional charge-density waves (see Sec. I). These materials are three-dimensional, and some important modifications of the present theory may be needed. There are other complications such as an extra dissipation due to uncondensed quasiparticles but those can be suppressed by lowering the temperature. In any case, it would be remarkable if narrow pinning modes could be produced in these materials simply by applying a sufficiently strong magnetic field. Our theory should literally apply to the pinning mode of a WC formed by electrons on solid hydrogen (see the second book cited under Ref. ). We hope that this paper will encourage further experiments on these or other numerous systems where the pinning manifests itself. ###### Acknowledgements. This research is supported by US Department of Energy Grant No. DE-FG02-90ER40542 and NSF Grant No. DMR-9802468. We thank R. Chitra, Mark Dykman, Lloyd Engel, Herb Fertig, Gabi Kotliar, Chi-Chun Li, Leonid Pryadko, Misha Raikh, Dan Tsui, and Valerii Vinokur for useful discussions. M. M. F. is grateful to Tito Williams for providng a copy of Ref. . ## A Model calculation for L-L scattering In order to evaluate the importance of the L-L scattering channel, let us consider the model action $`A_L(ϵ)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐪}{}}u_L^{}(𝐪)[H_L^0(𝐪)ϵ]u_L(𝐪)`$ (A1) $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐫}{}}u_L^{}(𝐫)S_L(𝐫)u_L(𝐫).`$ (A2) where $`S_L`$ is a local operator of a Gaussian white-noise type with parameters $`S_L=\mathrm{\Pi }_0,S_L(𝐫_1)S_L(𝐫_2)=\mathrm{\Pi }_0^2V^{\prime \prime }(0)\delta (𝐫_1𝐫_2).`$ Our objective is to calculate the disorder-averaged Green’s function $`D_L(𝐪,ϵ){\displaystyle \frac{i}{L^2}}{\displaystyle \frac{𝒟u_L𝒟u_L^{}u_L(𝐪)u_L^{}(𝐪)e^{iA_L(ϵ)}}{𝒟u_L𝒟u_L^{}e^{iA_L(ϵ)}}},`$ where an infinitesimal imaginary correction is assumed to be included via $`ϵϵ+i0`$ to make the integrals convergent. Using the “supersymmetry” technique, $`D_L`$ can be represented in the form $`D_L={\displaystyle \frac{i}{L^2}}{\displaystyle 𝒟\varphi 𝒟\varphi ^{}u_Lu_L^{}e^{iA_s[\varphi ]}}.`$ where $`\mathit{\varphi }_i^{}=[u_L^{}v_L^{}]`$ is a supervector, $`v_L(𝐪,\omega )`$ is an auxillary fermionic field, and $`A_s[\mathit{\varphi }]`$ is the sypersymmetric Eucledian action $`A_s`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐪}{}}\mathit{\varphi }^{}(𝐪)(Yq+\mathrm{\Pi }_0ϵ)\mathit{\varphi }(𝐪)`$ (A3) $`+`$ $`{\displaystyle \frac{i}{8}}V^{\prime \prime }(0){\displaystyle \underset{𝐫}{}}[\mathit{\varphi }^{}(𝐫)\mathit{\varphi }(𝐫)]^2.`$ (A4) Note that the coefficient in front of the quartic term is imaginary, which makes the action non-Hermitian. The sign of the imaginary part is such that excitations above the vacuum state $`\mathit{\varphi }0`$ decay after travelling a certain distance. This distance is simply the phonon mean free path. It is easy to establish some of the properties of the propagator $`D_L`$ right away. Since the field $`\mathit{\varphi }`$ is massive for $`ϵ\mathrm{\Pi }_0`$ at the tree level, the imaginary part of the propagator $`D_L`$ is small at such $`ϵ`$. However, $`\mathrm{Im}D_L`$ should grow rapidly as soon as $`ϵ`$ crosses the $`ϵ=\mathrm{\Pi }_0`$ threshold. Below we are going to verify this by explicit calculations. Although our methods do not work in the immediate vicinity of $`ϵ=\mathrm{\Pi }_0`$, the expressions obtained for $`ϵ<\mathrm{\Pi }_0`$ and $`ϵ>\mathrm{\Pi }_0`$ can be smoothly matched onto each other. Rescaling the fields in Eq. (A4) by $`\sqrt{2/Y}`$, we obtain the action with the Lagrangian $$=\mathit{\varphi }^{}\widehat{K}\mathit{\varphi }+\frac{M}{2}\mathit{\varphi }^{}\mathit{\varphi }i\frac{g}{4}(\mathit{\varphi }^{}\mathit{\varphi })^2,$$ (A5) where $`\widehat{K}`$ is the operator, which corresponds to multiplication by $`\|𝐪\|`$ in the momentum representation, $`M=2(\mathrm{\Pi }_0ϵ)/Y`$ plays the role of mass, and $`g2V^{\prime \prime }(0)/Y^2`$ is the bare coupling constant. The coupling constant is dimensionless, and it is possible to show that the field theory is renormalizable. Let us derive the Wilson-style renormalization group (RG) equations for the renormalized parameters $`g_R`$ and $`M_R`$. As usual, it is done by successive integrations over narrow momentum shells $`\mathrm{\Lambda }(l+\delta l)<q<\mathrm{\Lambda }(l)`$ where $`\mathrm{\Lambda }(l)=\mathrm{\Lambda }_ce^l`$ and $`\mathrm{\Lambda }_c1/a`$ are the running and the bare high-momentum cutoffs, respectively. Consider the effect of such an integration on $`g_R`$ first. Initially, $`g_R`$ remains small and we need to take into account only the three lowest order in $`g_R`$ diagrams (see Fig. 11a) to find $`g_R(l+\delta l)=g_R(l)+{\displaystyle \frac{2}{3}}g_R^2{\displaystyle \underset{e^{\delta l}\mathrm{\Lambda }<q<\mathrm{\Lambda }}{}}{\displaystyle \frac{d^2q}{(2\pi )^2}}{\displaystyle \frac{1}{q+M_R}}`$ (A6) $`\times \left[{\displaystyle \frac{1}{\|𝐪+𝐬\|+M_R}}+(𝐬𝐭)+(𝐬𝐮)\right],`$ (A7) where $`S=𝐬^2=(𝐪_1𝐪_2)^2`$, $`T=𝐭^2=(𝐪_1𝐪_3)^2`$, and $`U=𝐮^2=(𝐪_1𝐪_4)^2`$ are the Mandelstam variables with $`𝐪_i`$, $`i=1,\mathrm{},4`$ being the incoming momenta. Similarly, the tadpole diagram shown in Fig. 11b determines the variation of the renormalized mass $`M_R`$, $$M_R(l+\delta l)M_R(l)\frac{g_R}{2}\underset{e^{\delta l}\mathrm{\Lambda }<q<\mathrm{\Lambda }}{}\frac{d^2q}{(2\pi )^2}\frac{1}{\|𝐪\|+M_R}.$$ (A8) In addition to $`g_R`$ and $`M_R`$, the kinetic term $`\widehat{K}`$ also gets renormalized. We will neglect such an effect because it is of higher (second) order in $`g_R`$. Consequently, the spectral function $`\mathrm{Im}D_L`$ is determined just by the mass, $$\mathrm{Im}D_L(q,ϵ)=\frac{2}{Y}\frac{\mathrm{Im}M_R}{(q+\mathrm{Re}M_R)^2+(\mathrm{Im}M_R)^2}.$$ (A9) In this equation $`M_R=M_R(q,ϵ)`$ stands for the renormalized mass at such large negative $`l`$ where the RG flow eventually stops either because $`q\mathrm{\Lambda }(l)`$ or because $`\|M_R\|\mathrm{\Lambda }(l)`$. Below we consider exclusively $`q=0`$ case. From Eq. (A7) we find that in the limit $`\|𝐪_i\|,\|M_R\|\mathrm{\Lambda }`$ function $`\beta g_R/l`$ depends only on $`g_R`$, $$\beta =\frac{1}{\pi }g_R^2,$$ (A10) and so the RG flow equation is easy to solve: $$g_R=\frac{\pi }{l_Ll},l_L\frac{\pi }{g}.$$ (A11) At $`ll_L\pi `$, $`g_R`$ becomes of the order of one, and the theory enters the strong coupling regime. The corresponding spatial scale is $`L_{LL}\mathrm{\Lambda }_c^1\mathrm{exp}(\pi /g)`$. Since the strength of coupling maps to the strength of the effective disorder in the original formulation, $`L_{LL}`$ has the meaning of the L-phonon localization length (due to L-L scattering). Note that $`g_R`$ seems to exhibit a divergence (Landau pole) at $`l=l_L`$. This is an artifact of the one-loop RG. We expect that instead $`g_R`$ flows towards its fixed point value of the order of one. The other scaling equation, for $`M_R`$, can be deduced from Eq. (A8), $$\gamma _M\frac{M_R}{l}=\frac{g_R(l)}{4\pi }\frac{\mathrm{\Lambda }^2(l)}{\mathrm{\Lambda }(l)+M_Ri\delta }.$$ (A12) The solution cannot be not expressed in elementary functions but we can describe its asymptotics. 1. $`(g/4\pi )\mathrm{\Lambda }_c=M_{}\|M\|\mathrm{\Lambda }_c`$. Here the renormalization of the coupling constant is not important because the RG flow is effectively cut off at $`l\mathrm{ln}(\mathrm{\Lambda }_c/\|M_R\|)`$ before $`g_R`$ manages to get large. Equation (A12) reduces to the SCBA equation $$M_R=M\frac{g}{4\pi }\underset{\mathrm{\Lambda }}{\overset{\mathrm{\Lambda }_c}{}}\frac{dqq}{q+M_Ri\delta },$$ (A13) with the approximate solution $`\mathrm{Re}M_R`$ $``$ $`MM_{},`$ (A15) $`\mathrm{Im}M_R`$ $``$ $`(g/4)\mathrm{\Theta }(\mathrm{Re}M_R)\mathrm{Re}M_R,`$ (A16) where $`\mathrm{\Theta }(x)`$ is the Heaviside step-function. A word of caution is in order here. The imaginary part of $`M_R`$ vanishes at positive $`\mathrm{Re}M_R`$ only at the level of the one-loop RG. Nonperturbative methods indicate that $`\mathrm{Im}M_R`$ is nonzero albeit exponentially small (see below). 2. $`L_{LL}^1\|MM_{}\|M_{}`$. Now instead of Eq. (A13) we have $`\mathrm{Im}M_R{\displaystyle \frac{1}{4}}{\displaystyle \frac{\pi }{l_Ll_{}}}\mathrm{\Theta }(\mathrm{Re}M_R)\mathrm{Re}M_R,`$ (A18) $`\mathrm{Re}M_R\left({\displaystyle \frac{l_L}{l_Ll_{}}}\right)^{1/4}(MM_{}),`$ (A19) $`l_{}=\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }_c}{\mathrm{Re}M_R}}\right).`$ (A20) The above comment about small but nonzero $`\mathrm{Im}M_R`$ at $`\mathrm{Re}M_R>0`$ applies here as well. 3. $`\|MM_{}\|L_{LL}^1`$. Since the one-loop RG cannot be trusted beyong the point $`l=l_L\pi `$, we choose to terminate the RG procedure at this stage. Unfortunately, this leaves us with a theory which (a) does not have any small parameters and (b) includes not only the original quartic but also higher order interaction terms generated by the RG. However, on physical grounds we can expect that in this, essentially massless, limit, the mass scale must be set by the running cutoff $`\mathrm{\Lambda }(l_L)`$, i.e., $$\mathrm{Im}M_R\mathrm{\Lambda }_cg^{1/4}e^{\pi /g}.$$ (A21) This can be compared with the solution of the SCBA equation (A13) in the massless limit ($`\mathrm{Re}M_R=0`$), $$\mathrm{Im}M_R^{\mathrm{SCBA}}\mathrm{\Lambda }_ce^{4\pi /g}.$$ (A22) Clearly, the SCBA misses the factor of four in the exponential. As mentioned above, RG-enhanced perturbation theory formulas (A16) and (A18) do not work at low energies. Nonvanishing $`\mathrm{Im}M_R`$ can be detected using the nonperturbative method of constrained instantons. The details of the calculation will be presented elsewhere. Here we only quote the result, $$\mathrm{Im}M_R\frac{4\pi e\sqrt{6\pi }}{9g_R^2}\mathrm{ln}\left(\frac{\mathrm{\Lambda }_cg_R}{\mathrm{Re}M_R}\right)\mathrm{\Lambda }_ce^{\pi /g},$$ (A23) which is valid for $`\mathrm{Re}M_R1/L_{LL}`$, $`g_R1`$. Equations (A21) and (A23) match for $`\|\mathrm{Re}M_R\|1/L_{LL}`$. Therefore, we expect that the spectral function $`\mathrm{Im}D_L`$ exhibits a quasi-Lorentzian peak centered at $`ϵ_p=YM_{}\mathrm{\Pi }_0`$ of exponentially small width $$\mathrm{\Delta }ϵ_p=\frac{Y}{2}\mathrm{Im}M_R(0,ϵ_p)\mathrm{\Lambda }_cYg^\omega e^{\pi /g}.$$ (A24) ## B Dynamical response within the SCBA Within the SCBA the self-energy $`\mathrm{\Pi }_{\alpha \beta }(\omega ,q)`$ is diagonal, $`\mathrm{\Pi }_{\alpha \beta }(\omega ,q)=\delta _{\alpha \beta }\mathrm{\Pi }`$ and weakly $`q`$-dependent for $`qa1`$. $`\mathrm{\Pi }`$ is to be found from the self-consistency equation (38) which we reproduce here for convenience, $$\mathrm{\Pi }(\omega )=S_0+V^{\prime \prime }(0)\underset{𝐤}{}\mathrm{tr}D(𝐤,\omega ).$$ (B1) As discussed in Secs. III and IV, the SCBA applies only at relatively high frequencies. Moreover, in strong magnetic fields and near the pinning frequency $`\omega \omega _p`$, it is not even qualitatively correct. One may ask why we bother elaborating on this faulty approximation scheme here. The reason is as follows. We discovered that as far as the dynamical response is concerned, all alternative theories suggested so far are merely different forms of the SCBA. Thus, we deemed that it would be helpful to expose their interrelationship and correct a few calculational errors. Let us briefly summarize the results obtained by previous authors. Fukuyama and Lee concluded that $`\mathrm{\Delta }\omega _p\omega _p`$ but did not present the details of the calculation. Perturbative analysis of Fertig yields an exponentially small $`\mathrm{\Delta }\omega _p/\omega _p`$ (for weak disorder). Gaussian variational replica method (GVM) of Chitra et al. reduces to the SCBA at nonzero frequencies, and when analyzed further, predicts a power law dependence of $`\mathrm{\Delta }\omega _p`$ on the disorder strength. Finally, the SCBA analysis in the Appendix of Ref. leads to yet another dependence. One reason for such a disarray of conflicting results is the extreme sensitivity of the solution $`\mathrm{\Pi }(\omega )`$ of Eq. (B1) to the precise value of the dimensionless parameter $`CV^{\prime \prime }(0)/(4\pi \mu \mathrm{\Pi }_0)`$, where $`\mathrm{\Pi }_0\mathrm{\Pi }(0)`$. In principle, $`C`$ is fully determined by $`S_0`$, but $`S_0`$ is known only approximately. A strong partial cancellation of $`S_0`$ by the second term in Eq. (B1) leaves us only the order of magnitude estimate $`C1`$. One could argue, as Fukuyama and Lee did in analogous situation, that $`C=1`$ is the “best” choice because others lead to various physical absurdities. For example, if $`C>1`$, then $`\mathrm{\Pi }(\omega )`$ acquires a finite imaginary part at complex $`\omega `$, which means that the ground state is unstable; if $`C<1`$, then $`\mathrm{\Pi }(\omega )`$ has an imaginary part only above some threshold frequency $`\omega _{\mathrm{th}}>0`$, i.e., the phonon spectrum has a gap, which also appears to be unphysical. The truth is, of course, that the SCBA is simply unable to correctly describe the properties of the $`\omega 0`$ tail; therefore, it cannot be relied upon for selecting a “good” value of $`C`$. More sensible approach is to investigate a certain range of $`C`$ around unity, hoping that certain quantities depend on $`C`$ only weakly. Alas, the linewidth $`\mathrm{\Delta }\omega _p`$ is not one of such quantities. We found that the results of Chitra et al. are recovered for the “Fukuyama-Lee choice,” $`C=1`$. (Therefore, these two groups of authors should have obtained the same results). A slight reduction of $`C`$ from unity is sufficient to cross over to very different predictions of Fertig. Let us now demonstrate this in more detail. Using the definition of $`C`$ and Eq. (B1), we obtain $$\mathrm{\Pi }(\omega )\mathrm{\Pi }_0=2\mu C\mathrm{\Pi }_0\underset{0}{\overset{\mathrm{}}{}}𝑑kk\left[D_T(k,0)D_T(k,\omega )\right].$$ (B2) Let us now describe the solution $`\mathrm{\Pi }(\omega )`$ of this equation for different $`C`$ in the weak pinning regime $`R_ca`$, which corresponds to the inequality $`\alpha \sqrt{\mathrm{\Pi }_0\mu }/Y=\mu /\lambda 1`$. 1. $`0<1C\alpha `$. In this case $$\mathrm{\Pi }(\omega )C\mathrm{\Pi }_0+i\mathrm{\Pi }_0\sqrt{(\omega /\stackrel{~}{\omega }_p)^2(1C)^2},$$ (B3) where $`\stackrel{~}{\omega }_p=\omega _p/(\pi \alpha )^{1/2}`$ and $`\omega _p=\mathrm{\Pi }_0/\rho \omega _c`$. As one can see, the threshold frequency is $`\omega _{\mathrm{th}}=(1C)\stackrel{~}{\omega }_p`$. For $`C=1`$ Eq. (B3) coincides with that of Ref. . Strictly speaking, Eq. (B3) describes the solution of Eq. (B2) only for $`\omega \omega _p`$. Near the resonance, $`\omega \omega _p`$, it is off by a numerical factor of the order of unity. In principle, this factor can also be calculated. We limit ourselves to the order of magnitude estimate $`\mathrm{Im}\mathrm{\Pi }(\omega _p)\alpha ^{1/2}\mathrm{\Pi }_0`$, which leads to $`\mathrm{\Delta }\omega _p/\omega _p\alpha ^{1/2}=(\mu /\lambda )^{1/2}`$. Although this result is of the same form as Eq. (76), with $`s=1/2`$, we believe that $`s=3/2`$ is correct, see Secs. IV and V. 2. $`1C\alpha `$. Equation (B3) still applies as long as $`\omega _{p0}\omega `$ is positive and not too small. At larger frequencies, $`\mathrm{\Pi }(\omega )`$ receives additional contribution from the pole $`k_{}`$ of $`D_T`$ in the complex plane of $`k`$, which moves close to the real axis. The resonance linewidth $`\mathrm{\Delta }\omega _p`$ is exponentially small, $$\mathrm{\Delta }\omega _p=\frac{\mathrm{Im}\mathrm{\Pi }(\omega _p)}{\rho \omega _c}\omega _pe^{F/4\alpha ^2},$$ (B4) where $`F\mathrm{\Pi }(\omega _{\mathrm{th}})/(C\mathrm{\Pi }_0)1`$. The threshold frequency $`\omega _{\mathrm{th}}`$ is only slightly below $`\omega _p`$, by an amount of the order of $`\mathrm{\Delta }\omega _p`$. To be exact $`\omega _{\mathrm{th}}`$ can be found from the condition that the solution of Eq. (B2) is also the solution of the “derivative” of this equation with respect to $`\mathrm{\Pi }`$, $$1=2\mu C\mathrm{\Pi }_0\underset{0}{\overset{\mathrm{}}{}}\frac{dkk[ϵϵ_c+(\mathrm{\Pi }+Yk)^2]}{[(\mathrm{\Pi }+\mu k^2)(\mathrm{\Pi }+Yk)ϵϵ_c]^2}.$$ (B5) At frequencies somewhat higher than $`\omega _p`$, $`\mathrm{Im}\mathrm{\Pi }(\omega )`$ is dominated by the aforementioned pole at $`k_{}=(ϵϵ_c\mathrm{\Pi }^2)/\mathrm{\Pi }Y`$ near the real axis, $`\|\mathrm{Re}k_{}\|\|\mathrm{Im}k_{}\|`$. $`\mathrm{Im}\mathrm{\Pi }`$ is given simply by the residue of this pole, $$\mathrm{Im}\mathrm{\Pi }(\omega )2\pi \alpha ^2C\mathrm{\Pi }_0\left(\frac{\omega ^2}{\omega _p^2}1\right).$$ (B6) In conclusion, we would like to reiterate that neither of Eqs. (B3), (B4), or (B6) describes the dynamical response of the system correctly. We derived them here just to facilitate the comparison with the previous work on the subject. The correct expressions for $`\mathrm{Im}\mathrm{\Pi }`$ and $`\sigma _{xx}`$ are given in Sec. IV.
warning/0004/quant-ph0004053.html	ar5iv	text	# Quantum computing with trapped ions, atoms and light ## 1 Introduction This paper will discuss various issues in the physics of ion traps and quantum information. Our purpose is to address some quite general questions about quantum information physics and ion (or atom) traps, with the aim of identifying useful directions for theoretical and experimental research in the near future and the longer term. We begin in section 2 by considering the major requirements of a quantum computer without assuming any particular technology from the outset. Rather, we try to identify physical phenomena which appear to be intrinsically well-suited to quantum computing, using arguments based as much as possible on fundamental physical principles. We find that there is a natural link with methods in quantum optics, such as ion and atom trapping and cooling. In section 3 we focus our attention on currently achieved experimental results, considering the particular strengths of ion trap methods. Section 4 briefly surveys the main experiments in quantum information physics which are feasible in the near future using ion traps. These include several fundamental quantum information effects which have not yet been observed in any area of physics. In section 5 we examine how far ion trap methods can go: we estimate the major limitations to the gate rate, at given precision, for two different methods to implement the gates. These are the coupling via the vibrational degree of freedom which has been used up till now, and coupling via cavity quantum electrodynamics (CQED) methods. In section 6 we sketch a design for a moderately large quantum computer: one which could perform $`10^6`$ Toffoli gates on 100 logical qubits. This computer, based on atomic physics and CQED concepts, is conceivable using current technology, and the quantum optics methods it is based on are probably necessary in any case for quantum communication links. It illustrates the power of these methods, and underlines the interest of further experiments, and theoretical studies, in this area. ## 2 An ideal physical system for quantum computing We would like to consider the question, what might be the ideal system for a future quantum computer? Here, we do not intend to restrict attention to any one branch of physics (or other science). Rather, we want to know what system we might choose, if we are guided only by basic physical principles and the properties of systems which, in some useful sense, Nature provides. We would like our ideal quantum computer to have the highest quality memory, logic gates, and read-out. We note that the read-out gives automatically the ability to prepare simple initial states such as product states. By “highest quality” we mean primarily reliable operation, but if the system or the gate can also be small or fast, then so much the better. For a quantum memory, we want a system of qubits which does not interact with anything else, while for logic gates we want a coupling between qubits which is fast, and not coupled to anything other than the bits. These two demands are almost, but not quite, contradictory. They imply that a quantum computer should be composed of entities $`q`$ which are, in their passive state when no gates are switched on, almost completely isolated, and yet which can be coupled rapidly. This means they must have a strong coupling to something, $`𝒢`$. The conflicting demands are met if $`𝒢`$ has the property that it can be made to be wholly absent when it is not wanted, and introduced rapidly when it is wanted. Furthermore we would like $`𝒢`$ to interact only with the entities $`q`$ and with nothing else. It turns out that Nature does provide a physical entity $`q`$ which meets the contradictory demands of memory and logic gates better than we might have imagined possible. This entity is the nuclear spin. The advantages of nuclear spins for quantum computing are already well recognised. A spin has a smaller coupling to its environment than any degree of freedom based on charge or the motion of particles; a nuclear spin has a particularly small magnetic moment, and this tiny magnet comes ready-packaged in an electron cloud with highly useful properties for logic gates. The atomic electrons provide the means to take hold of the atom and place it where we wish, and they also provide a ready-made very strong and very stable magnetic field on the nucleus. This results in the hyperfine splitting. The stability of this splitting for isolated atoms or ions is well documented, and is in fact used to provide our standards of time and frequency. The existence of hyperfine structure makes it possible to couple to the nuclear spin via the electronic state. This provides the handle whereby logic gates can be achieved. The next question is, what is the best way to grasp this handle? The existing proposals which are based on the nuclear spin and/or hyperfine splitting differ in the way the coupling $`𝒢`$ is brought about. These proposals include bulk nuclear magnetic resonance (NMR) , ion trap and other atomic-physics-based methods , and a proposal for future solid state devices based on nuclear spins of dopant atoms implanted in a semiconductor . In bulk liquid-state NMR, the electronic ‘handle’ is almost entirely ignored, and the method relies instead on the tiny direct spin-spin interaction between neighbouring nuclei in a molecule. This permits logic gates but not a direct measurement of the spin state. The fact that all neighbouring qubits are permanently coupled in such methods has both advantages and disadvantages. The other methods all use the electronic ‘handle’. Ion trap experiments up till now have used a light-induced coupling between the electronic state and the vibrational motion (phonon) of relatively heavy charged particles (ions) in a trap in vacuum . There are proposals in which light alone is used to couple the electronic state of one atom and another , and a realisation of this concept (in an experiment not based on nuclear spin or hyperfine interaction) using a beam of neutral atoms . The solid state proposal is to use low-mass charged particles (electrons) moving in a solid to provide the coupling . Elaboration on the above summary would enable us to see various strengths and weaknesses of current proposals. However, our purpose here is to seek a method which appears to be in some sense natural, that is, which makes use of physical principles which are naturally suited to the task. We suggest that the natural, and possibly in the long term the best, choice for the entity $`𝒢`$ which provides rapid controllable coupling, is light. Light is in any case the fundamental coupling between charges. It will travel at the fastest possible speed between qubits, it can be made to appear and disappear at will and, perhaps most importantly, it does not couple to extraneous electromagnetic fields in the computer, which greatly reduces possible decoherence mechanisms. Furthermore, photons provide a natural way to couple quantum information out of the computer and down a quantum communication link. To ensure there is no light when it is not wanted, we should use frequencies well above the thermal spectrum at the temperature of the computer, but otherwise the main consideration is ease and precision of manipulation of the light (including the ability to select individual qubits). This suggests near-infra-red or optical frequencies. When we wish to couple resonantly to the hyperfine splitting, which is in the microwave domain, we use Raman transitions. Note that the typical frequency scale for hyperfine splitting, i.e. GHz, is attractive because electronic techology allows the most precise control in this frequency regime. This is likely to remain true in the future, owing to basic properties of electromagnetism and conduction in metals. At this point in the argument, we may envisage an ideal quantum computing system as based on an array of nuclear spins inside atoms, the spins coupling to the electrons of their atoms, and the electrons coupling to photons which ferry information around the computer, appearing and disappearing at the behest of the controlling machinery. The only remaining question is, how is the electron-photon coupling to be both strong enough, and under sufficient control? To achieve a strong enough coupling, the light must be confined in a small volume, and to permit coherent coupling we require a long photon storage time in the confining cavity, as well as accurate positioning of the atoms. These considerations will be addressed in section 5. It is possible to imagine that the atoms might be held in place by any one of a number of methods, including attaching them to long chain molecules or fabricating nanoscale structures to hold them. However, the additional atoms and electrons which form the body of any such structures will introduce new degrees of freedom which may cause decoherence, or weaken the light-atom coupling. One possibility is to situate the atoms on the surface, or perhaps under the surface, of a highly transparent solid (see section 6). In this paper we will concentrate on the case that the atoms are held in an r.f. Paul trap (ion trap) or an optical dipole force trap, and consider the possibility of placing the atoms on a surface in the final section. We note that whereas we have advocated using the nuclear spin alone as quantum memory, current experiments designed to achieve quantum information processing in ion traps are not operating in this regime. In the work of Wineland et al. with the beryllium ion the qubit involves both nuclear and electron spin, since its energy level separation is a sum of hyperfine and first-order Zeeman effects. A pair of electronic states (Coulomb interaction with the nuclear charge) with first order Zeeman effect is adopted in and the electron spin alone in . We envisage that all these experiments will contribute to the overall development of the field, and it will be a relatively small step to adapt them to hyperfine transitions with no electron spin component. ## 3 Strengths of ion trap technology Before discussing future prospects, we will highlight in this section the strengths of current ion trap technology. We consider the three requirements for quantum information processing, which are quantum memory, quantum logic gates, and measurement of quantum states. The primary consideration for all of these is precision and reliability, simply because we need the computer to work; we are willing to sacrifice both speed of operation, and ease of construction, if it makes the difference between a computer which works and one which does not<sup>1</sup><sup>1</sup>1For a large computer this will, of course, only be possible if the speed does not fall, nor the system size increase, exponentially with the number of qubits.. Note that all the three requirements are equally significant. In particular, measurement is at the heart of error correction protocols , therefore it is a central consideration during the whole operation of the computer, not merely at the final step where the computation result is measured. For some purposes, it is sufficient that a dissipation process can be applied to chosen quantum bits at chosen times . The dissipation process forces the qubit to a known final state, no matter what its initial state. This can be easier to implement, so will be considered also. ### 3.1 Quantum memory and single-qubit gates We discussed the attractive features of nuclear spins for quantum computing in section 2. Experiments with trapped atoms and ions offer the most precise methods known for manipulation of the nuclear spin, via the hyperfine interaction. Indeed, time and frequency standards throughout the world are based on optical manipulation of atoms trapped in high vacuum, and ion trap frequency standards now rival those based on neutral atoms. This is the first advantage of ion trap methods. From a practical point of view it means that the quantum memory and single-qubit gates are, broadly speaking, solved problems, in that we can envisage trapped ions whose nuclear spin state is as accurately preserved and manipulated as anything which current technology allows. ### 3.2 Read-out The read-out is also, broadly speaking, a solved problem for experiments with trapped atoms or ions. The measurement of the hyperfine state can be carried out rapidly by the electron shelving (or ‘quantum jump’) method, which offers close to 100% reliability . The timescale for such a measurement is set by the need to scatter a few thousand photons on an allowed atomic transition, requiring of order a few hundred $`\mu `$s. For dissipation, we can use the method of optical pumping. Here, only a few photons need to be scattered before we can be confident the system has relaxed, so the time scale for controlled dissipation is of order $`0.1\mu `$s. A further point about measurement is significant: as long as separate atoms or ions can be resolved by an optical imaging system (implying a separation of at least a few wavelengths) then they can be measured simultaneously. The central problem for ion trap quantum computing is, then, the question of implementing the 2- or more-bit logic gates. This will be discussed in section 5. ### 3.3 Quantifying qubits, gates and entanglement The standard way to quantify the complexity of an algorithm on any computer, whether quantum or classical, is to count the number of bits in the memory and the number of 2-bit or 3-bit logic gates used. In the case of quantum computing, it makes sense also to have a measure of the degree to which an algorithm involves highly non-classical effects. A useful measure is to ask whether $`n`$-particle entangled states, such as the “Schrödinger cat” state $`\|000\mathrm{}0+\|111\mathrm{}1`$, can be produced. So far no ion trap experiment has combined all the necessary features to allow general processing on more than one ion. However, all the ingredients of general processing have been demonstrated in separate experiments, and highly entangled states have been produced. The definition of entanglement requires some comment. If two or more separate spin-half particles are in a joint state $`\|\varphi `$, then the existence of a non-zero overlap with an entangled state does not necessarily imply the presence of entanglement. For example, the overlap between the 4-qubit separable state $`\|++++`$ (where $`\|+=(\|0+\|1)/\sqrt{2}`$) and the cat state $`(\|0000+\|1111)/\sqrt{2}`$ is $`1/8`$. Also, a superposition $`\|M=+3/2+\|M=3/2`$ of the two stretched states of a spin-$`3/2`$ particle is not entangled, even though it could be written $`\|0\|0+\|1\|1`$ by a suitable choice of state labels. The latter point is important: the mere fact that a state can be written $`\|0\|0+\|1\|1`$ in some basis is not enough to mean that it is entangled in any sense which is significant to quantum information physics. A strict definition of the term “entanglement” would restrict its use to refer only to degrees of freedom which could in principle be located in separate spatial locations (so that entanglement-enhanced communication could be realised), or which could be used to gain the reduction in computational complexity offered by quantum computation (compared to classical computation) for certain algorithms. In that case the state of spin and motion of an electron emerging from a Stern-Gerlach apparatus would not be regarded as entangled. However, it has become quite common to broaden this strict definition slightly, so as to include the case of “entanglement” between the internal and motional degree of freedom of a single particle. Such entanglement has been achieved between the internal state of a single trapped ion and its motional state, with a high degree of precision and control, in at least two laboratories . A measure of the degree of entanglement is the size of the Hilbert space in which coherent evolution is demonstrated in the experiment. For example, the “Schrödinger cat” states realised in involve a superposition of coherent states. Each coherent state has a Poissonian distribution over vibrational levels, characterised by a parameter $`\alpha `$ which had the value $`\alpha =2.97\pm 0.06`$ in the experiments. The mean vibrational quantum number $`n=\|\alpha \|^29`$, and standard deviation $`\sigma _n\|\alpha \|3`$. The size of the motional Hilbert space in which coherent evolution must take place in order to observe the interference is of order $`\mathrm{log}_2(n+1+\sigma _n)3.7`$ qubits. Adding the internal degree of freedom, this is a “cat state” of $`4.7`$ qubits. A true multi-particle entanglement is very rare in physics, and indeed for more than three particles it has only been achieved, to our knowledge, in a single experiment. This is the 4-particle entanglement recently demonstrated in an ion trap experiment by Sackett et al. . The strength of ion trap experiments which is underlined by this achievement is that the generation of entanglement is under complete experimental control: it is deterministic, rather than being the result of a process which relies on an essentially random event (e.g. spontaneous parametric down-conversion, or velocity selection of atoms from a thermal source). This is a significant distinction because the amount of $`n`$-qubit entanglement produced by a random process falls off exponentially with the number $`n`$ of qubits, and therefore results in a system which cannot exhibit some of the essential defining features of quantum computation, such as the breaking of the classical hierarchy of complexity classes. It is noteworthy that in all experimental “realisations” of quantum algorithms so far reported, the size of the apparatus, or the duration of the experiment, has scaled with the number of qubits required to define the problem at the same rate or worse than a classical computer or information channel would scale with the number of classical bits. This does not mean that randomly produced entanglement is uninteresting, since it can be used to demonstrate some of the basic principles of quantum mechanics and quantum information. However, one might draw an analogy with the properties of light sources: a thermal source, with a sufficiently narrow filter in front of it, can produce radiation with just as narrow a bandwidth as is available from a laser, but there remains a qualitative, and practically significant, difference between thermal radiation and laser radiation. For light sources, a useful parameter which emphasizes that bandwidth is not the only consideration is the number of photons per mode. It would be useful to have a comparable measure for entanglement, such as “the number of singlets per 2-qubit Hilbert space” (a singlet being the 2-qubit entangled state $`(\|01\|10)/\sqrt{2}`$). The difficulty in forming such a measure is that the Hilbert space size, unlike the modes of a radiation cavity, depends on which parts of the system we choose to focus our attention on. For example, we may consider all the atoms in a thermal beam, or just those selected by a velocity selector. The least ambiguous measure is arguably that implicit in : we define the entangling efficiency $`ϵ`$ to be > the probability that, starting from initial conditions of no entanglement, a singlet can be caused to be present in a predetermined Hilbert space at a predetermined time. The predetermined Hilbert space means we indicate which systems (eg atoms, spins, photons) will contain the singlet, without the need to check by measuring them, and the predetermined time means we decide beforehand at which moment we want the singlet, without reference to the details of the experimental apparatus (thus ruling out statements such as “1 ms after detector D clicks”, if we can’t predict when detector D will click). The purpose of the quantity defined is to enable us to assess rapidly the slow-down to be expected when the same apparatus is used to form 3-particle, 4-particle and higher forms of entanglement. In parametric down-conversion experiments reported to date , $`ϵ`$ was of order $`10^4`$, and in cavity QED experiments using a thermal atomic beam, it was $`ϵ3\times 10^3`$. For thermal ensembles such as those in current liquid state NMR experiments it is zero . A related quantity, the entangling rate (number of successful singlet-generating runs per unit time) was approximately 8000 s<sup>-1</sup> and 2 s<sup>-1</sup> respectively for the down-conversion and CQED experiments. The first observation of an entangling efficiency of order 1 was in the experiment of Turchette et. al. . The internal state of two trapped beryllium ions was driven to a singlet state with reliability approximately 70% (with entangling rate approximately 30000 s<sup>-1</sup>). The recent report of 4-qubit entanglement arose from further work in the same laboratory . These remain the only demonstrations of a high entangling efficiency in any area of physics. ## 4 Experiments feasible in the short term We have noted that the main strengths of current ion trap experiments, compared with other quantum information experiments, are that they allow rapid and reliable measurement, and deterministic entanglement. The following list concentrates on experiments which exploit these strengths, identifying goals which are either not realisable at all in other systems, or for which the ion trap may be the system of choice. A single trapped ion allows the experimental exploration of two important avenues: the vibrational degrees of freedom, and cavity QED . The interest of the vibrational degrees of freedom is illustrated by the Schrödinger cat and environment engineering experiments which further our understanding and control of decoherence. With 2 ions in the same trap, some standard quantum information ideas can be demonstrated, such as the EPR experiment , “dense coding” and a simple “algorithm” such as Grover’s search algorithm . Of these, the EPR experiment is the most significant, since the detector efficiency problem can be avoided ; however the close spacing of the ions makes impractical a test involving space-like separated measurement processes. A demonstration of dense coding would be the first time this idea had been implemented without needing post-selection, and hence allowing an unambiguous increase in the capacity of the quantum channel to transmit classical information. However since the ‘channel’ involved only covers a distance of some tens of microns in vaccum, it is of no practical use. With only 2 qubits it is debatable whether the most significant features of Grover’s algorithm can be demonstrated, but the algorithm would provide a useful way of showing that quite general manipulations of a two-ion system had been achieved. With 3 ions three highly significant experiments could be done. These are “teleportation” , entanglement-enhanced communication , and quantum error correction . In addition, a thorough (though still very simple) demonstration of Grover’s algorithm would be possible. Quantum teleportation is significiant not only in the context of quantum communication, but also as an essential ingredient of fault-tolerant quantum processing. A reliable teleportation experiment within a small quantum processor would therefore be a significant development. The most accessible example of entanglement-enhanced communication is the “Guess my number” protocol , in which three parties use shared entanglement and classical communication to learn the answer to a simple mathematical problem. In order to obtain a result which breaks the classical limits on communication, an experiment of overall reliability above 50% is needed. The simplest example of quantum error correction requires 3 qubits, which are used to protect a single logical qubit against a restricted class of errors . The set of correctable errors could be, for example, phase errors on single bits. The most striking result is obtained, however, if the errors are not merely unitary precession of the qubits themselves, but non-unitary relaxation processes where information leaks away into the environment. For example, optical pumping could be used to cause a relaxation of one qubit, where, after tracing over the environmental degrees of freedom, the qubit has ‘collapsed’ into a mixed state, with density matrix $`P_0\|00\left\|+P_1\right\|11\|`$. After this, the correction network is applied, and it would still recover the exact encoded state in the three qubits. Furthermore, the process of random error followed by correction, could be repeated many times on the same encoded state. This process is remarkable from several points of view. After it is repeated a few times, the environment would have had a chance to “measure” all the qubits in the processor, thus causing, one would think, a large perturbation to the state, and yet the qubit of information is perfectly preserved. Alternatively, if we drive optical pumping continuously but weakly on all the qubits, then the loss of fidelity of the qubits is linear with time, for small times, while after correction it becomes quadratic with time, therefore allowing the Zeno effect to be implemented in the case of a relaxation process . The 3-qubit quantum error correction has been investigated in an NMR experiment . Some aspects of the expected behaviour were demonstrated, but owing to the limitations of the pseudo-pure state method, a genuine error correction was not available, since the entropy could not be extracted from the system. This is seen most clearly in two aspects of the experiment. First, in encoding from one qubit into three, the signal size fell by a factor 8, and no subsequent error correction can make up for the increased sensitivity to errors in this situation. Secondly, it was not possible to apply correction repetitively. Note that all the experiments we have listed for three ions rely on the ability to perform not just unitary processing operations, but also strong measurements of one or more chosen qubits. Successful realisation of the “Guess my number” or the repeatable quantum error correction protocols would be landmarks in quantum information science. Obviously, there are more and more experiments which are possible as the number of qubits increases, even before useful quantum computation becomes possible. ## 5 Logic gate methods in ion traps The first general method proposed to implement quantum logic gates between trapped ions was that discovered by Cirac and Zoller . This is based on using the motional degree of freedom to ferry quantum information from one ion to another. All ion trap quantum processing experiments so far have been based on this idea. A significant further insight was provided by Mølmer and Sørensen who showed how to make better use of the motional degree of freedom, and the recent experiments of are based on these further insights. Another method to couple separate atoms or ions coherently is to use light to ferry the information around, using the proposal of based on concepts in cavity quantum electrodynamics (CQED). In this section we will compare the two methods—motional and photon-based gates. We assume some familiarity with both methods on the part of the reader. At present the motional methods are easier to achieve experimentally, but the CQED methods allow, in principle, higher gate rates, and also quantum communication between separate ion- or atom-traps. There is a subtlety regarding the hyperfine interaction and the optical transitions involved in these methods. The electric dipole optical transitions which we will use do not couple directly to the nuclear spin. In the motional coupling, this implies that the change of internal state of the ion must involve the electronic wavefuntion, so it is not purely a nuclear spin rotation. As a result, the relevant hyperfine levels will typically have a first-order Zeeman effect. In the CQED method, 4 states in the ground hyperfine manifold of a single ion are used. In either case, during the action of the gate, the quantum information is stored in electronic not nuclear degrees of freedom. However, the quantum memory can remain a wholly nuclear spin system: we swap quantum information between Zeeman levels $`\|M_F`$ just before and just after each gate, by driving a Raman transition in the internal state of the ions involved in the gate. We assume this transition can be fast compared to the gate operations to be discussed. ### 5.1 Motional coupling The Cirac-Zoller method to implement 2-qubit gates such as ‘controlled not’ between separate ions is to couple the internal state of chosen ions to the vibrational degree of freedom, by driving Rabi flopping on a vibrational sideband of the atomic transition. The phenomena which cause the main limitations of this method are relaxation and/or heating, and off-resonant driving of unwanted transitions. #### 5.1.1 Relaxation There are two main sources of relaxation in the ion trap. These are the spontaneous decay of excited states of the ion, and the heating or relaxation of the vibrational degree of freedom. To minimise the effects of these, a compromise between fast and slow operation of the processor is needed. The quantum gates between ions involve the excitation of the motional degree of freedom, so we consider driving the first red vibrational sideband of a resonant Raman transition between hyperfine levels in a trapped ion , see figure 1. The Raman transition is driven by a pair of lasers detuned by $`\mathrm{\Delta }\mathrm{\Gamma }`$ from an allowed single-photon transition whose natural width is $`\mathrm{\Gamma }`$ (full width half maximum in angular frequency units). The Rabi frequencies of the relevant single-photon transitions are $`\mathrm{\Omega }`$ and $`g=\eta \mathrm{\Omega }`$, where $`\eta `$ is the Lamb-Dicke parameter. In this situation the pair of hyperfine levels connected by the Raman transition form an effective two-level system; the two-level transition has effective Rabi frequency $`\mathrm{\Omega }_{\mathrm{eff}}=\mathrm{\Omega }g/2\mathrm{\Delta }`$. A two-ion gate such as a state-swapping operation requires two $`\pi `$ pulses on a vibrational sideband, so the total time is $`T=2\pi /\mathrm{\Omega }_{\mathrm{eff}}`$. During each pulse the mean population of the unstable excited state of the ion is $`\mathrm{\Omega }^2/4\mathrm{\Delta }^2`$ (assuming $`g\mathrm{\Omega }`$), therefore the mean number of photons scattered is $$p_1=\frac{\mathrm{\Omega }^2}{4\mathrm{\Delta }^2}\mathrm{\Gamma }T=\frac{\pi \mathrm{\Gamma }\mathrm{\Omega }}{\mathrm{\Delta }g}.$$ (1) We assume that at all times the vibrational state of the ion suffers a non-unitary heating process, characterised by a rate $`\kappa `$ which is the rate of heating (or relaxation) from one vibrational state to an orthogonal one. Therefore the probability of relaxation by this process, during the two pulses, is $$p_2=\kappa T=\frac{4\pi \kappa \mathrm{\Delta }}{\mathrm{\Omega }g}.$$ (2) The total probability of failure is $$p=p_1+p_2=\frac{\pi \mathrm{\Gamma }}{g}\frac{\mathrm{\Omega }}{\mathrm{\Delta }}+\frac{4\pi \kappa }{g}\frac{\mathrm{\Delta }}{\mathrm{\Omega }}.$$ (3) In this equation $`\mathrm{\Gamma }`$ is constant for a given atom, and $`\kappa `$ is characteristic of a given experimental apparatus, while $`\mathrm{\Omega }/\mathrm{\Delta }`$ can be adjusted to minimise $`p`$. This minimisation gives $`\mathrm{\Omega }/\mathrm{\Delta }=2(\kappa /\mathrm{\Gamma })^{1/2}`$, and $$p_{\mathrm{min}}=4\pi \sqrt{\frac{\kappa \mathrm{\Gamma }}{g^2}},\frac{1}{T}=\frac{g^2}{\mathrm{\Gamma }}\frac{p_{\mathrm{min}}}{8\pi ^2}.$$ (4) #### 5.1.2 Off-resonant coupling The vibrational levels in an ion trap are typically closely spaced compared to all other energy level separations in the system, so the transitions driven off-resonantly are primarily the carrier transitions (those which don’t change the vibrational state), which are off-resonant by the vibrational frequency $`\omega _z`$. This problem is studied in detail in . The conclusion is that after two $`\pi `$ pulses, the amount of population which leaks into unwanted states due to off-resonant coupling is $$p_3=\frac{\mathrm{\Omega }_{\mathrm{eff},0}^2}{\omega _z^2}=\left(\frac{\mathrm{\Omega }^2}{2\mathrm{\Delta }}\right)^2\frac{1}{\omega _z^2}.$$ (5) where $`\mathrm{\Omega }_{\mathrm{eff},0}`$ is the Rabi frequency for carrier transitions. #### 5.1.3 Discussion Of the processes we have considered which limit the motional coupling, the atomic relaxation $`\mathrm{\Gamma }`$ and the off-resonant excitation are intrinsic to the physics of the system, while the motional heating $`\kappa `$ could in principle be made arbitrarily small (in recent experiments values of $`\kappa /\omega _z`$ as low as $`2\times 10^7`$ for a single ion and $`2\times 10^6`$ for the stretch mode of two ions were reported.) Therefore equations (1) and (5) give the main limitations. The gate rate is limited by (5) since for a given value of $`\mathrm{\Omega }_{\mathrm{eff},0}/\omega _z`$ we can make $`p_1p_3`$ by increasing the laser intensity and its detuning $`\mathrm{\Delta }`$ (until the detuning becomes comparable to further energy-level separations in the ion, such as the fine structure). In the case that motional heating is significant, the Mølmer-Sørensen approach may be advantageous, since it permits gates of high fidelity in the presence of motional heating (at the expense of reduced gate rate). However, in the limit of small $`\kappa `$, the gate rate produced by this method in its standard form is limited by off-resonant excitation and is the same as that given in equation (5). In conclusion, the gate rate at given failure probability $`p`$ for a 2-qubit swap gate via the motional state is (from (5)) $$\frac{1}{T}=p^{1/2}\eta \frac{\omega _z}{2\pi },$$ (6) assuming $`\kappa p/T`$ and $`\mathrm{\Delta }\pi \mathrm{\Gamma }/\eta p`$, and using either the Cirac-Zoller or the Mølmer-Sørensen methods. It is notable that this limit, imposed by off-resonant carrier transtions, could be exceeded by exciting the ion in the node of a laser standing wave . That is technically very difficult, but it shows that the physics of the system can allow a faster switching rate. Recently, less demanding methods to gain such a faster rate have been proposed. The Mølmer-Sørensen approach can probably be made to produce faster gates than (6) by a careful choice of parameters , and recently a new approach has been put forward in which such a speed up is thoroughly analysed . The latter uses a light-shift-induced resonance, yielding a gate rate $`\eta \omega _z/2\pi `$ and gates of fidelity approximately $`1\eta ^2/2`$. Therefore the speed increase compared to eq. (6) is significant, of order $`1/\eta `$. ### 5.2 Photon-based coupling We will now consider coupling qubits via the excitation of a mode of a high-finesse optical cavity. We will assume an allowed electric dipole transition, so the Rabi frequency describing the coupling is $`g=Ed/\mathrm{}`$ where $`E`$ is the electric field of the light, and $`d`$ is the electric dipole matrix element. We will use this coupling to exchange quantum information between an atom and a light field by absorption or emission of single photons, therefore we are interested in the value of $`g`$ when the electric field is that of a single photon. If the photon has angular frequency $`\omega `$ and occupies a mode of volume $`V`$, then its energy is $`\mathrm{}\omega =ϵ_0E^2V/2`$, hence $$g=d\sqrt{\frac{2\omega }{ϵ_0\mathrm{}V}}$$ (7) The strongest electric dipole matrix elements in atoms are all of order $`ea_0`$ where $`e`$ is the charge on the electron and $`a_0`$ is the Bohr radius. The spontaneous decay rate $`\mathrm{\Gamma }`$ of an atom on a strong transition varies as $`1/\lambda ^3`$: $$\mathrm{\Gamma }=\frac{\omega ^3d^2}{3\pi ϵ_0\mathrm{}c^3}$$ (8) so for $`g\mathrm{\Gamma }`$ the long wavelength region is best. However, we will need the logic gates to be fast, setting a premium on large $`g`$, hence small wavelengths. The other major source of decoherence is decay of the photon mode owing to the finite finesse of the cavity which contains it. This decay rate is $$\kappa =\frac{c\pi }{L}$$ (9) for a cavity of length $`L`$ and finesse $``$. #### 5.2.1 Dark state and adiabatic passage Since we need very precise gates, there is interest in any method to implement them which has reduced sensitivity to the relaxation of the atom and the cavity photon. Such a method is the adiabatic passage, as described in . Briefly, a state-swapping operation is carried out between any two atoms in the cavity by shining laser pulses on the two atoms. The pulses are not exactly simultaneous, but overlap in time, coupling ground states $`\|a_i`$ to each atom’s excited state with Rabi frequencies $`\mathrm{\Omega }_i`$, for atoms $`i=1,2`$ respectively. The cavity mode produces strong coupling between the excited state and a metastable level $`\|b_i`$ simultaneously for both atoms. In our case, $`\|b_i=\|F^{},M_F^{}_i`$ is in the ground state manifold, separated from $`\|a_i=\|F,M_F_i`$ by the hyperfine interaction. The system of two atoms plus cavity photon exhibits the phenomenon of dark states, i.e. superpositions of states which by quantum interference are decoupled from the excited states. For up to one photon in the cavity, there are two dark states , $`\|D_0`$ $`=`$ $`\|b,b,0\|b_1\|b_2\|0_c,`$ (10) $`\|D_1`$ $``$ $`\mathrm{\Omega }_1g\|b,a,0+\mathrm{\Omega }_2g\|a,b,0\mathrm{\Omega }_1\mathrm{\Omega }_2\|b,b,1.`$ The swap operation carries one qubit of information from one atom 1 into atom 2. A general operation such as controlled-not is then carried out within atom 2 using four of its states (2 Zeeman components of each of 2 hyperfine levels), then the information is swapped back. Any other atoms in the cavity do not participate because they are not illuminated by the laser pulses, and they are in the states $`\|a`$ which are not coupled to the cavity photon. To ensure the off-resonant coupling is sufficiently small, we will require the hyperfine splitting to be much larger than $`\mathrm{\Omega },g`$. The method of adiabatic passage is limited by two considerations. First we need to preserve adiabaticity, and second we need to avoid populating states which suffer non-unitary relaxation. Numerical solution of the master equation for the complete system is discussed in . Here we will make rough estimates in order to identify the best operating regime, for given parameters $`g`$, $`\mathrm{\Gamma }`$, $`\kappa `$. To preserve adiabaticity the rate of change of the conditions must be slow compared to the frequency separation between the state we wish the system to remain in (here, the dark state) and any other state (here, the nearest bright state). For example, if the frequency separation $`\omega _{ij}`$ and the rate of change of the state are constant in time, then the probability to make an (unwanted) transition from the desired state $`\|i`$ to some other state $`\|j`$ after time $`t`$ is $$P_{ij}\left\|\frac{dH_{ji}}{dt}\frac{1}{\mathrm{}\omega _{ji}^2}\right\|^22\left(1\mathrm{cos}\omega _{ji}t\right).$$ (11) In our case, we will assume the laser pulses have the form $`\mathrm{\Omega }_2(t)=\{0,\mathrm{\Omega }t/T,\mathrm{\Omega }\}`$ for $`\{t0,0<t<T,tT\}`$ respectively, and $`\mathrm{\Omega }_1=\mathrm{\Omega }\mathrm{\Omega }_2`$, therefore $`dH_{ji}/dt=\mathrm{}\mathrm{\Omega }/T`$. The frequency separation between the dark state and the nearest bright state is a complicated function of the Rabi frequencies, which we simplify to $`\omega _{ji}(t)(\mathrm{\Omega }_1^2+\mathrm{\Omega }_2^2)^{1/2}`$ (it will emerge that for the cavities we will consider, we will need $`\mathrm{\Omega }g`$ in order to minimise the effects of relaxation of the cavity mode). For $`P_{ij}1`$ the oscillating term in (11) averages to zero during any interval of time small compared to $`T`$. We find the probability to make a transition out of the dark state by using the derivative of (11) with respect to time, and then integrating from $`t=0`$ to $`T`$, to obtain $$p_1\frac{4}{T^2\mathrm{\Omega }^2}.$$ (12) During the adiabatic passage, there is population in the state $`\|b_1\|b_2\|1`$ which decays at a rate $`\kappa `$. We model this by assuming that at any time $`t`$ the system is in a mixture containing dark state population $`1P_{bb1}\kappa 𝑑t`$, and non-dark state population $`P_{bb1}\kappa 𝑑t`$, where $`P_{bb1}`$ is the population of the $`\|b_1\|b_2\|1`$ component in $`\|D_1`$, which, from (10), is approximately $`\mathrm{\Omega }_1^2\mathrm{\Omega }_2^2/g^2(\mathrm{\Omega }_1^2+\mathrm{\Omega }_2^2)`$, for $`\mathrm{\Omega }g`$. The non-dark part of this mixture will be strongly coupled to the environment by photon scattering, but the dark part is unaffected by any process involving $`\mathrm{\Gamma }`$. Therefore the net loss of fidelity is given by the integral of $`P_{bb1}\kappa `$ over the switching time, which is $$p_2\frac{\mathrm{\Omega }^2}{2g^2}\kappa T.$$ (13) The net failure probability of the gate is $`p=p_1+p_2`$. We choose $`T`$ to minimise $`p`$, and then express the gate rate in terms of $`p`$ and the coupling parameters, obtaining: $$\frac{1}{T}\frac{1}{9}p^2\frac{g^2}{\kappa }.$$ (14) Note that we can obtain a gate of as high fidelity as we wish (against decoherence by the mechanisms we have considered) by slowing down the processor. The laser intensity must be chosen to give $`\mathrm{\Omega }=g^2(p/3)^{3/2}/\kappa `$ to meet the conditions of maximum fidelity. We used the approximation $`\mathrm{\Omega }g`$, which will be valid for this choice of $`\mathrm{\Omega }`$ when $`g/\kappa (3/p)^{3/2}`$. #### 5.2.2 Rabi flopping The method of adiabatic passage suffers from the problem that the gate rate given in equation (14) scales as the square of the failure probability, in contrast to the better scaling properties of equations (4) and (6). To avoid this poor scaling, the single-photon mode can be used in a way analogous to that adopted for the vibrational mode in section 5.1: instead of adiabatic passage in a dark state, we drive Rabi flopping on a chosen atom $`A`$, using a $`\pi `$ pulse to place a single photon in the cavity mode, and then a similar pulse on another atom $`B`$ swaps the information from the cavity mode into atom $`B`$. The analysis of this process is exactly as in equations (1)-(4) except that now $`g`$ is the coupling between atom and cavity mode, $`\kappa `$ is the relaxation rate of the cavity mode, and since the cavity mode only decays (it does not heat), the relaxation probability given in equation (2) is halved, so we replace $`\kappa `$ by $`\kappa /2`$ in equations (2)-(4), obtaining $$p=2\pi \sqrt{\frac{2\kappa \mathrm{\Gamma }}{g^2}},\frac{1}{T}=\frac{g^2}{\mathrm{\Gamma }}\frac{p}{8\pi ^2}.$$ (15) The gate rate is now limited by $$\frac{g^2}{\mathrm{\Gamma }}=\frac{3c\lambda ^2}{2\pi V},$$ (16) where we have used (7) and (8). From general physical principles, the minimum cavity mode volume $`V`$ might be expected to scale with wavelength as $`\lambda ^3`$, which would imply the processor runs faster at shorter wavelengths. In practice the cavity dimensions are also limited by technological considerations. The state of the art for high finesse Fabry Perot cavities in the optical domain is indicated by (see also ). A pair of mirrors of radius of curvature $`10`$ cm gave a finesse $`=4.2\times 10^5`$ at $`\lambda =852`$ nm, and was used to form a cavity of length $`L=44.6\mu \mathrm{m}`$, Gaussian mode waist $`w_0=20\mu \mathrm{m}`$, hence cavity field decay rate $`\kappa /(2\pi )=8`$ MHz. The mode volume $`V=Lw_0^2`$ yields $`g/(2\pi )=70`$ MHz for coupling to the D<sub>2</sub> line of atomic caesium, linewidth $`\mathrm{\Gamma }/(2\pi )=5.3`$ MHz. Equation (15) then gives $`p0.8`$. Another important type of optical cavity is provided by the whispering gallery modes of silica microspheres . Mabuchi and Kimble give the theory of coupling between the whispering gallery mode and an atom positioned at or near the surface of the sphere. On the surface of a $`50\mu \mathrm{m}`$ radius sphere, for example, the coupling to the D<sub>2</sub> line of neutral caesium atoms ($`\lambda =852`$ nm) is $`g/\mathrm{\Gamma }6`$ and for quality factor $`Q2\times 10^9`$ which has been reported , $`g/\kappa 174`$, giving $`p0.3`$. Combining (9), (16) and (15) we obtain $`p=(4\pi ^2w/\lambda )(3)^{1/2}`$, where $`w`$ is the average diameter of the mode, so that $`V=Lw^2`$. This implies that it will remain very difficult to obtain $`p1`$ for a considerable time, since great improvements in finesse will be needed, as well as a reduction in the mirror or sphere radius of curvature (to reduce the mode volume). #### 5.2.3 Discussion The advantage of the adiabatic passage method is that it allows precise gates even in the presence of relaxation of both the atomic excited state and the cavity mode. However, it is generally true that adiabatic methods achieve their greater degree of noise tolerance at the expense of processor speed. In the present case, if a cavity of sufficiently low decay rate can be built, then the Rabi flopping method may be preferable in that it will be faster in the limit of small $`p`$, and may perform better against further considerations, such as driving of off-resonant transitions. Since the physics of the vibration of an ion string is similar to that of the excitation of a cavity mode, it ought to be possible to apply a method analogous to the Mølmer-Sørensen one to the case of photon coupling. The essential result of the Mølmer Sørensen method is that the sensitivity to relaxation of the degree of freedom providing coupling ($`\kappa `$) is reduced by a factor $`M`$, while the gate time gets longer by the factor $`M`$. Examining (15) we see that in order to halve $`p`$, we would need to multiply the factor $`M`$ by four, increasing the gate time by a factor 4. The gate rate thus scales again as $`p^2`$, as it does in the adiabatic passage method. ### 5.3 Gate time per ion In the cases both of motional coupling and of photon coupling, the gate rate is reduced when the number of ions $`N`$ in the trap increases. For the motional coupling, this is partly because the ion string gets heavier and so has a reduced recoil frequency, and partly because if we wish to allow individual ions to be resolved (for single-ion addressing) the trap confinement must be reduced as more ions are added. For the photon coupling, the slow-down arises because the mode volume must be large enough to enclose all the ions, therefore reducing $`g`$. We will assume that the ions are in a linear string, or else a rectangular array, each separated from its neighbours by $`s=5\lambda `$. This means that if individual addressing is achieved by directing separate Gaussian laser beams on each ion, then the cross-talk between ions would be at the level $`10^4`$ if each beam had a waist $`w=s/(\mathrm{log}100)^{1/2}2.3\lambda `$. Such a beam can be produced by optics of modest numerical aperture. An alternative way to achieve the individual addressing is discussed by Leibfried . The scaling with $`N`$ of the motional coupling using the Cirac-Zoller method is discussed in . If we require the closest ions in a string to be separated by at least $`s=s(\lambda )`$, then the gate time increases approximately as $`N^{0.93}`$. Approximating this as proportional to $`N`$, and adopting the breathing mode (vibrational frequency $`\sqrt{3}`$ times that of the centre of mass mode) for the gates, we obtain from equation (6) a gate time per ion which depends only on the failure probability $`p`$ and properties of the ion such as its mass and recoil frequency. For candidate ions such as beryllium and calcium, this time is of order $`2`$ and $`10\mu `$s respectively for $`p=0.01`$, when $`s=5\lambda `$. The faster method of produces a gate rate $`\eta \omega _z/2\pi `$ limited through the limit on $`\omega _z`$ imposed by the need to keep ions separated by $`s`$. We take $`N=140`$ as an example, which will be useful in section 6. For the 313 nm transition in the beryllium ion the requirement $`s=5\lambda `$ leads to $`\omega _z/(2\pi )=418`$ kHz, hence $`\eta =0.088`$ and gate rate 37 kHz. The 397 nm transition in the calcium ion gives $`\omega _z/(2\pi )=141`$ kHz, $`\eta =0.056`$ and gate rate 8 kHz. Expressed as a gate time per ion, these examples are $`0.2`$ and $`0.9\mu `$s respectively, with gate failure probabilities of order $`p\eta ^2/20.004`$ and $`0.0016`$ respectively. For photon coupling we take as an example the cavity used in whose properties are summarised in section 5.2.2, with two changes: we reduce the mirror radius of curvature by a fifth (to $`2`$ cm), and we assume mirrors could be polished to provide the same finesse at half the wavelength. To be precise, we choose $`\lambda =493`$ nm which is appropriate for the barium ion. This ion can be readily laser cooled, and has the right kind of hyperfine structure for the adiabatic passage approach. The cavity mirrors are placed $`L=100\mu \mathrm{m}`$ apart, yielding a mode waist $`w_0=12.5\mu \mathrm{m}`$ and $`\kappa =(2\pi )3.6`$ MHz. The mode can therefore accommodate about 200 ions and equation (14) gives a gate time per ion of $`66`$ ns for $`p=0.01`$ (using $`\mathrm{\Gamma }/(2\pi )=11`$ MHz for the $`D`$ lines in the barium ion, we calculate $`g/(2\pi )=62`$ MHz). The cavity is illustrated in figure 2. The conclusion is that for optical cavities which are currently accessible or which may be expected in the near future, the CQED and motional methods have similar speeds (which one is faster will depend on $`p`$). It would be a lot harder to build the combined optical/atomic system compared to an ion trap alone. However, in principle the optical method could be much faster if suitable cavities could be made, and one possibility for this is a silica microsphere cavity. ## 6 Design for a quantum computer In view of the large number of technical problems still to be investigated in the laboratory, it is premature to try to design or build a large quantum computer. However, by sketching the main features of a possible design, we can learn about the issues, and identify avenues for further investigation. We aim to sketch a design for a quantum computer which could perform algorithms requiring $`10^6`$ Toffoli (controlled-controlled-not) gates on $`100`$ logical qubits. These numbers are chosen on the basis that about 100 qubits are likely to be needed to allow computations which could not be done on a classical computer. For example, the input to the computation might require 20 qubits, and the further 80 are needed as workspace. A problem needing less than 20 qubits input can probably be solved more easily on a large classical computer. It is important to count Toffoli (or equivalent) gates, not 2-qubit gates such as controlled-not, because Toffoli gates are a major component in any quantum algorithm which cannot be efficiently simulated classically, and, in contrast to controlled-not, it is non-trivial to implement them in a fault-tolerant manner . We take $`10^6`$ gates since Shor’s algorithm requires $`O(k^3)`$ gates for a $`k`$-qubit problem, and a similar scaling is likely to be involved for any useful algorithm. The computer will rely on fault-tolerant methods and quantum error correction. To be specific, we will adopt the methods described in : the quantum computer consists of blocks of 127 physical qubits, where each block encodes 29 logical qubits in a 7-error-correcting BCH code, and each data block is accompanied by several ancilliary blocks. We place each block in a separate processor, and link processors together by CQED and optical fibre methods . In order to allow low-level error correction methods in addition to those acting at the level of the encoded blocks, each processor will contain 13 extra physical qubits, making 140 in all. These also serve to implement communication protocols between processors, and for other tasks such as probing the local magnetic field. We will calculate the processor speed for two designs. The first is a linear ion trap containing calcium ions as described in section 5.3. The gates between ions in a given trap use the vibrational motion, and we adopt the fast gates offered by the lightshift-based concept of Jonathan et al. or by other methods which can be tailored in a similar fashion. To network between processors each ion trap has around it a Fabry-Perot optical cavity, with a mode shape overlapping several ions at the centre of the trap (see figure 3). Note that this cavity should have parameters optimised for quantum communication between traps , not for the optical gates discussed in section 5.2. The relevant transitions in calcium have wavelengths around 400 nm. We will also consider a more speculative idea: an all-optical method involving no ion traps. Instead, each processor is a single silica microsphere, with 140 caesium atoms positioned around the circumference of the sphere, either on the silica surface or trapped near it by a dipole force optical trap . The coupling between spheres is by further optical cavities whose design is left unspecified (it is not easy to see how to design them). The characteristics of the sphere are as given in section 5.2.2, except that to accommodate 140 atoms spaced by $`5\lambda `$ we require a sphere of radius $`63\mu \mathrm{m}`$ (for this calculation we use the wavelength in silica, since the atoms can be addressed by directing laser beams through the sphere). This reduces $`g`$ and $`\kappa `$ by a factor $`50/630.8`$ from the values quoted in section 5.2.2. In it was assumed that multiple controlled-not operations, in which there is one control qubit and many target qubits, could be performed in a single time-step. This will not be assumed possible here (but note the comments in ), so we must modify the results of accordingly. The preparation of ancillas is slowed down, which will increase the effect of memory noise. In order to reduce this problem, we provide the computer with more ancilla blocks. It can then prepare them in parallel, but staggered in time, so that enough ancillas are always available when they are needed. Providing 40 ancillas per data block (instead of 4 as in ) reduces the memory noise requirements by an order of magnitude. The whole computer then needs 138 ion traps and associated optical cavities for the data and ancilla blocks. Further optical cavities, with a few ions in each, may be useful for switching information paths, so that each block can communicate with most other blocks. This brings the total number of ion traps or microspheres and associated cavities to around 200. Using an analysis along the lines of that in , we find that the quantum algorithm can be stabilised as long as the failure rate per gate is $`\gamma 10^4`$, and the memory noise $`10^6`$ per bit per timestep. To obtain the gate failure probability $`10^4`$ it is unlikely that the best policy is to use the methods of section 5 alone. Instead, we will assume we can implement a low-level error correction tailored to the physical error process, as for example in . We assume the gates of precision $`10^4`$ can be built from gates of precision $`2\times 10^3`$ at a cost of a factor 10 in speed<sup>2</sup><sup>2</sup>2These ratios of speeds and failure probabilities are typical for a first-order (“single error”) correction process which can use reliable measurements.. For the ion trap, the gate time is therefore approximately $`10/8000=1250\mu `$s, using the figures given in section 5.3. The trap would use an axial vibrational frequency of 418 kHz for the centre of mass mode, and frequencies above 20 MHz for radial confinement. Vibrational heating would need to be at the level $`\kappa /\omega _z<10^6`$, which is within current achievements. For the microsphere the gate time is $`10\times 35=350\mu `$s (from (14) with $`g=4.8\mathrm{\Gamma }=160`$ MHz, $`\kappa (g/170)0.9`$ MHz). To prepare and verify an ancilla takes $`5000`$ time steps , but since we prepare ten in parallel for each one we need, the delay before the next one is available is approximately 500 time steps. We envisage that this time is also sufficient to carry out one inter-block controlled-not: that is, 127 controlled-nots at the physical qubit level between ions in separate traps (via the cavities and fibres). Therefore the time per correction of the whole computer is of order $`0.2`$ to $`0.6`$ s. We need about 8 such corrections per Toffoli gate in the logical algorithm , so the whole algorithm would take two to eight weeks to run. There are several methods which could be adopted to reduce this run time for the motional gates in the ion trap processor. The trap could be made tighter (increased $`\omega _z`$): the ions at the centre of the string would then be too close to be addressed individually, therefore one would use every other ion in this part of the string. There must be more ions in the trap, which off-sets the speed-up, but overall a factor 2 speed-up is readily obtained, and a higher factor if the spacer ions are of a lighter element. More demanding techniques offer further speed-up, for example by using more than one vibrational mode simultaneously, thus partially parallelizing the gates within a trap. Furthermore, fault-tolerant methods are based on synthesis of highly entangled states and make much use of controlled-multiple-not operations. Such operations can be generated in an ion trap in a time which scales as $`N^{1/2}`$ rather than $`N`$ as assumed above . Fabrication of the electrodes for hundreds of ion traps is fairly straightforward, using microfabrication methods or otherwise, as are the low-noise r.f. electronics to provide the trapping fields. For a detailed analysis of experimental issues for processing within each trap we refer the reader to and references therein. Whereas we have not given a thorough treatment of such issues here, we believe the parameters for trap tightness, optical addressing, and heating rates which we have assumed are reasonable. Less well understood is the phenomenon of charge build-up on the optical cavity mirrors, which will influence the operation of an ion trap, and techniques to prevent this may be essential. To build the mirrors for the optical cavities would be a major undertaking, but a possible one. One problem would be slow degradation of the mirrors during assembly or during loading of the traps. To place the mirrors and traps accurately together, and include associated optics such as optical fibres, would be taxing but feasible. The construction and operation of such a quantum computer would have more in common with the construction of a detector in high energy physics than with the manufacture of a classical computer chip; it is a lengthy, expensive and intricate process, but one whose results might merit the investment of resources, if serious uses for a 100-qubit computer can be found. The most important point is that it is conceivable. In conclusion, ion trap methods currently offer the only way to achieve multi-particle entangled states in a controllable way. They offer the prospect in the fairly near term of achieving various fundamental principles of quantum information physics, such as entanglement-enhanced communication and repeatable quantum error correction. It is clear that the controlled coupling of a trapped ion or neutral atom to a single-photon field in a high-finesse cavity also merits investigation for both quantum communication and quantum computing experiments. The rough sketch of a quantum computer design which we have given has two significant features: first, it is based on simple physical systems whose behaviour, including decoherence mechanisms, is well understood, and secondly it assumes only currently accessible levels of technology: all the components could be built now. The major unknowns are the detailed experimental issues which may arise when small optical cavities are combined with ion traps, whether the approximate error correction analysis gives a fair estimate of the noise tolerance, and whether suitable low-level error correction protocols, together with high-stability laser systems and magnetic and electric field noise suppression, can provide the $`10^4`$ gate precision and $`10^6`$ memory precision required throughout the computer. Further experiments, and more thorough feasibility studies, are certainly called for in this area. This work was supported by EPSRC, by Christ Church, Oxford and by the European Community network “QUBITS”. We thank D. Stacey for helpful comments on the manuscript.
warning/0004/cond-mat0004292.html	ar5iv	text	# Adsorption on carbon nanotubes: quantum spin tubes, magnetization plateaus, and conformal symmetry ## Abstract We formulate the problem of adsorption onto the surface of a carbon nanotube as a lattice gas on a triangular lattice wrapped around a cylinder. This model is equivalent to an XXZ Heisenberg quantum spin tube. The geometric frustration due to wrapping leads generically to four magnetization plateaus, in contrast to the two on a flat graphite sheet. We obtain analytical and numerical results for the magnetizations and transition fields for armchair, zig-zag and chiral nanotubes. The zig-zags are exceptional in that one of the plateaus has extensive zero temperature entropy in the classical limit. Quantum effects lift up the degeneracy, leaving gapless excitations which are described by a $`c=1`$ conformal field theory with compactification radius quantized by the tube circumference. Monolayer adsorption of noble gases onto graphite sheets has proven to be an interesting problem both theoretically and experimentally. Many of the observed features can be understood within a lattice gas model, where the underlying hexagonal substrate layer forms a triangular lattice of preferred adsorption sites. An equivalent formulation is in the language of spin models on a triangular lattice, where the repulsion between adsorbed atoms in neighboring sites translates into an antiferromagnetic Ising coupling. The frustration of the couplings by the triangular lattice leads to the rich phase diagram of the monolayer adsorption problem. Introducing hopping adds quantum fluctuations, further enriching the phase diagram. In this paper we address what happens if, in addition to the triangular lattice frustration, one has an extra geometric frustration due to periodic boundary conditions. In fact, such a system is physically realized by a single walled carbon nanotube, which may be viewed as a rolled graphite sheet. In this context, adsorption has been the subject of growing experimental and theoretical interest spurred by potential applications. Stan and Cole have considered the limit of non-interacting adatoms at low density, finding that they are localized radially near a nanotube’s surface at a distance comparable to that in flat graphite ($`3\text{Å}`$). In that work, it was sufficient to omit the hexagonal structure of the substrate. However, the corrugation potential selects the hexagon centers as additional commensurate localization points. In view of the similarity to flat graphite, we include both the substrate lattice and adatom interactions and consider a wider range of densities. In fact, very recently, it has been shown that the adsorbate stays within a cylindrical shell for fillings less than $`0.1/\text{Å}^2`$ (or $`0.5`$ adatom/hexagon), justifying the densities studied here. The adsorption sites thus form a triangular lattice wrapped around a cylinder. The geometrically inequivalent ways in which the wrapping is realized are labeled by two integers, $`(N,M)`$; the case $`(N,N)`$ is known as the armchair tube, $`(N,0)`$ is the zig-zag, and all others are chiral. The geometric frustration is present whenever the wrapping destroys the tripartite nature of the triangular lattice, which occurs when $`(NM)\mathrm{mod}\mathrm{\hspace{0.33em}3}`$ is non-zero (this criterion is familiar in the context of electronic conductivity). Thus, both zig-zag and chiral tubes can be frustrated geometrically, whereas armchair tubes cannot. In Fig. 1 we show the tube obtained from the $`(7,0)`$ zig-zag. In terms of the equivalent spin systems, this is similar to recent models of “spin tubes”. When the adsorbed gas is a hard-core boson, the lattice gas is defined by the Bose-Hubbard Hamiltonian $$=t\underset{ij}{}b_i^{}b_j+b_j^{}b_i+V\underset{ij}{}n_in_j\mu \underset{i}{}n_i,$$ (1) where $`n_i`$ is the boson density at site $`i`$, $`V`$ is the nearest neighbor repulsion and $`t`$ is the hopping amplitude. In the equivalent Heisenberg spin representation, $$=2t\underset{ij}{}S_i^xS_j^x+S_i^yS_j^y+V\underset{ij}{}S_i^zS_j^zH\underset{i}{}S_i^z,$$ (2) where $`S_i^z=n_i1/2`$ and $`H=\mu 3V`$ is an effective external magnetic field. Throughout the paper we will use the spin and density representations interchangeably. The Ising limit, in which hopping is not allowed, already contains many interesting features. We start the analysis in this regime, obtaining the phase diagram as a function of the magnetic field, and then consider quantum fluctuations perturbatively in $`t/V`$. We summarize our results first. The phase diagram in the temperature-magnetic field plane of a typical tube is shown in Fig. 2. When the index $`q=(NM)\mathrm{mod}\mathrm{\hspace{0.33em}3}`$ is $`1`$ or $`2`$, we find four lobes (solid lines), corresponding to two plateaus with magnetizations $`m_{}<1/3`$ and $`m_+>1/3`$. Here, we use the standard Ising notation in which spin is $`\pm 1`$. Note that the plateaus are real phases only at zero temperature because the tube is one-dimensional. At finite temperature, the boundaries should be interpreted as crossovers. Nonetheless, deep within a lobe, at $`k_BTV`$, the magnetizations are well-defined. Specifically, for $`q=1`$, we obtain the exact expressions $`m_+={\displaystyle \frac{1}{3}}\left(1+{\displaystyle \frac{2}{2M+N}}\right)m_{}={\displaystyle \frac{1}{3}}\left(1{\displaystyle \frac{2}{2N+M}}\right)`$ (3) $`H_c=\left(4{\displaystyle \frac{2M}{N+M}}\right)V`$ (4) The complementary case of $`q=2`$ is obtained by interchanging $`NM`$. On the other hand, those tubes without geometric frustration ($`q=0`$) behave similarly to the flat sheet (dotted lines) which has only two lobes with magnetizations $`\pm 1/3`$ . As the tube perimeter approaches the flat sheet limit, one expects that the geometric frustration becomes irrelevant. Indeed, as $`N`$ or $`M\mathrm{}`$, $`m_+`$ and $`m_{}`$ squeeze $`1/3`$ as the inverse of the tube diameter and become indistinguishable. Beyond the lobes, where the field is strong enough to overcome all nearest neighbor bonds ($`\|H\|/V>6`$ at $`k_BT=0`$), the tube is fully polarized. The filling fractions are obtained from the magnetizations by $`m=2(n1/2)`$. The phase diagram, however, is more easily visualized in terms of spin since spin reversal, $`mm`$, corresponds to particle-hole symmetry, $`n1n`$. We have verified this prediction numerically by transfer matrix methods for zig-zag tubes up to $`N=11`$ and for the chiral tubes up to $`N+M=7`$. Although $`7`$ is probably too small to be physical, we believe that the arguments in this paper generalize to any tube. In Fig. 3 we display sample data for two zig-zag tubes with different $`q`$: $`(7,0)`$ and $`(8,0)`$. The magnetization curves show clear plateaus whose values and transition fields match those predicted by Eq. (4). By increasing the temperature and following the evolution of the plateaus, we generate the phase diagram above. We find that a rather interesting feature of the zig-zag $`(N,0)`$ tubes emerges, making them exceptional. The insets in Fig. 3 indicate an extensive entropy at zero temperature, which has plateaus, too. Upon enumerating the degenerate space explicitly, we shall show that the entropy is exactly $`s=(\mathrm{ln2})/N`$ and that it occurs in $`m_+`$ for $`q=1`$ and in $`m_{}`$ for $`q=2`$. In the presence of hopping, the non-degenerate plateaus retain their gaps, whereas the degenerate ones become correlated states with a unique ground state and gapless excitations. More precisely, conformal invariance develops and the effective theory has central charge $`c=1`$ with a compactification radius, $`R`$, quantized by the tube circumference, $`R=N`$. In order to understand the magnetizations and nature of the geometric frustration, it is more intuitive to use the original bosonic picture. As a result of hard-core repulsion on the infinite graphite sheet, the $`m=1/3`$ plateau corresponds to filling one of the three sublattices, $`A,B`$ or $`C`$, of the triangular lattice. This configuration minimizes the repulsion, $`Vn_in_j`$, while maximizing the filling, $`\mu n`$. It is natural to try the same for nanotubes, as we illustrate in Fig. 4 for $`(5,0)`$. Upon wrapping, however, the thick vertical lines are identified and the lattice is no longer tripartite. In fact, the number of sublattice sites is no longer equal, and there is a mismatch along the thick line, which we term the “zipper”. On the left we fill the $`A`$ sublattice, obtaining the filling fraction $`n_+=2/5`$, and on the right either $`B`$ or $`C`$ may be filled with the result that $`n_{}=3/10`$. For general $`(N,0)`$ there are $`2N`$ hexagons in the unit cell, and the filling fractions are $`n_+=2N/3/2N`$ and $`n_{}=2N/3/2N`$, where $`x`$ and $`x`$ denote the larger and smaller of the two bounding integers of $`x`$, respectively. The magnetizations in Eq. (4) follow directly by using the correspondence $`m=2(n1/2)`$. Furthermore, due to the sublattice mismatch, the number density of adjacent particles, $`n_b`$, may be non-zero. In the case of $`(5,0)`$, there are two broken bonds per unit cell in $`n_+`$, and none in $`n_{}`$. This result generalizes to any $`q=2`$ zig-zag tube: $`n_{b+}=2/2N`$ and $`n_b=0`$. For $`q=1`$, the argument goes through as before, except that $`n_{b+}=1/2N`$. We summarize this compactly by $`n_{b+}=q/2N`$. Substituting these fillings into the Hamiltonian (1) yields two energies per site, $`e_\pm (\mu )=Vn_{b\pm }\mu n_\pm `$. The transition occurs when these levels cross: $`e_+=e_{}`$, or $$\frac{q}{2N}\mu \frac{2N/3}{2N}=\mu \frac{2N/3}{2N}$$ (5) Solving for $`\mu `$ and using the correspondence $`H=\mu 3V`$ gives precisely the critical field in Eq. (4). In particular, this explains why there are exactly two independent plateaus. Note that, for the special case of the zig-zags, the critical field depends only on $`q`$ and not on $`N`$ per se. In the above analysis, we have made only one assumption, namely that the zipper runs parallel to the tube axis. In general, the zipper may wind helically around the tube or wiggle sideways. However, in all the cases that we considered, the straight zipper has the lowest energy, and moreover, our transfer matrix computations, which are blind to this assumption, are consistent with our analysis. The chiral tubes are different. Due to their geometry the zipper is forced to wind, but, again, we find that the choice of the straightest possible zipper reproduces our numerics for $`N+M`$ up to $`7`$. The determination of the fillings and level crossings is much more involved than that of the zig-zag, and we leave it for a more detailed paper. In any case, our analysis reveals that the plateaus in a chiral tube are not macroscopically degenerate, so that the zig-zags are at a special degenerate point. Having understood in detail the Ising limit, we now turn on a small hopping, $`tV`$, that introduces quantum fluctuations. Deep within a plateau, the substrate is maximally filled since adding a particle increases $`n_b`$. Consequently, all plateaus begin with a classical gap of order $`V`$, and we work in the Hilbert space of the classical ground states. Those plateaus which have only a discrete symmetry must retain their gaps, but the macroscopically degenerate plateaus are more complicated. Let us reconsider the $`n_+`$ filling of the $`(5,0)`$ tube in Fig. 4. Notice that a particle may hop laterally by one site without changing $`n_b`$, as we illustrate in Fig. 5, left. Imagine building a typical $`n_+`$ state layer-by-layer from top to bottom, with a total of $`L`$ layers. Each new layer must add exactly two filled sites and one nearest-neighbor bond ($`n_b=1/5`$). This constraint implies that no two adjacent sites may be occupied within a layer; if they were, then, to conserve $`n_b`$, two adjacent sites must be occupied in the next, and so on up the tube. However, this state is not connected to any other by a single hop. Similarly, the particles cannot hop from layer to layer because this adds another intra-layer bond. An allowed state can be represented as a string of occupied sites, $`\{\sigma _i\}`$, $`i=1,\mathrm{},L`$, which in our example is $`\{\mathrm{}(5,3)(5,2)(3,1)(2,4)(1,4)(1,4)\mathrm{}\}`$. At each layer, there are exactly two possibilities for the following one. For example, $`(1,4)`$ can be followed by $`(1,4)`$ or by $`(2,4)`$. However, the total number of possibilities at any given level is five. Fig. 5 (right) summarizes this structure succinctly as a square lattice wrapped on the cylinder. A typical state, then, is a lattice path along the tube. Generalizing to $`(N,0)`$, we find $`N`$ possible states in each layer and two in the succeeding one, and the structure of states is again that of a wrapped square lattice with $`N`$ squares along the circumference. The dimension of the Hilbert space is the number of lattice paths, $`N2^L`$, so that in an infinitely long tube, the entropy per site is exactly $`(\mathrm{ln2})/N`$, as claimed earlier. Notice that constrained paths introduce correlations along the length of the tube, despite the absence of inter-layer hopping. The matrix elements of the projected Hamiltonian connect only those states that differ by a single hop: $$\{\tau \}\|\|\{\sigma \}=\{\begin{array}{c}2t\mathrm{if}_i\delta _{\sigma _i\tau _i}=L1\\ 0\mathrm{otherwise}\end{array}$$ (6) We diagonalize this Hamiltonian numerically with periodic boundary conditions for system sizes up to $`N=11`$ and $`L=10`$. Additionally, we can obtain the ground state energy up to $`L=16`$ due to the sparseness of $``$. We will fix $`2t=1`$ in what follows. We find that the degeneracy is lifted and the ground state becomes unique and uniform. The ground state energy, $`E_0(L)`$, follows $`E_00.607L\pi c/6`$ with $`c1.005`$. The lowest $`N1`$ excited states are given by $`\mathrm{\Delta }_a=a^2\mathrm{\Delta }/(N^2L)`$, with $`\mathrm{\Delta }=12.9\pm 0.5`$, which is shown in Fig. 6 for $`a=1,2,3`$. All of these levels are doubly degenerate. This ground state energy and spectrum are in perfect agreement with a conformally invariant bosonic theory with central charge $`c=1`$ compactified on a radius $`R=\zeta N`$. We take the Lagrangian density $`=\frac{1}{8\pi }[v^1(_t\varphi )^2v(_x\varphi )^2]`$, where $`v`$ is a velocity, and the compactification is defined by $`\varphi \varphi +2\pi R`$. The zero mode energies of this theory are $$E_{a,b}^0=\frac{2\pi v}{L}\left(\frac{a^2}{R^2}+\frac{b^2R^2}{4}\right),$$ (7) where $`a,b`$ are integers that label topological momenta and windings of $`\varphi `$. Right- and left-moving oscillator modes of energy $`\omega _n=vk_n`$, where $`k_n=2\pi n/L`$, also appear in the spectrum, but for $`a<N`$ the zero modes are the lowest. Our spectrum in Fig. 6 corresponds to $`E_{a,b}^0`$ with $`b=0`$. To fix $`\zeta `$, we look at higher low-lying levels (which also scale like $`1/L`$). We find that the $`N`$’th excitation energy is independent of $`N`$ and quadruply degenerate. This can happen only if the $`N`$’th zero mode, $`E_{\pm N,0}^0=2\pi v/\zeta ^2L`$, is degenerate with the lowest oscillator mode, $`\omega _{\pm 1}=2\pi v/L`$, which fixes $`\zeta =1`$. Thus, the compactification radius is $`R=N`$. The velocity can be read off from the slopes in Fig. 6 as $`v=\mathrm{\Delta }/2\pi `$. The rest of our spectrum is consistent with these parameters. One observable consequence of conformal symmetry is that the low temperature heat capacity is fixed by $`c`$: $$C=c\frac{\pi k_B^2}{3}T=\frac{\pi k_B^2}{3}T$$ (8) It is noteworthy that, even though the dispersion of the oscillator modes is independent of $`N`$, the spectrum remembers, via the zero-modes, the finite radius of the nanotube. Furthermore, $`R`$ is quantized by $`N`$; in the language of Luttinger liquids, this means that the Luttinger parameter is fixed by topology, similarly to the case of edge states in a fractional quantum Hall fluid, and in contrast to quantum wires (where the Luttinger parameter can vary continuously). Because there is no inter-layer hopping, $`\varphi `$ is tied to transverse, rather than to longitudinal, density fluctuations along the tube. We will present a detailed analytical derivation of the effective theory from the lattice in Fig. 5 elsewhere. Let us briefly view the spin tube as a quantum spin ladder to see if it yields the zero gap. A standard approach is to use a Lieb-Schultz-Mattis (LSM) argument, in which the spins are deformed slowly along the length. Applying it to our tube, we find that a plateau is gapless if $`SM`$ is not an integer, where $`S`$ is the total spin and $`M`$ the magnetization per layer. Using $`S=N/2`$ and the magnetizations from Eqn. (4), we find that $`SM`$ is an integer in the macroscopically degenerate plateaus, so that the LSM argument is insufficient in this case. A conclusive argument must take the geometric frustration into account. Before concluding, we should point out that the geometry of the $`(2,0)`$ tube is special; all sites in adjacent layers are interconnected. As a result, all of its plateaus have an extensive entropy, and we find that hopping opens a gap in both plateaus. In fact, this tube can be written as a spin chain that has been studied at isotropic coupling, $`2t=V`$. Two plateaus were found in this case, and it is tempting to speculate whether the two regimes are connected adiabatically. In conclusion, we have studied the problem of monolayer adsorption on carbon nanotubes and identified several interesting filling fraction plateaus. Since the difference between the plateaus decreases slowly, as the inverse of the tube diameter, experimental measurement should be feasible for large enough tubes. We have identified the zig-zag tubes as exceptional, in which the geometric frustration together with quantum fluctuations lead to conformal symmetry. This system is a physical realization of quantum spin tubes. The authors wish to thank C. Buragohain, M. El-Batanouny, E. Fradkin, N. Read, C. Nayak, M. Vojta, and X.-G. Wen for helpful comments. Support was provided by the NSF Grant DMR-98-18259(D. G.), DMR-98-76208 and the Alfred P. Sloan Foundation (C. C.).
warning/0004/hep-ex0004033.html	ar5iv	text	# Measurement of the product branching fraction ℬ(𝑐→Θ_𝑐𝑋)⋅ℬ(Θ_𝑐→Λ𝑋) ## I Introduction Inclusive measurements of charmed baryon decay products provide essential information on the relative contributions of different decay processes (e.g., external W-emission, internal W-emission, W-exchange) to the weak $`csW^+`$ transition in baryons. Difficulties in distinguishing direct charm decay products from jet fragmentation particles have hampered such inclusive measurements. An example is the measurement of $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)`$ (although $`\mathrm{`}\mathrm{`}\mathrm{\Lambda }_c^{\prime \prime }`$ here designates $`\mathrm{\Lambda }_c^+`$, charge conjugation is implicit throughout). Using the total $`\mathrm{\Lambda }`$ yield at $`\sqrt{s}`$ = 10 GeV/c to measure $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)`$ requires separating the $`\mathrm{\Lambda }`$ component due to light quark fragmentation from production via $`c\mathrm{\Lambda }_c\mathrm{\Lambda }X`$. The difficulty of separating fragmentation $`\mathrm{\Lambda }`$’s from those resulting from $`\mathrm{\Lambda }_c`$ decays can be overcome by using a tagged sample of $`e^+e^{}c\overline{c}`$ events. In this analysis, we use charm-event tagging to measure the product of the likelihood for a charm quark to materialize as a charmed baryon $`\mathrm{\Theta }_c`$ times the branching fraction for a charmed baryon to decay into a $`\mathrm{\Lambda }`$: $`(c\mathrm{\Theta }_cX)(\mathrm{\Theta }_c\mathrm{\Lambda }X)`$. From JETSET 7.3 Monte Carlo simulations, using the default LEP-tuned control parameters, we expect that $``$88% of $`\mathrm{\Theta }_c`$ particles produced in $`e^+e^{}c\overline{c}`$ fragmentation at $`\sqrt{s}`$=10 GeV will be $`\mathrm{\Lambda }_c`$’s. Previous measurements of $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)`$ have either measured the increase in the $`\mathrm{\Lambda }`$ production rate as the $`e^+e^{}\mathrm{\Lambda }_c\overline{\mathrm{\Lambda }}_c`$ threshold is crossed, topologically tagged $`\mathrm{\Lambda }_c`$ decays, or derived a value for $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)`$ based on measurements of baryon production in B-meson decay. Knowing, for example, that $`(Bp/\overline{p}(direct)+anything)=5.5\pm 0.5`$%, $`(B\mathrm{\Lambda }/\overline{\mathrm{\Lambda }}+anything)=4.0\pm 0.5`$%, and assuming that $`\overline{B}\mathrm{\Lambda }_c\overline{N}`$X dominates baryon production in B-decay, with $`\overline{N}`$ equally likely to be $`\overline{p}`$ or $`\overline{n}`$, one can estimate: $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)\frac{4}{5.5\times 2}`$=36%. This simple-minded estimate, however, needs to be modified to take into account many corrections, among them the recent result that $`\frac{(\overline{B}\overline{\mathrm{\Lambda }}X)}{(\overline{B}\mathrm{\Lambda }X)}=0.43\pm 0.09\pm 0.07`$. That latter result implies that only $``$2/3 of the inclusive ($`\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}`$) yield in $`\overline{B}`$-decay come from decays of charmed baryons; the remainder is presumably due to associated production or decay of anticharm baryons. In this Article, we use a new technique to determine the product of the probability for a charm quark to produce a charmed baryon $`\mathrm{\Theta }_c`$ times the probability that the charmed baryon will decay into a $`\mathrm{\Lambda }`$: $`_\mathrm{\Lambda }=(c\mathrm{\Theta }_cX)(\mathrm{\Theta }_c\mathrm{\Lambda }X`$), using continuum $`e^+e^{}`$ annihiliation events at $`\sqrt{s}`$=10.55 GeV. A sample cc̄ event is schematically shown below, showing some of the particles relevant to our measurement. In the event depicted below, the fragmentation of the original cc̄ quark-antiquark results in a $`\mathrm{\Lambda }`$ recoiling in one hemisphere opposite the anti-charm tag (either the soft pion or the electron) in the other hemisphere. For this analysis, event “hemispheres” are defined using the axis which minimizes the momentum (charged plus neutral) transverse to that axis (the “thrust” axis). Note that both the soft pion and the electron are of charge opposite to the $`p`$ daughter of the $`\mathrm{\Lambda }`$. In addition to the tags depicted below, we also tag cc̄ events with fully reconstructed $`\overline{D}^0K^+\pi ^{}`$ or $`D^{}K^+\pi ^{}\pi ^{}`$ events, in which the $`\overline{D}`$ daughter kaon contains an s̄-quark, in contrast to the $`\mathrm{\Lambda }`$. For all four tags, we will therefore refer to our signal as an opposite hemisphere, opposite sign (OH/OS) correlation. $`\begin{array}{cccccccc}& & & c& \overline{c}& & & \\ & & \mathrm{\Lambda }_c& & & D^{}& & \\ & \mathrm{\Lambda }X& & & & & \pi _{soft}^{}+\overline{D^0}& \\ p+\pi ^{}& & & & & & & e^{}+\overline{\nu }_e+K^+\end{array}`$ We note that the $`\mathrm{\Lambda }_c`$ in this example can, in practice, be any charmed baryon and will from here on be denoted as “$`\mathrm{\Theta }_c`$”. ## II Apparatus and Event Selection This analysis was performed using the CLEO II detector operating at the Cornell Electron Storage Ring (CESR) at center-of-mass energies $`\sqrt{s}`$ = 10.52–10.58 GeV. The CLEO II detector is a general purpose solenoidal magnet spectrometer and calorimeter designed to trigger efficiently on two-photon, tau-pair, and hadronic events . Measurements of charged particle momenta are made with three nested coaxial drift chambers consisting of 6, 10, and 51 layers, respectively. These chambers fill the volume from $`r`$=3 cm to $`r`$=1 m, with $`r`$ the radial coordinate relative to the beam ($`\widehat{z}`$) axis. This system is very efficient ($`ϵ`$98%) for detecting tracks that have transverse momenta ($`p_T`$) relative to the beam axis greater than 200 MeV/c, and that are contained within the good fiducial volume of the drift chamber ($`\|\mathrm{cos}\theta \|<`$0.94, with $`\theta `$ defined as the polar angle relative to the beam axis). This system achieves a momentum resolution of $`(\delta p/p)^2=(0.0015p)^2+(0.005)^2`$ ($`p`$ is the momentum, measured in GeV/c). Pulse height measurements in the main drift chamber provide specific ionization resolution of 5.5% for Bhabha events, giving good $`K/\pi `$ separation for tracks with momenta up to 700 MeV/c and separation of order 2$`\sigma `$ in the relativistic rise region above 2 GeV/c. Outside the central tracking chambers are plastic scintillation counters, which are used as a fast element in the trigger system and also provide particle identification information from time of flight measurements. Beyond the time-of-flight system is the electromagnetic calorimeter, consisting of 7800 thallium-doped CsI crystals. The central “barrel” region of the calorimeter covers about 75% of the solid angle and has an energy resolution which is empirically found to follow: $$\frac{\sigma _\mathrm{E}}{E}(\%)=\frac{0.35}{E^{0.75}}+1.90.1E;$$ (1) $`E`$ is the shower energy in GeV. This parameterization includes effects such as noise, and translates to an energy resolution of about 4% at 100 MeV and 1.2% at 5 GeV. Two end-cap regions of the crystal calorimeter extend solid angle coverage to about 95% of $`4\pi `$, although energy resolution is not as good as that of the barrel region. The tracking system, time of flight counters, and calorimeter are all contained within a superconducting coil operated at 1.5 Tesla. Flux return and tracking chambers used for muon detection are located immediately outside the coil and in the two end-cap regions. The event sample used for this measurement is comprised of 3.1 $`fb^1`$ of data collected at the $`\mathrm{{\rm Y}}`$(4S) resonance and 1.6 $`fb^1`$ of data collected about 60 MeV below the $`\mathrm{{\rm Y}}`$(4S) resonance. Approximately $`5\times 10^6`$ continuum cc̄ events are included in this sample. For our analysis, we select continuum hadronic events which contain either a low-momentum pion $`\pi _{soft}^{}`$ emitted at small angles relative to the event thrust axis (from $`D^{{}_{}{}^{}}\overline{D}^0\pi _{soft}^{}`$), a high momentum electron (from $`\overline{c}\overline{s}e^{}\nu _e`$), or a fully reconstructed $`\overline{D}^0K^+\pi ^{}`$ or $`D^{}K^+\pi ^{}\pi ^{}`$ as a tag of $`e^+e^{}c\overline{c}`$ events. In order to suppress background and enrich the hadronic fraction of our event sample, we impose several event requirements. Candidate events must have: (1) at least five detected, good quality, charged tracks; (2) an event vertex consistent with the known $`e^+e^{}`$ interaction point; (3) a total measured visible event energy, equal to the sum of the observed charged plus neutral energy $`E_{vis}(=E_{chrg}+E_{neutral}`$), greater than 110% of the single beam energy, $`E_{vis}`$ $`>`$ 1.1 $``$ $`E_{beam}`$. In addition, when using an electron to tag a cc̄ event we require that either the beam energy $`E_{beam}`$ be less than 5.275 GeV or that the event be well collimated. Specifically, the ratio of Fox-Wolfram event shape parameters $`H2/H0`$ can be used to quantify the “jettiness” of an event. For a perfectly spherical flow of event energy, this ratio equals 0; for a perfectly jetty event, this ratio equals 1.0. For our electron tags, we require this ratio to be greater than 0.35. This requirement is necessary to remove contamination from B$`\overline{\mathrm{B}}`$ events. Similarly, when using $`\overline{D}^0K^+\pi ^{}`$ or $`D^{}K^+\pi ^{}\pi ^{}`$ as our charm tags, we eliminate B$`\overline{\mathrm{B}}`$ background by requiring that the reconstructed $`\overline{D}`$ momentum exceed 2.3 GeV/c. ## III Tag Identification ### A Electron Tags To suppress background from fake electrons, as well as true electrons not necessarily associated with $`e^+e^{}`$ cc̄ events, we require that our electron-tag candidates satisfy the following criteria: (a) The electron must pass a strict “probability of electron” identification criterion. This identification likelihood combines measurements of a given track’s specific ionization deposition in the central drift chamber with the ratio of the energy of the associated calorimeter shower to the charged track’s momentum. True electrons have shower energies approximately equal to their drift chamber momenta; hadrons tend to be minimum ionizing and have considerably smaller values of shower energy relative to their measured momenta. We require that the logarithm of the ratio of a charged track’s electron probability relative to the probability that the charged track is a hadron $`𝙻_e`$ be greater than 7 ($`𝙻_e7`$). (b) The momentum of the electron must be greater than 1 GeV/c. This criterion helps eliminate kaon and pion fakes and also suppresses electrons from photon conversions ($`\gamma e^+e^{}`$) and $`\pi ^0`$ Dalitz decays ($`\pi ^0\gamma e^+e^{}`$). (c) The electron must have an impact parameter relative to the event vertex less than 4 mm along the radial coordinate and no more than 2 cm along the beam axis. This provides additional suppression of electrons resulting from photon conversions. ### B Soft pion tags Our soft pion tag candidates must pass the following restrictions: (a) The pion must have an impact parameter relative to the event vertex less than 5 mm along the radial coordinate and no more than 5 cm along the beam axis. (b) The pion must pass a 99% probability criterion for pion identification, based on the associated charged track’s specific ionization measured in the drift chamber. (c) The pion’s measured momentum must be between 0.15 GeV/c and 0.40 GeV/c. (d) The pion’s trajectory must lie near the trajectory of the parent charm quark, as expected for pions produced in $`D^{}\overline{D}^0\pi _{soft}^{}`$. Experimentally, this is checked using the variable sin$`{}_{}{}^{2}\theta `$, with $`\theta `$ the opening angle between the candidate soft pion and the event thrust axis. Assuming that the thrust axis approximates the original cc̄ axis, true $`\pi _{soft}^{}`$ should populate the region $`\mathrm{sin}^2\theta 0`$. Figure 1 displays the sin$`{}_{}{}^{2}\theta `$ distribution for candidates passing our event and track selection criteria. The excess in the region $`\mathrm{sin}^2\theta `$ near 0 constitutes our charm-tagged sample. The fit includes a signal contribution, the shape of which is determined from Monte Carlo simulations, and a lower-order polynomial to fit the background. The technique for determining the signal shape and background follows that of an earlier CLEO analysis, which used this method to measure the branching fraction $`(D^0K^{}\pi ^+)`$. ### C $`\overline{D}^0K^+\pi ^{}`$ and $`D^{}K^+\pi ^{}\pi ^{}`$ tags The $`\overline{D}`$-tagged analysis was performed independent of the $`\pi _{soft}^{}`$ and electron-tagged analyses. For this latter analysis, we take advantage of improved track and particle reconstruction algorithms, which were unavailable when the $`\pi _{soft}^{}`$ electron-tagged analyses were conducted. Also, in order to compensate for the intrinsically smaller efficiency of $`\overline{D}^0K^+\pi ^{}`$ and $`D^{}K^+\pi ^{}\pi ^{}`$ reconstruction, we also use a three-fold larger data sample (13.1 fb<sup>-1</sup>, including the CLEO II.V data set) for this analysis. Note that, for the purpose of this analysis, which does not utilize the precision vertexing afforded by the CLEO II.V silicon vertex system, the essential detector performance characteristics are the same as for the CLEO II data sample. $`\overline{D}^0K^+\pi ^{}`$ and $`D^{}K^+\pi ^{}\pi ^{}`$ tags are fully reconstructed from kaon and pion candidates as follows: (a) The kaon and pion candidates must have impact parameters relative to the event vertex less than 5 mm along the radial coordinate and no more than 5 cm along the beam axis. (b) Both the pion and kaon tracks must be consistent with their assumed particle identities at the level of 2.5 standard deviations ($`\sigma `$), using the available specific ionization and time-of-flight particle identification information. (c) Both the pion and kaon must have momentum greater than 0.3 GeV/c. (d) The fully reconstructed $`\overline{D}`$ meson tag must have momentum greater than 2.3 GeV/c to eliminate BB̄ backgrounds. ## IV Lambda Detection After finding a charm tag, we reconstruct $`\mathrm{\Lambda }p\pi `$ in the hemisphere opposite the tag. In addition to a 99% particle identification probability requirement placed on both the daughter proton and pion used in reconstructing the $`\mathrm{\Lambda }`$, we also require that candidate $`\mathrm{\Lambda }`$ particles have momenta greater than 1 GeV/c and that the lambda vertex be located at least 2 cm away from the $`e^+e^{}`$ collision point in the radial direction. According to Monte Carlo simulations (Figure 2), imposing the minimum $`\mathrm{\Lambda }`$ momentum requirement (p$`{}_{\mathrm{\Lambda }}{}^{}>1.0`$ GeV) in a charm-tagged event passing our event selection requirements results in a $`\mathrm{\Lambda }`$ sample which is $`>`$95% pure $`\mathrm{\Theta }_c\mathrm{\Lambda }X`$, with the remaining $`\mathrm{\Lambda }`$’s due to light quark fragmentation. ## V Yields To extract our signal yields in the lepton-tagged sample, we plot the proton-pion invariant mass for $`\mathrm{\Lambda }`$ candidates in electron-tagged events. Figure 3 shows the candidate $`\mathrm{\Lambda }`$ mass separated into each of the four possible sign/hemisphere correlations. Our candidate signal $`\mathrm{\Lambda }`$’s are contained in Figure 3d (lower right). The number of signal $`\mathrm{\Lambda }`$’s is extracted by fitting a Gaussian $`\mathrm{\Lambda }`$ signal function plus a second-order Chebyschev polynomial background. To determine the $`\mathrm{\Lambda }`$ yield in $`\pi _{soft}^{}`$ tagged events, we plot the $`\mathrm{sin}^2\theta `$ of the $`\pi `$/thrust axis angle for each candidate $`\mathrm{\Lambda }`$ found vs. the candidate $`p^+\pi ^{}`$ mass. We then project the resulting histogram onto the $`\mathrm{sin}^2\theta `$ axis and fit the peak at $`\mathrm{sin}^2\theta 0`$ for the case where the $`p^+\pi ^{}`$ invariant mass is in the $`\mathrm{\Lambda }`$ region (signal) versus the case where the $`p^+\pi ^{}`$ invariant mass is in the $`\mathrm{\Lambda }`$ sidebands. After performing a sideband subtraction in $`\mathrm{\Lambda }`$ mass of the two $`\mathrm{sin}^2\theta `$ distributions, we obtain Figure 4. A fit to the $`\mathrm{sin}^2\theta `$ distribution for all pion candidates (Figure 1) determines our total number of $`\pi _{soft}^{}`$-tagged cc̄ events (the denominator in our ratio). Monte Carlo simulations indicate that both (i) the fraction of non-$`c\overline{c}`$ tags and (ii) the fraction of candidate signal $`\mathrm{\Lambda }`$’s that do not originate from $`\mathrm{\Theta }_c`$ decays but pass our selection criteria are small ($`<`$3%, see Figure 2). We test the overall accuracy of the Monte Carlo by comparing same-sign, opposite hemisphere correlation events (“SS/OH”, i.e. opposite of the sign correlation expected for signal) in simulations compared to data. We find that the ratio of SS/OH electron-$`\mathrm{\Lambda }`$ correlation events to the number of “right sign” (OS/OH) signal events is $`0.19\pm 0.05`$ in Monte Carlo and $`0.19\pm 0.07`$ in data. The corresponding values for the $`\pi _{soft}^{}\mathrm{\Lambda }`$ correlations are 0.09 and 0.16, respectively. Within statistics, the Monte Carlo reproduces the “wrong sign” (SS/OH) fractions observed in data. Nevertheless, we conservatively assign a relative systematic error of 10% (19%/2) to reflect our confidence in the simulations. This value is entered in the final systematic errors table (Table IV) as “Event generator mismodeling”. For the $`\overline{D}`$-tagged sample, the signal is extracted from a two-dimensional plot of $`M_{\mathrm{\Lambda }_c}`$ (the mass of the $`\mathrm{\Lambda }_c`$ candidate) vs. $`M_{\overline{D}}`$ (the mass of the $`\overline{D}`$-tag candidate, either $`M_{K^+\pi ^{}}`$ or $`M_{K^+\pi ^{}\pi ^{}}`$), as indicated in Figures 5 and 6. Our signal comprises events which contain both a fully reconstructed $`\overline{D}`$ and also a $`\mathrm{\Lambda }`$ (“double-tags”). The double-tag signal yields are determined by a two-dimensional sideband subtraction technique, similar to that used to determine the signal yield for the soft-pion tagged sample. Here, we subtract the scaled $`\mathrm{\Lambda }`$ yield in the $`\overline{D}`$ sideband region from the $`\mathrm{\Lambda }`$ yield in the $`\overline{D}`$ signal region. The resulting excess, is, by definition, our double-tag signal. As a check of the signal extraction, the yield for the “wrong-sign” double-tag signal (i.e., $`D^0K^{}\pi ^+`$ vs. $`D^0K^{}\pi ^+`$) is similarly extracted using the same subtraction. In such events, the expected true correlated signal should be negligible and non-zero only through doubly Cabibbo suppressed decays; in fact, we find $`15\pm 25`$ events from a two-dimensional $`M(K^{}\pi ^+)`$ vs. $`M(K^{}\pi ^+)`$ plot, and $`115\pm 157`$ events from the two-dimensional $`M(K^{}\pi ^+\pi ^+)`$ vs. $`M(K^{}\pi ^+\pi ^+)`$ plot. ## VI Calculations The product branching fraction $`(c\mathrm{\Theta }_c)(\mathrm{\Theta }_c\mathrm{\Lambda }X`$) can be derived from the fraction of times that an event containing a charm tag in one hemisphere contains a $`\mathrm{\Lambda }`$ in the opposite hemisphere. This, effectively, is the fraction of $`\mathrm{\Lambda }`$ particles per cc̄ event and should equal the probability of a $`c`$ quark fragmenting to produce a $`\mathrm{\Theta }_c`$ multiplied by the probability of the $`\mathrm{\Theta }_c`$ to decay to a $`\mathrm{\Lambda }`$ multiplied by the efficiency for detection of a $`\mathrm{\Lambda }`$ in our charm-tagged event sample. In equation form, defining $`\frac{N(\mathrm{\Lambda })}{c\overline{c}}`$ as the ratio of the number of reconstructed $`\mathrm{\Lambda }`$’s in tagged $`c\overline{c}`$ events to the total number of $`c\overline{c}`$ event tags, we have: $$\frac{N(\mathrm{\Lambda })}{c\overline{c}}=(c\mathrm{\Theta }_c)(\mathrm{\Theta }_c\mathrm{\Lambda }X)ϵ_{\mathrm{\Lambda },tagged}$$ (2) for both data and Monte Carlo. Assuming that the Monte Carlo simulation accurately reproduces the efficiency for finding a $`\mathrm{\Lambda }`$ in a tagged event ($`ϵ_{\mathrm{\Lambda },tagged}`$ in this equation), the yield of non cc̄ tags ($`<`$4%), and the fraction of non-signal $`\mathrm{\Lambda }tag`$ correlations ($`<`$5%), we can then calibrate our observed value of $`\mathrm{\Lambda }`$’s per cc̄ in data to Monte Carlo: $$\frac{\frac{N(\mathrm{\Lambda })}{c\overline{c}}^{Data}}{\frac{N(\mathrm{\Lambda })}{c\overline{c}}^{MC}}=\frac{(c\mathrm{\Theta }_cX)(\mathrm{\Theta }_c\mathrm{\Lambda }X)^{Data}}{(c\mathrm{\Theta }_cX)(\mathrm{\Theta }_c\mathrm{\Lambda }X)^{MC}};$$ (3) the Monte Carlo values for $`(c\mathrm{\Theta }_cX)`$ and $`(\mathrm{\Theta }_c\mathrm{\Lambda }X)`$ are 0.0667 and 0.369, respectively. (A recent measurement by the ALEPH collaboration at $`\sqrt{s}`$=90 GeV has determined $`(c\mathrm{\Lambda }_cX)=0.079\pm 0.008\pm 0.004\pm 0.020`$ \[the last systematic error represents the uncertainty in the $`\mathrm{\Lambda }_cpK^{}\pi ^+`$ branching fraction\], although it is not clear how appropriate this value is for $`\sqrt{s}`$=10 GeV.) Note that the efficiency $`ϵ_{\mathrm{\Lambda },tagged}`$ is tag-dependent - due to geometric and momentum correlations from hemisphere to hemisphere, we expect the highest efficiency for the $`\overline{D}^0`$ and $`D^{}`$ tags, followed by soft pion tags and electron tags. A summary of our yields and calculations for $``$($`c\mathrm{\Theta }_cX`$) $``$ $``$($`\mathrm{\Theta }_c\mathrm{\Lambda }X`$) is presented in Tables I and II. Presented in those Tables are our raw yields, the number of true electrons which do not tag cc̄ events (‘fake tags’), and the number of $`\mathrm{\Lambda }`$’s reconstructed in the opposite hemisphere for our electron-tagged, soft pion-tagged, and $`\overline{D}`$-tagged samples. Backgrounds in the electron-tagged sample from BB̄ and $`\tau \overline{\tau }`$ events are estimated from a large sample of Monte Carlo events, using a CLEO event generator for B decays, and KORALB for $`\tau \overline{\tau }`$ decays. The electron background from $`\gamma \gamma `$ events is estimated from the forward-backward excess of positrons versus electrons, compared to the expectation from QED. Our yields correspond to $`_\mathrm{\Lambda }`$ = (1.62 $`\pm `$ 0.10)% using electrons to tag cc̄ events, $`_\mathrm{\Lambda }`$ = (1.53 $`\pm `$ 0.06)% using $`\pi _{soft}^{}`$ to tag cc̄ events, $`_\mathrm{\Lambda }`$ = (2.12 $`\pm `$ 0.09)% using $`\overline{D}^0`$ to tag cc̄ events, and $`_\mathrm{\Lambda }`$ = (2.09 $`\pm `$ 0.13)% using $`D^{}`$ to tag cc̄ events (statistical errors only). ## VII Cross Checks We have conducted two cross-checks to verify the accuracy of our derived result for $`_\mathrm{\Lambda }`$. We emphasize that these are not measurements in themselves (and therefore have no quoted systematic errors), but are presented only to verify our $`_\mathrm{\Lambda }`$ measurement. ### A $`D^0K^{}\pi ^+`$ decays. As a first cross-check, we compare the data- versus Monte Carlo-derived values for the product branching fraction: $`(cD^0)(D^0K^{}\pi ^+)`$, using charm-tagging. Since the branching fraction for $`D^0K^{}\pi ^+`$ is known precisely, and since the fractional uncertainty in $`(cD^0)`$ is expected to be smaller than the corresponding uncertainty in $`(c\mathrm{\Lambda }_c)`$, we can compare the value of $`(cD^0)(D^0K^{}\pi ^+)`$ measured with charm-tagging in data versus Monte Carlo simulations and thereby verify the method used in the $`\mathrm{\Lambda }`$ measurement. Using the same charm-tagged sample as before, we therefore search for the decay $`D^0K^{}\pi ^+`$ (using the same requirements mentioned before) opposite the tag rather than $`\mathrm{\Lambda }p\pi ^{}`$. As before, we perform a sideband subtraction to determine the number of $`D^0K^{}\pi ^+`$ decays in our $`\pi _{soft}^{}`$ charm tagged sample. We thus use the same equation as with our $`\mathrm{\Lambda }`$ analysis, only modified for the $`D^0K^{}\pi ^+`$ decay mode: $$\frac{N(K^{}\pi ^+)}{c\overline{c}}=(cD^0)(D^0K^{}\pi ^+)ϵ_{D^0K^{}\pi ^+}$$ (4) Again, assuming that the Monte Carlo accurately reproduces the efficiency for finding a $`D^0`$ decay in a tagged event, we calibrate our observed value of $`D^0`$’s per cc̄ in data to Monte Carlo: $$\frac{\frac{N(K\pi )}{c\overline{c}}^{Data}}{\frac{N(K\pi )}{c\overline{c}}^{MC}}=\frac{(cD^0)(D^0K\pi )^{Data}}{(cD^0)(D^0K\pi )^{MC}}$$ (5) The results of our $`D^0`$ cross-check are presented in Table III. The Monte Carlo adequately reproduces the $`D^0K^{}\pi ^+`$ yield per $`\overline{c}`$-tag. Based on the consistency between these values and the known $`D^0K^{}\pi ^+`$ branching fraction, a scale factor is applied to the data and a systematic error is added which reflects only the statistical precision of this cross check. In all cases, the scale factor (Table III) is consistent with unity. ### B $`\mathrm{\Lambda }_c\mathrm{\Lambda }e\nu _e`$ decays. A second cross-check is afforded by our $`\mathrm{\Lambda }`$-electron correlation sample. We note that the same-hemisphere, same-sign events are expected to be dominated by $`\mathrm{\Lambda }_c\mathrm{\Lambda }e^+\nu _e`$ decays. Since the $`\mathrm{\Lambda }_c\mathrm{\Lambda }e\nu _e`$ branching fraction has been measured, we can use the relative ratio of the same-hemisphere, same-sign $`\mathrm{\Lambda }`$-electron events, compared to the opposite-hemisphere, opposite-sign events to estimate the $`\mathrm{\Lambda }_c\mathrm{\Lambda }`$X branching fraction. This estimate is “internally normalizing”; i.e., we do not need to measure the fraction of our total charm tags which contain $`\mathrm{\Lambda }`$’s. We can relate the branching fractions $`(\mathrm{\Lambda }_c\mathrm{\Lambda }e\nu _e)`$ and $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)`$ (and their corresponding efficiencies $`ϵ`$) to the number of observed same-hemisphere, same-sign events ($`N_{SH/SS}`$), the number of observed opposite hemisphere, opposite sign events $`N_{OH/OS}`$, and their production fractions in cc̄ events. Without an explicit fake electron subtraction to the observed yields, we have: $$\frac{N_{SH/SS}}{N_{OH/OS}}\frac{(c\mathrm{\Lambda }_c)(\mathrm{\Lambda }_c\mathrm{\Lambda }e\nu _e)ϵ(\mathrm{\Lambda }_c\mathrm{\Lambda }e\nu _e)}{(c\mathrm{\Lambda }_c)(\overline{c}eX)(\mathrm{\Lambda }_c\mathrm{\Lambda }X)(ϵ(\mathrm{\Lambda }_c\mathrm{\Lambda }X)(\overline{c}eX))}$$ Note that the efficiency in the numerator of this equation refers to the correlated efficiency of having both the $`\mathrm{\Lambda }`$ and the electron in $`\mathrm{\Lambda }_c\mathrm{\Lambda }e^+\nu _e`$ pass all our selection criteria ($`ϵ(\mathrm{\Lambda }_c\mathrm{\Lambda }e\nu _e)=0.023\pm 0.002`$); the efficiency in the denominator refers to the efficiency for having a $`\mathrm{\Lambda }`$ from a $`\mathrm{\Lambda }_c`$ decay pass our selection criteria in one hemisphere, and also an electron from a generic charm decay pass our selection requirements in the opposite hemisphere ($`ϵ(\mathrm{\Lambda }_c\mathrm{\Lambda }X)(\overline{c}eX)=0.043\pm 0.002`$). The efficiency is lower in the numerator due to the presence of momentum correlations between the $`\mathrm{\Lambda }`$ and the electron, resulting in a reduced efficiency for both particles to simultaneously pass the minimum momentum requirement $`p>`$1 GeV/c. The value for $`(\overline{c}eX)`$ (0.091$`\pm `$0.008) is taken from data at $`\sqrt{s}`$=10 GeV. Using the current Particle Data Group value for $`(\mathrm{\Lambda }_c\mathrm{\Lambda }e\nu _e)=0.021\pm 0.006`$ and our measured values for $`N_{SH/SS}`$ ($`445\pm 26`$) and $`N_{OH/OS}`$ ($`743\pm 32`$), we obtain an inferred value of $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)0.23\pm 0.07`$, where the error is statistical only. This is consistent with our measured product branching fraction for $`_\mathrm{\Lambda }=(c\mathrm{\Theta }_c)(\mathrm{\Theta }_c\mathrm{\Lambda }X`$) if $`(c\mathrm{\Theta }_c`$)=6.67%, and assuming that $`\frac{(c\mathrm{\Lambda }_c)}{(c\mathrm{\Theta }_c)}1.0`$. ## VIII Systematic Errors In order to determine additional systematic errors, we varied each of our individual particle and event requirements and noted the corresponding variation in our derived values for $`(c\mathrm{\Theta }_c)(\mathrm{\Theta }_c\mathrm{\Lambda }X`$). Typical variations were of order $``$20%. Using this approach, we summarize our systematic dependencies in Table IV. The default values of, e.g., our kinematic cuts are defined, as well as the variation used to assess systematic dependencies. For the pion and the electron tags, our largest systematic errors are due to uncertainties in the Monte Carlo event generation modeling, as determined using the “wrong-sign” yields. For the $`\overline{D}`$-tags, among the largest errors are the errors associated with signal extraction - this is assessed by determining the difference in the calculated final result when the signal and sideband regions are varied from their default values by $`\pm `$30%. We also take the r.m.s. spread in the values for $`_\mathrm{\Lambda }`$ obtained with the four tags (14%) as an additional systematic error, reflecting the differences in the lepton-tagged vs. $`\pi _{soft}^{}`$-tagged vs. D-tagged samples. This variance is the dominant contributor to our overall quoted systematic error. ## IX Summary and Discussion Using four different $`e^+e^{}c\overline{c}`$ tags, we measure the product branching fraction $`_\mathrm{\Lambda }`$ = $`(c\mathrm{\Theta }_cX)(\mathrm{\Theta }_c\mathrm{\Lambda }X`$): $`(1.62\pm 0.10\pm 0.32)\%(electrontags)`$ $`(1.53\pm 0.06\pm 0.30)\%(softpiontags)`$ $`(2.12\pm 0.09\pm 0.30)\%(\overline{D}^0tags)`$ $`(2.09\pm 0.13\pm 0.42)\%(D^{}tags);`$ these results sum over the charmed baryons $`\mathrm{\Theta }_c`$ produced at $`\sqrt{s}`$=10 GeV. Separating common from independent systematic errors, and weighting each result by the quadrature sum of its statistical error plus independent systematic error, we combine these four numbers to obtain a weighted product branching fraction: $`_\mathrm{\Lambda }=(1.87\pm 0.03\pm 0.33)\%`$ In obtaining this result, we have not corrected for the statistical overlap between the four tag samples. Correcting for this would tend to slightly reduce the overall quoted statistical error. We can convert this result into a contour in the plane: $`(\mathrm{\Theta }_c\mathrm{\Lambda }X)`$ vs. $`(c\mathrm{\Theta }_c)`$, as shown in Figure 7. Using the Monte Carlo value for $``$($`c\mathrm{\Theta }_cX`$) of 6.67% (this is consistent with the tabulated product cross-section $`(e^+e^{}c\overline{c})(c\mathrm{\Lambda }_c)(\mathrm{\Lambda }_cpK^{}\pi ^+)`$ using $`(\mathrm{\Lambda }_cpK^{}\pi ^+)`$=5.0%), and taking the results from our four tags, we can infer a weighted average value: $`(\mathrm{\Theta }_c\mathrm{\Lambda }X)=(28\pm 1\pm 5)\%`$It is important to note that this measurement is independent of the $`\mathrm{\Lambda }_cpK^{}\pi ^+`$ normalization but is dependent on the Monte Carlo estimated value for $`(c\mathrm{\Theta }_c)`$. This measurement is the first of its kind at $`\sqrt{s}`$ = 10 GeV. In the simplest picture, a charmed baryon such as a $`\mathrm{\Lambda }_c`$ decays weakly through external W-emission. Neglecting fragmentation at the lower vertex, this produces either a $`\mathrm{\Sigma }^0`$ or, if isospin does not change, a $`\mathrm{\Lambda }`$. Since all $`\mathrm{\Sigma }^0`$’s decay into $`\mathrm{\Lambda }`$, we therefore expect that $`\mathrm{\Lambda }_c\mathrm{\Lambda }X100\%`$ if the external spectator diagram dominates. This simple-minded prediction is expected to be obeyed in semileptonic decays; i.e., $`(\mathrm{\Lambda }_c\mathrm{\Lambda }l\nu )/(\mathrm{\Lambda }_cXl\nu )1`$. Present data, however, give a value of approximately 50% for this ratio, albeit with large errors. External and internal W-emission, as well as W-exchange can lead to $`NKX`$ final states ($`\mathrm{\Lambda }_cpK_\mathrm{s}^0`$, e.g.). The fact that the $`\mathrm{\Lambda }_c`$ lifetime is only half that of the $`D^0`$ meson suggests that internal W-emission and W-exchange processes may comprise a large fraction of the total $`\mathrm{\Lambda }_c`$ width. Although internal W-emission may be suppressed in decays of charmed mesons due to the color-matching requirement (which would predict $`(DW_{int}X)/(DW_{ext}X)=1/9`$), the larger number of degrees of freedom in baryon decays may mitigate this suppression, leading to a potentially large fraction of $`pKX`$ final states. In the case of the $`\mathrm{\Lambda }_c`$, W-exchange decays can produce either $`\mathrm{\Lambda }`$’s or $`NK`$ in the final state, depending on the quark configuration. The naive expectation that the absence of exchange diagrams in $`\mathrm{\Xi }_c^+`$ decays will lead to a longer lifetime for $`\mathrm{\Xi }_c^+`$ compared to $`\mathrm{\Xi }_c^0`$ and $`\mathrm{\Lambda }_c`$ is consistent with current experimental data. The current world average for $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)`$ (35 $`\pm `$ 11)% is consistent with the notion that the simple-minded external W-emission picture does not saturate $`\mathrm{\Lambda }_c`$ decays. The value for $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)`$ therefore has implications for the external versus internal spectator fractions in charmed baryon decay. Our results are therefore qualitatively consistent with a possibly substantial internal spectator contribution to charmed baryon decay. Exclusive $`\mathrm{\Lambda }_c\mathrm{\Lambda }X`$ channels have also been measured; normalized to an estimate of $`(\mathrm{\Lambda }_cpK^{}\pi ^+)=(5.0\pm 1.3)\%`$ , the sum of the observed exclusive modes account for the bulk of the presently tabulated inclusive $`\mathrm{\Lambda }_c\mathrm{\Lambda }X`$ rate ($`\frac{\mathrm{\Sigma }(\mathrm{\Lambda }_c\mathrm{\Lambda }+X)_{exclusive}}{(\mathrm{\Lambda }_cpK^{}\pi ^+)}`$5, where the sum includes a contribution of $`\frac{\mathrm{\Sigma }(\mathrm{\Lambda }_c\mathrm{\Sigma }^0+X)_{exclusive}}{(\mathrm{\Lambda }_cpK^{}\pi ^+)}`$1.5). We note that the difference between our inferred value for $`(\mathrm{\Theta }_c\mathrm{\Lambda }X)`$ and the sum of the exclusive $`\mathrm{\Lambda }_c`$ modes to $`\mathrm{\Lambda }`$’s suggests that most of the inclusive $`\mathrm{\Lambda }_c\mathrm{\Lambda }`$X rate has been accounted for. If charmed baryons produced in $`e^+e^{}`$ events are predominantly $`\mathrm{\Lambda }_c`$’s, and if the JETSET expectation for $`f(c\mathrm{\Lambda }_c)`$ is accurate, then our results are in agreement with the current world average for $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)`$. We note that the methodology of this analysis differs substantially from the previous CLEO analysis, which relied on a model of charmed baryon production in B-decay to derive $`(\mathrm{\Lambda }_c\mathrm{\Lambda }X)`$. Naively, one might expect fragmentation and decay of charmed baryons to be similar to bottom baryons. Using vertex tagging techniques, the OPAL collaboration has determined $`(b\mathrm{\Theta }_b)(\mathrm{\Theta }_b\mathrm{\Lambda }X)=(3.50\pm 0.32\pm 0.35)\%`$ . Perhaps the simplest way to reconcile the two numbers is to assume (ad hoc) that $`(b\mathrm{\Theta }_b)`$ at $`\sqrt{s}`$90 GeV is approximately twice as large as $`(c\mathrm{\Theta }_c)`$ at $`\sqrt{s}`$=10 GeV, and that $`(\mathrm{\Lambda }_b\mathrm{\Lambda }_cX)`$1.0. However, the fact that the $`\mathrm{\Lambda }_b`$ lifetime is only 2/3 that of the B-mesons , coupled with the fact that $`\mathrm{\Lambda }_b`$ has already been observed through $`\mathrm{\Lambda }_b\psi \mathrm{\Lambda }`$ imply that $`(\mathrm{\Lambda }_b\mathrm{\Lambda }_cX)<1`$. Correspondingly, we expect an enhancement of $`(b\mathrm{\Theta }_b)`$ at OPAL relative to $`(c\mathrm{\Theta }_c)`$ at CLEO. Finally, we stress that our final central value for $`_\mathrm{\Lambda }`$ averages over the specific mix of charm-tags that we use in this analysis. The composition of our $`\overline{D}`$-meson tags will not be the same as the composition of our electron tags, insofar as the lepton-tagged sample represents a weighted sum of $`\overline{\mathrm{\Theta }}_cl^{}X`$, $`D_s^{}l^{}X`$, $`\overline{D}^0l^{}X`$ and $`D^{}l^{}X`$. Our quoted final result can be general only if the two hemispheres in an $`e^+e^{}c\overline{c}`$ event fragment independently. Although not yet measured, it is possible that there may be correlated $`\overline{\mathrm{\Theta }}_c\mathrm{\Theta }_c`$ production, in which case the likelihood of observing $`\mathrm{\Theta }_c\mathrm{\Lambda }X`$ would be larger for $`\overline{\mathrm{\Theta }}_c`$ tags than for $`\overline{D}`$ tags, and the assumption of independent fragmentation would be invalid. A study of correlated $`\overline{\mathrm{\Theta }}_c\mathrm{\Theta }_c`$ production, presently in progress, will be the subject of a forthcoming publication. ###### Acknowledgements. We gratefully acknowledge the effort of the CESR staff in providing us with excellent luminosity and running conditions. J.R. Patterson and I.P.J. Shipsey thank the NYI program of the NSF, M. Selen thanks the PFF program of the NSF, M. Selen and H. Yamamoto thank the OJI program of DOE, J.R. Patterson, K. Honscheid, M. Selen and V. Sharma thank the A.P. Sloan Foundation, M. Selen and V. Sharma thank the Research Corporation, F. Blanc thanks the Swiss National Science Foundation, and H. Schwarthoff and E. von Toerne thank the Alexander von Humboldt Stiftung for support. This work was supported by the National Science Foundation, the U.S. Department of Energy, and the Natural Sciences and Engineering Research Council of Canada.
warning/0004/hep-th0004047.html	ar5iv	text	# Note on irreducible approach to reducible second-class constraints ## Acknowledgment Two of the authors (C.B. and S.O.S.) acknowledge financial support from a Romanian National Council for Academic Scientific Research (CNCSIS) grant.
warning/0004/hep-ph0004090.html	ar5iv	text	# 1 Introduction ## 1 Introduction The Standard Model (SM), while being extremely successful in describing interactions of quarks and leptons at low energies, still has many unanswered questions. Among these one problem is predominant: the unification of all elementary forces within the framework of a simple gauge theory and the ensuing hierarchy of mass scales in particle physics at large and small distances. Possible solutions to this problem are commonly related to supersymmetry (SUSY) and Grand Unified Theories (GUT) . At present they receive some indirect (largely qualitative) experimental support from the apparent lightness of the Higgs boson, the values of gauge couplings given by precision measurements and the heavy top quark mass. At low energies the SUSY GUT turns into the Minimal Supersymmetric Standard Model (MSSM). Of course the MSSM can certainly be considered on its own as a simple SUSY extension of the Standard Model, leaving aside for the moment the question of unification. However, no matter which level of theory is considered, there is one point that crucially distinguishes SUSY from non-SUSY models. This is that they do not contain the automatic accidental symmetries, corresponding to baryon (B) and lepton (L) number conservation, which are present in the ordinary Standard Model. Does this mean that B and L number can be violated in the SUSY context or must some special protecting symmetry be postulated in the MSSM? Usually, the requirement of B and L number conservation in the MSSM is indeed satisfied by postulating the existence of some multiplicative discrete symmetry called R-parity . An exact R-parity ($`RP`$) implies that SUSY particles should be produced in pairs and that the lightest SUSY particle (LSP) is stable. On the other hand, there is no fundamental reason to prefer models with exact $`RP`$ over those with broken $`RP`$ in the framework of the supersymmetric SM, where not only fermions but also their scalar superpartners automatically become the carriers of lepton and baryon numbers. Thereby, among the basic renormalisable couplings in the low–energy MSSM superpotential,one would generally expect to find the lepton and baryon number violating ones $$\mathrm{\Delta }W=\mu _iL_iH_u+\lambda _{ijk}L_iL_j\overline{E}_k+\lambda _{ijk}^{}L_iQ_j\overline{D}_k+\lambda _{ijk}^{\prime \prime }\overline{U}_i\overline{D}_j\overline{D}_k.$$ (1) Here, $`i`$, $`j`$, $`k`$ are generation indices and a summation is implied (colour and weak isospin indices are suppressed); $`L_i`$($`Q_j`$) denote the lepton (quark) $`SU(2)`$–doublet superfields and $`\overline{E}_i`$($`\overline{U}_i`$, $`\overline{D}_i`$) are $`SU(2)`$–singlet lepton ($`up`$–quark, $`down`$–quark) superfields; $`\mu _i`$ are mass parameters which mix lepton superfields with the $`up`$–type Higgs superfield $`H_u`$, while $`\lambda _{ijk}`$ ($`\lambda _{ijk}=\lambda _{jik}`$), $`\lambda _{ijk}^{}`$ and $`\lambda _{ijk}^{\prime \prime }`$ ($`\lambda _{ijk}^{\prime \prime }=\lambda _{ikj}^{\prime \prime }`$) are dimensionless couplings. The first three terms in (1) violate lepton number, while the last violates baryon number. While SUSY-inspired B number violation (BNV) leads in general to unacceptably fast proton decay and must be highly suppressed, SUSY-inspired L number violation (LNV) could readily occur at a level consistent with present experimental constraints, but large enough for the observation of some of its spectacular manifestations at present or future colliders. Remarkably enough, instead of $`RP`$, another (gauge) discrete symmetry could appear in the MSSM: superstring-inherited $`Z_3`$ baryon parity , which strongly protects B number and allows for L number violation only. Thus, as long as it is not in conflict with any phenomenology, SUSY-inspired lepton number violation merits further detailed investigation, both theoretically and experimentally. From the theoretical point of view, the principal question concerns the search for a Grand Unified framework within which, while treating quarks and leptons equally, L-violation should be allowed at the same time as B-conservation. Unfortunately, the discrete symmetries acceptably protecting B-conservation while allowing L-violation in the MSSM, such as the above mentioned $`Z_3`$ baryon parity, transform quarks and leptons differently and, thereby, are incompatible with the known GUTs. Nevertheless several extended GUT models have been constructed , where the coexistence of lepton number violation and baryon number conservation can in principle be arranged. This is typically achieved by introducing at the Planck scale high-dimensional operators, involving Higgs and matter multiplets in some exotic representations of the underlying GUT symmetry, and then imposing additional custodial symmetries to ensure that only the required set of LNV high-order operators is allowed. These operators become the renormalisable LNV couplings (1) at lower energies after the GUT symmetry breaks at the GUT scale $`M_{GUT}`$. Despite some progress, one has the uneasy feeling that such a solution to this problem looks rather artificial, as it is generically correlated neither with the nature of the GUT nor with its breaking pattern. Instead, we suggest that it is just the breaking pattern of the underlying GUT symmetry which could give a fundamental reason for the difference in treatment of the baryon and lepton numbers of the matter particles involved in GUTs. We show that a suppression of baryon number violating interactions in the superpotential (1) naturally occurs in some $`SU(N)`$ SUSY GUTs where a missing VEV vacuum configuration develops, which also gives a solution to the doublet-triplet splitting problem . We construct explicit examples of $`RP`$-violating $`SU(7)`$ and $`SU(8)`$ GUTs where the effective LNV couplings immediately evolve from the GUT scale, while the baryon number non-conserving ones are safely projected out by the missing VEV vacuum configuration breaking the GUT symmetry down to that of the MSSM. However, at the next stage when SUSY breaks, radiative corrections shift the missing VEV to some nonzero value of order $`M_{SUSY}`$ and induce BNV violating couplings with hierarchically small coupling constants $`\lambda _{ijk}^{\prime \prime }=O(M_{SUSY}/M_{GUT})`$, which appear to be of phenomenological interest . ## 2 Missing VEV solutions in $`SU(N)`$ GUTs The most elegant solution to the gauge hierarchy problem in supersymmetric $`SU(N)`$ GUTs could well be related to the existence of a missing VEV vacuum configuration , according to which the basic adjoint scalar $`\mathrm{\Sigma }_j^i`$ ($`i,j=1,\mathrm{},N`$) does not develop a VEV in some of the directions in $`SU(N)`$ space. Through its coupling with a pair of Higgs fields $`H`$ and $`\overline{H}`$, their masses are split in a hierarchical way so as to have light weak doublets breaking electroweak symmetry and giving masses to up and down quarks, on the one hand, and superheavy colour triplets mediating proton decay, on the other. However, it is well known that a missing VEV solution can not appear in $`SU(N)`$ GUTs in the ordinary one-adjoint scalar case. This is due to the presence of a cubic term $`\mathrm{\Sigma }^3`$ in the general Higgs superpotential $`W`$ leading to the unrealistic trace condition $`Tr\mathrm{\Sigma }^2=0`$ for the missing VEV vacuum configuration, unless there is a special fine-tuned cancellation between $`Tr\mathrm{\Sigma }^2`$ and driving terms stemming from other parts of the superpotential $`W`$. ### 2.1 The two-adjoint alternative So, it seems the only way to obtain a natural missing VEV solution in $`SU(N)`$ theories is to exclude the cubic term $`\mathrm{\Sigma }^3`$ from the superpotential, by imposing some extra reflection symmetry on the adjoint supermultiplet $`\mathrm{\Sigma }`$ $$\mathrm{\Sigma }\mathrm{\Sigma }$$ (2) On its own the elimination of the $`\mathrm{\Sigma }^3`$ term leads to the trivial unbroken symmetry case. However the inclusion of higher even-order $`\mathrm{\Sigma }`$ terms (supposedly inherited from superstrings or induced by gravitational corrections) in the effective superpotential leads to an all-order missing VEV solution, as was shown in recent papers . Alternatively one can introduce another adjoint scalar $`\mathrm{\Omega }`$ with only renormalisable couplings appearing in $`W`$. Let us consider briefly the high-order term case first. The $`SU(N)`$ invariant superpotential for an adjoint scalar field $`\mathrm{\Sigma }`$ conditioned also by the gauge $`Z_2`$ reflection symmetry (2) $$W_A=\frac{1}{2}m\mathrm{\Sigma }^2+\frac{\lambda _1}{4M_P}\mathrm{\Sigma }^4+\frac{\lambda _2}{4M_P}\mathrm{\Sigma }^2\mathrm{\Sigma }^2+\mathrm{}$$ (3) contains, in general, all possible even-order $`\mathrm{\Sigma }`$ terms scaled by inverse powers of the (conventionally reduced) Planck mass $`M_P=(8\pi G_N)^{1/2}2.410^{18}`$ GeV. It is readily shown that the necessary condition for any missing VEV solution to appear in the $`SU(N)`$ $`Z_2`$ invariant superpotential $`W_A`$ is the tracelessness of all the odd-order $`\mathrm{\Sigma }`$ terms $$Tr\mathrm{\Sigma }^{2s+1}=0\text{ , }s=0,1,2,\mathrm{}$$ (4) This condition uniquely leads to a missing VEV pattern of the type $`Nk\text{ }k/2\text{ }k/2`$ (5) $`<`$ $`\mathrm{\Sigma }>=\sigma Diag(\stackrel{}{0\mathrm{}0},\stackrel{}{1\mathrm{}1},\stackrel{}{1\mathrm{}1})\text{,}`$ where the VEV value $`\sigma `$ is calculated using the $`\mathrm{\Sigma }`$ polynomial taken in $`W_A`$ (3). The vacuum configuration (5) gives rise to a particular breaking channel of the $`SU(N)`$ GUT symmetry $$SU(N)SU(Nk)SU(k/2)SU(k/2)U(I)_1U(I)_2\text{ ,}$$ (6) which we will discuss in some detail a little later. So we conclude from Eqs. (5, 6) that a missing VEV solution could actually exist, with the ordinary MSSM gauge symmetry $`SU(3)_CSU(2)_WU(1)_Y`$ surviving at low energies, provided that $`N7`$. The superpotential (3) can be viewed as an effective one, following from an ordinary renormalisable two-adjoint superpotential with the second heavy adjoint scalar integrated out. Hereafter, although both approaches are closely related, we deal for simplicity with the two-adjoint case. So let us consider a general $`SU(N)`$ invariant renormalisable superpotential for two adjoint scalars $`\mathrm{\Sigma }`$ and $`\mathrm{\Omega }`$, also satisfying the gauge-type $`Z_2`$ reflection symmetry ($`\mathrm{\Sigma }\mathrm{\Sigma }`$, $`\mathrm{\Omega }\mathrm{\Omega }`$) inherited from superstrings: $$W_A=\frac{1}{2}m\mathrm{\Sigma }^2+\frac{1}{2}M_P\mathrm{\Omega }^2+\frac{1}{2}h\mathrm{\Sigma }^2\mathrm{\Omega }+\frac{1}{3}\lambda \mathrm{\Omega }^3.$$ (7) Here the second adjoint $`\mathrm{\Omega }`$ can be considered as a state originating from a massive string mode with the Planck mass $`M_P`$. The basic adjoint $`\mathrm{\Sigma }`$ may be taken at another well motivated scale $`mM_P^{2/3}M_{SUSY}^{1/3}O(10^{13})`$ GeV where, according to many string models, the adjoint moduli states $`(1_c,1_w)`$, $`(1_c,3_w)`$ and $`(8_c,1_w)`$ (in a self-evident $`SU(3)_CSU(2)_W`$ notation) appear. In the present context these states can be identified as just the non-Goldstone remnants $`\mathrm{\Sigma }_{0,}`$ $`\mathrm{\Sigma }_3`$ and $`\mathrm{\Sigma }_8`$ of the relatively light adjoint $`\mathrm{\Sigma }`$ which breaks $`SU(N)`$ in some way. However, all our conclusions remain valid for any reasonable value of $`m`$, which is the only mass parameter (apart from $`M_P`$) in the model considered. As a general analysis of the superpotential $`W_A`$ (7) shows , there are just four possible VEV patterns for the adjoint scalars $`\mathrm{\Sigma }`$ and $`\mathrm{\Omega }`$: (i) the trivial unbroken symmetry case, $`\mathrm{\Sigma }=\mathrm{\Omega }=0`$; (ii) the single-adjoint condensation, $`\mathrm{\Sigma }=0`$, $`\mathrm{\Omega }0`$; (iii) the <sup>′′</sup>parallel<sup>′′</sup> vacuum configurations, $`\mathrm{\Sigma }\mathrm{\Omega }`$ and (iv) the <sup>′′</sup>orthogonal<sup>′′</sup> vacuum configurations, $`Tr(\mathrm{\Sigma }\mathrm{\Omega })=0`$. The Planck-mass mode $`\mathrm{\Omega }`$, having a cubic term in $`W_A`$, in all non-trivial cases develops a standard <sup>′′</sup>single-breaking<sup>′′</sup> VEV pattern $`Nkk`$ (8) $`<`$ $`\mathrm{\Omega }>\text{ }=\omega Diag(\stackrel{}{1\mathrm{}1},\stackrel{}{{\displaystyle \frac{Nk}{k}}\mathrm{}{\displaystyle \frac{Nk}{k}}})\text{,}`$ which breaks the $`SU(N)`$ GUT symmetry to $$SU(N)SU(Nk)SU(k)U(I)\text{ .}$$ (9) However, in case (iv), the basic adjoint $`\mathrm{\Sigma }`$ develops the radically new missing VEV vacuum configuration (5), thus giving a <sup>′′</sup>double breaking<sup>′′</sup> of $`SU(N)`$ to (6). Using the approximation $`\frac{h}{\lambda }>>\frac{m}{M_P}`$, which is satisfied for any reasonable values of the couplings $`h`$ and $`\lambda `$ in the generic superpotential $`W_A`$ (7), the VEV values are given by $$\omega =\frac{k}{Nk}\frac{m}{h}\text{}\sigma =\left(\frac{2N}{Nk}\right)^{1/2}\sqrt{mM_P}/h$$ (10) respectively. Surprisingly, just the light adjoint $`\mathrm{\Sigma }`$ develops the largest VEV in the model which, for a properly chosen adjoint mass $`m`$ and coupling constant $`h`$, can easily come up to the string scale $`M_{str}`$ (see ). Furthermore, as concluded above, one must consider $`SU(N)`$ GUTs with $`N7`$, in order to have the standard gauge symmetry $`SU(3)_CSU(2)_WU(1)_Y`$ remaining after the breaking (6). As is easily seen from Eqs. (5, 6), there are two principal possibilities: the weak-component and colour-component missing VEV solutions respectively. If it is granted that the ”missing VEV subgroup” $`SU(Nk)`$ in (6) is just the weak symmetry group $`SU(2)_W`$, as is traditionally argued , one is led to $`SU(8)`$ as the minimal GUT symmetry ($`Nk=2,k/2=3`$) . Another, and in fact the minimal, possibility is to identify $`SU(Nk)`$ with the colour symmetry group $`SU(3)_C`$ in the framework of an $`SU(7)`$ GUT symmetry ($`Nk=3,`$ $`k/2=2`$) . The higher $`SU(N)`$ GUT solutions, if considered, are also based on just those two principal possibilities: the weak-component or colour-component missing VEV vacuum configurations respectively. Let us see now how this missing VEV mechanism works to solve the doublet-triplet splitting problem in $`SU(8)`$ or $`SU(7)`$ GUT with the superpotential $`W_A`$ (7). It is supposed that there is a reflection-invariant coupling of the ordinary MSSM Higgs-boson containing supermultiplets $`H`$ and $`\overline{H}`$ with the basic adjoint $`\mathrm{\Sigma }`$, but not with $`\mathrm{\Omega }`$, in the superpotential $`W_H`$ $$W_H=f_0\overline{H}\mathrm{\Sigma }H+W_H^{}(\mathrm{\Sigma }\mathrm{\Sigma },\overline{H}H\overline{H}H)$$ (11) The second part $`W_H^{}`$contains possible mixings with other scalar fields, which are inessential for the moment. The superfields $`H`$ and $`\overline{H}`$ do not develop VEVs during the first stage of the symmetry breaking (6). Thereupon the first term in $`W_H`$ turns into a mass term for $`H`$ and $`\overline{H}`$ determined by the missing VEV pattern (5). This vacuum, while giving generally heavy masses (of the order of $`M_{GUT}`$) to $`H`$ and $`\overline{H}`$, leaves their weak components strictly massless. To be certain of this, we must specify the multiplet structure of $`H`$ and $`\overline{H}`$ for both the weak-component and the colour-component missing VEV vacuum configurations, that is in $`SU(8)`$ and $`SU(7)`$ GUTs respectively. In the $`SU(8)`$ case $`H`$ and $`\overline{H}`$ are fundamental octets whose weak components (ordinary Higgs doublets) do not get masses from the basic coupling (11). In the $`SU(7)`$ case $`H`$ and $`\overline{H}`$ are 2-index antisymmetric $`21`$-plets in which, after projecting out the extra heavy states (see Section 3.1), there is left just one pair of massless Higgs doublets as a consequence of the coupling (11). Thus, there is a natural doublet-triplet splitting in both cases and we also have a vanishing $`\mu `$ term at this stage. However, radiative corrections generate a $`\mu `$ term of the right order of magnitude at the next stage when SUSY breaks . ### 2.2 Projection to low energies Missing VEV vacua, which ensure the survival of the MSSM at low energies, only appear in $`SU(N)`$ GUTs with a higher symmetry group than the standard $`SU(5))`$ model. In order not to spoil gauge coupling unification, the extra gauge symmetry should also be broken, $`SU(N)SU(5)`$, at the GUT scale. Then the following question arises: how can the missing VEV survive this extra symmetry breaking with at most a shift of order the electroweak scale? This requires, in general, that the superpotential (7) be strictly protected from any large influence from the $`N5`$ scalars $`\phi ^k`$ ($`k=1,\mathrm{},N5`$) providing the extra symmetry breaking (or from uncontrollable gravitational corrections). Technically, such a custodial symmetry may be a superstring-inherited anomalous $`U(1)_A`$ , which can naturally keep two sectors of the total superpotential separate and then induce a high-scale extra symmetry breaking through the Fayet-Iliopoulos (FI) $`D`$-term : $$D_A=\xi +Q_A^k<\phi ^k>^2,\xi =\frac{TrQ_A}{192\pi ^2}g_{str}^2M_P^2.$$ (12) Here the sum runs over all ”charged” scalar fields in the theory, including those which do not develop VEVs and which contribute to $`TrQ_A`$ only. For realistic or semi-realistic models, $`TrQ_A`$ has turned out to be quite large, $`TrQ_A=O(100)`$ (see for a recent discussion). Therefore, the spontaneous breaking scale of the $`U(1)_A`$ symmetry and of the related extra gauge symmetry is naturally located at the string scale. The protecting anomalous $`U(1)_A`$ symmetry is needed to keep the scalars $`\phi ^{(k)}`$ and $`\overline{\phi }^{(k)}`$ essentially decoupled from the basic adjoint superpotential $`W_A`$ (7), so as not to strongly influence its missing VEV vacuum configuration (5). Otherwise potentially dangerous couplings could appear of the type $`\overline{\phi }^{(k)}\mathrm{\Sigma }\phi ^{(k)}`$, where the $`\phi ^{(k)}`$ and $`\overline{\phi }^{(k)}`$ scalar superfields are taken in pairs of conjugate fundamental representations ($`N`$ and $`\overline{N}`$) of $`SU(N)`$. If these couplings actually appeared, they would give rise to shifts in the missing VEV components of the adjoint scalar $`\mathrm{\Sigma }_B^A`$ of the order $`\mathrm{\Sigma }_B^A\frac{\delta _B^A}{m}\overline{\phi }^{(k)}\phi ^{(k)}O(M_{GUT})`$, as directly follows from the minimisation condition for the scalar potential. So the presence of a protecting symmetry is essential for the missing VEV mechanism to function properly. We will now enlarge on this key point in order to gain a better understanding of the missing VEV approach. The symmetry protected separation of the adjoint scalar and the $`\phi ^k`$ scalar sectors in the total superpotential implies the appearance of an accidental global symmetry $`SU(N)_{\mathrm{\Sigma }\mathrm{\Omega }}U(N)\phi `$ in the $`SU(N)`$ $`U(1)_A`$ gauge theory considered. This global symmetry is in turn radiatively broken, resulting in a set of pseudo-Goldstone (PG) states of the type $$5+\overline{5}+SU(5)\mathrm{singlets}$$ (13) which gain a mass at the TeV scale where SUSY softly breaks . There can be a maximum of $`N5`$ families of PG states of the type (13), corresponding to the case where the scalars $`\phi ^{(k)}`$ and $`\overline{\phi }^{(k)}`$ are only allowed to appear in the Higgs potential through the basic $`SU(N)`$ and $`U(1)_A`$ $`D`$-terms. In this case the $`U(N)\phi `$ global symmetry is increased to $`U(N)_{\phi ^{(1)}}\mathrm{}.U(N)_{\phi ^{(N5)}}`$. This case would occur if the $`U(1)_A`$ charges of the bilinears $`\overline{\phi }^k\phi ^k^{}`$ were all positive (or negative), so that they could not appear in the $`SU(N)U(1)_A`$ invariant superpotential in any order. However, in a properly extended model it is possible for the adjoint and fundamental scalar sectors in the superpotential to overlap without disturbing the adjoint missing VEV configuration. This naturally occurs when the scalars $`\phi ^{(k)}`$ are conditioned by the $`U(1)_A`$ symmetry to develop orthogonal VEVs along the <sup>′′</sup>extra<sup>′′</sup> directions $$\phi _A^{(k)}=\delta _{A,5+k}V_k,k=1,\mathrm{},N5$$ (14) As a result, some safe non-diagonal couplings $`\overline{\phi }^{(m)}\mathrm{\Sigma }\phi ^{(n)}`$ are generated between the two sectors, giving contributions to the pseudo-Goldstone masses which leave only one light PG family (12). Let us consider this possibility in some detail. The least restrictive choice of such safe mixing terms for the general $`SU(N)`$ case is achieved by introducing two sets of new singlet scalar superfields fields, $`S_{mn}`$ and $`T_{mn}`$, with non-diagonal couplings of the type $$W_{mix}=\underset{m<n}{\overset{N5}{}}\overline{\phi }^{(m)}[a_{mn}S_{mn}+b_{mn}T_{mn}\mathrm{\Sigma }]\phi ^{(n)}$$ (15) which are also invariant under the reflection symmetry $`\mathrm{\Sigma }\mathrm{\Sigma }`$, $`T_{mn}T_{mn}`$. The coupling constants $`a_{mn}`$ and $`b_{mn}`$ are all of order $`O(1)`$ and $`O(1/M_P)`$ respectively, and the $`(N5)(N6)/2`$ singlet scalars $`S_{mn}`$ and $`T_{mn}`$ ($`m<n`$) get their VEVs through the FI $`D`$-term (12), as do all the $`\phi `$ and $`\overline{\phi }`$ scalars. One can consider the fields $`S_{mn}`$ as the basic carriers of the $`U(1)_A`$ charges $`Q_{mn}`$ which are all taken positive in the model (the fields $`T_{mn}`$ carry the same charges $`Q_{mn}`$). The $`U(1)_A`$ charges of the $`\overline{\phi }^{(m)}\phi ^{(n)}`$ bilinears ($`m<n`$) appearing in $`W_{mix}`$ are then determined to be $`Q_{mn}`$, while the charges of all the other bilinears, diagonal $`\overline{\phi }^{(m)}\phi ^{(m)}`$ and non-diagonal $`\overline{\phi }^{(n)}\phi ^{(m)}`$, can always be chosen positive. This implies that any terms containing $`\phi `$ and $`\overline{\phi }`$ scalars can only appear in the superpotential if they also include the bilinears $`\overline{\phi }^{(m)}\phi ^{(n)}`$ so as to properly compensate the $`U(1)_A`$ charges. However, for a vacuum configuration where the orthogonality conditions $`\overline{\phi }^{(m)}\phi ^{(n)}=0`$ naturally arise, such terms do not lead (in any order) to the dangerous $`\overline{\phi }^{(k)}\mathrm{\Sigma }\phi ^{(k)}`$ couplings, although they can can contribute to the pseudo-Goldstone masses. In fact these orthogonality conditions are satisfied at the SUSY invariant global minimum of the Higgs potential, as follows from the vanishing $`F`$-terms of the superfields $`S_{mn}`$ ($`T_{mn}`$), $`\overline{\phi }^{(m)}`$ and $`\phi ^{(n)}`$ involved in (15): $$\overline{\phi }^{(m)}\phi ^{(n)}=0,a_{mn}S_{mn}=b_{mn}T_{mn}\mathrm{\Sigma }_{5+n}^{5+n},m<n$$ (16) (no summation is implied). Here the orthogonal VEV values of the scalars $`\phi ^{(n)}`$ (14) have been used. One can now readily see that non-diagonal mass terms appear for the PG states related to the multiplets $`\overline{\phi }^{(m)}`$ and $`\phi ^{(n)}`$ $$M_{mn}[W_{mix}^{^{\prime \prime }}]_{\overline{\phi }^{(m)}\phi ^{(n)}}=b_{mn}T_{mn}(\mathrm{\Sigma }\mathrm{\Sigma }_{5+n}^{5+n}I),m<n$$ (17) where $`I`$ is the $`N\times N`$ unit matrix. Diagonalisation of the mass matrix (17) explicitly shows that one PG superposition $`5+\overline{5}`$ (13) is left massless, while the others become heavy<sup>1</sup><sup>1</sup>1This mass matrix is in fact a ”triangular” $`(N5)\times (N5)`$ matrix with zeros on the main diagonal, $`M_{mn}=0`$ for $`mn`$. Such a matrix has in general only one zero eigenvalue.. This is in fact a general consequence of the symmetry breaking pattern involved. The point is that neither of the other mass-terms $`M_{mm}`$ and $`M_{nm}`$ can be allowed by $`U(1)_A`$ symmetry for any generalisation of the superpotential $`W_{mix}`$ (15). Otherwise the dangerous $`\overline{\phi }^{(k)}\mathrm{\Sigma }\phi ^{(k)}`$ couplings inevitably appear as well. So, one can conclude that even in the general case one PG family of the type (13) always exists. Together with the ordinary quarks and leptons and their superpartners these PG states, both bosons and fermions, determine the particle spectrum at low energies. In most of what follows the existence of just one family of PG states at the sub-TeV scale will be assumed. We consider below both of the minimal possible GUTs, $`SU(7)`$ and $`SU(8)`$, with the missing VEV solution naturally allowing the survival of the MSSM down to low energies. Whereas the $`SU(7)`$ model is taken as an ordinary one-family unifying GUT , the $`SU(8)`$ model can include unification of the quark-lepton families as well . ## 3 One-family unifying GUT: $`SU(7)`$ By analogy with the standard $`SU(5)`$ model, we take the simplest anomaly–free set of matter fields, consisting of the combination of the fundamental and 2-index antisymmetrical representations of the $`SU(7)`$ gauge group $$\left[2\overline{\mathrm{{\rm Y}}}^A+\overline{\mathrm{\Psi }}^A+\mathrm{\Psi }_{[AB]}\right]_i$$ (18) ($`A,B=1,\mathrm{},7`$ are the $`SU(7)`$ indices) for each of the three quark–lepton families or generations ($`i=1,2,3)`$. The quarks and leptons belong to the multiplets $`\overline{\mathrm{\Psi }}^A(\overline{7})+`$ $`\mathrm{\Psi }_{[AB]}(21)`$ , while the extra multiplets $`\overline{\mathrm{{\rm Y}}}^A`$ are specially introduced in (18) for anomaly cancellation. There is also a set of Higgs superfields among which are the two already mentioned adjoint Higgs multiplets $`\mathrm{\Sigma }_B^A`$ and $`\mathrm{\Omega }_B^A`$, responsible for the breaking (6, 9) of $`SU(7)`$, and a conjugate pair of multiplets $`H_{[AB]}`$ and $`\overline{H}^{[AB]}`$ (being the $`21`$-plets of the $`SU(7)`$) where the ordinary electroweak doublets $`H_u`$ and $`\overline{H}_d`$ reside. Besides, as in the general $`SU(N)`$ case (see Section 2.2), there should be extra-symmetry breaking scalar superfields $`\phi ^{(p)}`$ and $`\overline{\phi }^{(p)}`$ ($`p=1,2`$) which are taken to be fundamental septets and anti-septets respectively. They are supposed to develop their string-scale order VEVs along the ”extra” directions $$\phi _A^{(1)}=\delta _{A6}V_{1,}\phi _A^{(2)}=\delta _{A7}V_2$$ (19) only through the (FI) $`D`$-term related to the $`U(1)_A`$ symmetry (12). The protecting anomalous $`U(1)_A`$ symmetry keeps the $`\phi `$ scalars decoupled from the basic adjoint superpotential $`W_A`$ (7), so as not to strongly influence the missing VEV solution (5) through dangerous couplings of the type $`\overline{\phi }^{(p)}\mathrm{\Sigma }\phi ^{(p)}`$. With the given assignment of matter and Higgs superfields, the particle spectrum at low energies looks as if one had just the standard SUSY $`SU(5)`$ as a starting GUT symmetry, except that one family of PG states of type (13) appears, when a missing VEV vacuum configuration develops in the $`SU(7)`$ GUT. With this exception, all the other $`SU(7)`$ inherited states in matter and Higgs multiplets aquire GUT scale masses due to symmetry breaking, thus completely decoupling from low-energy physics. We demonstrate this for the Higgs sector in the next sub-section. ### 3.1 Higgs sector We now show that all the states, except for one pair of weak doublets in the basic Higgs multiplets $`H_{[AB]}`$ and $`\overline{H}^{[AB]}`$, become superheavy. Firstly one substitutes the colour-component missing VEV solution, obtained from the general case (5) by setting $`N=7`$ and $`k=4`$, into the superpotential (11). Superheavy masses are thereby generated for most of the components of the $`H`$ and $`\overline{H}`$ multiplets. However, the following states (weak, colour and extra symmetry components are explicitly indicated) $$H_{w6},\overline{H}^{w6},H_{w7},\overline{H}^{w7},H_{[cc^{}]},\overline{H}^{[cc^{}]}$$ (20) still remain massless at this stage of $`SU(7)`$symmetry breaking (6). Therefore one of the two pairs of weak doublets in (20), as well as the colour triplets, must further become heavy in order to get the ordinary picture of MSSM at low energies. This happens as a result of mixing $`H`$ and $`\overline{H}`$ with the specially introduced heavy scalar supermultiplets $`\mathrm{\Phi }_{[ABC]}`$ and $`\overline{\mathrm{\Phi }}^{[ABC]}`$ (being $`35`$-plets of $`SU(7)`$) in the basic Higgs superpotential $$W_H^{}=fH\overline{\mathrm{\Phi }}\phi ^{(1)}+\overline{f}\overline{H}\mathrm{\Phi }\overline{\phi }^{(1)}+yS\overline{\mathrm{\Phi }}\mathrm{\Phi },$$ (21) ($`f`$, $`\overline{f}`$ and $`y`$ are dimensionless coupling constants) when the scalars $`\phi `$ get their VEVs, thus breaking the extra gauge symmetry. The presence of the <sup>′′</sup>conjugated<sup>′′</sup> $`\overline{\mathrm{\Phi }}H`$ and $`\mathrm{\Phi }\overline{H}`$ mixings in $`W_H^{}`$ could allow the dangerous $`\overline{\phi }\mathrm{\Sigma }\phi `$ terms, destroying the missing VEV solution, unless the bilinear term $`\overline{\mathrm{\Phi }}\mathrm{\Phi }`$ has nonzero $`U(1)_A`$ charge. Therefore, this term appears in $`W_H^{}`$ together with the singlet scalar superfield $`S`$, the basic $`U(1)_A`$ charge carrier introduced earlier in $`W_{mix}`$ (15) in a general $`SU(N)`$ context (for $`SU(7)`$ there appears only one pair of such singlets, $`S`$ and $`T`$). It should be clear now that the $`W_H^{}`$ couplings (21) will rearrange the mass spectrum of the states (20), so as to leave just one pair of massless weak-doublets as needed for the MSSM. By diagonalising the $`2\times 2`$ mass matrix for the states $`H_{[cc^{}]}`$ and $`\overline{H}^{[cc^{}]}`$ and the double-coloured components $`\mathrm{\Phi }_{[cc^{}6]}`$ and $`\overline{\mathrm{\Phi }}^{[cc^{}6]}`$, the mass of the colour triplet components in (20) is found to be of order $$M_{}\frac{f\overline{f}}{y}\frac{<\phi ><\overline{\phi }>}{S}M_{GUT}$$ (22) where the combination of the primary coupling constants $`f,`$ $`\overline{f}`$and $`y`$, can be taken $`O(1)`$ in general. There is a $`3\times 3`$ mass matrix for the weak doublet states, corresponding to the mixing of the states $`H_{w6}`$, $`H_{w7}`$ and $`\mathrm{\Phi }_{[w67]}`$ and their <sup>′′</sup>conjugates<sup>′′</sup> $`\overline{H}^{w6}`$, $`\overline{H}^{w7}`$ and $`\overline{\mathrm{\Phi }}^{[w67]}`$ respectively. After diagonalisation this matrix leaves just one pair of weak-doublets $`H^{w6}`$ and $`\overline{H}^{w6}`$ strictly massless, while the other pair $`H_{w7}`$ and $`\overline{H}^{w7}`$ aquires a mass $`M_{}`$ (22) of order $`M_{GUT}`$. In much the same way all the additional states in the $`SU(7)`$ matter multiplets (18) become superheavy during the starting GUT symmetry breaking $`SU(7)SU(5)`$ . ### 3.2 Yukawa couplings The usual dimension-4 trilinear Yukawa couplings are forbidden by $`SU(7)`$ gauge invariance. So we suppose that all the generalized Yukawa couplings, the $`RP`$-conserving (ordinary up and down fermion Yukawas) as well as the $`RP`$-violating ones allowed by the $`SU(7)U(1)_A`$ symmetry, are given by a similar set of dimension-5 operators of the form ($`i,j,k=1,2,3`$ are the generation indices, the $`SU(7)`$ indices $`A,B,C=1,\mathrm{},7`$ are hereafter omitted): $$𝒪_{ij}^{up}=\frac{G_{ij}^u}{M_P}(\mathrm{\Psi }_i\mathrm{\Psi }_j)(H\phi ^{(2)})$$ (23) $$𝒪_{ij}^{down}=\frac{G_{ij}^d}{M_P}(\overline{\mathrm{\Psi }}_i\mathrm{\Psi }_j)(\overline{H}\phi ^{(1)})$$ (24) $$𝒪_{ijk}^{rpv}=\frac{G_{ijk}}{M_P}(\overline{\mathrm{\Psi }}_i\mathrm{\Psi }_j)(\overline{\mathrm{\Psi }}_k\mathrm{\Sigma }).$$ (25) Further, substituting the VEVs of the scalars $`\mathrm{\Sigma }`$ (5) and $`\phi `$ (19) into the basic operators (2325), one obtains at low energies the effective renormalisable Yukawa and LNV interactions with coupling constants $$Y_{ij}^u=G_{ij}^u\frac{<\phi ^{(2)}>}{M_P},Y_{ij}^d=G_{ij}^d\frac{<\phi ^{(1)}>}{M_P},\mathrm{\Lambda }_{ijk}=G_{ijk}\frac{<\mathrm{\Sigma }>}{M_P}.$$ (26) At the same time the baryon number non-conserving couplings $`\lambda _{ijk}^{\prime \prime }`$ completely disappear. The crucial point is that the adjoint field $`\mathrm{\Sigma }`$ develops a VEV configuration with strictly zero colour components (5) in the SUSY limit. When SUSY breaks, radiative corrections will shift the missing VEV components of $`\mathrm{\Sigma }`$ to nonzero values of order $`M_{SUSY}`$, thus inducing the ordinary $`\mu `$-term of the MSSM, on the one hand, and baryon number violating interactions with hierarchically small coupling constants of the order $`M_{SUSY}/M_{GUT}`$, on the other. The effective dimension-5 interactions (2325) could be generated by the exchange of some heavy states, such as massive string modes. When generated by the exchange of the same superheavy multiplet (that is a vector-like pair of fundamental septets $`7+\overline{7}`$), the resulting operators (24) and (25) have effective coupling constants (26) aligned in flavour space : $$\mathrm{\Lambda }_{ijk}=Y_{ij}^dϵ_k^{}$$ (27) The parameters $`ϵ_k^{}`$ ($`k=1,2,3)`$ include some known combination of the primary dimensionless coupling constants and a ratio of the VEVs of the scalars $`\mathrm{\Sigma }`$ and $`\phi `$. This relation (27) further splits into the ones for charged lepton ($`cl`$) and down quark ($`dq)`$ LNV couplings respectively, $$\lambda _{ijk}=Y_{ij}^{cl}ϵ_k,\lambda _{ijk}^{}=Y_{ij}^{dq}ϵ_k,$$ (28) when evolved from the $`SU(7)`$ scale down to low energies. So we see that the possible common origin of all the generalised Yukawa couplings, both $`RP`$-conserving and $`RP`$\- violating, at the GUT scale results in some minimal form of lepton number violation, provided that the appropriate heavy-state mediator exists. As a result, we are driven to a simple picture where the flavour structure, as well as the hierarchies of the trilinear LNV couplings in $`\mathrm{\Delta }W`$ (1), are essentially aligned with the down quark and charged lepton mass and mixing hierarchies. At the same time, the effective bilinear LNV terms appear to be generically suppressed by the custodial $`U(1)_A`$ symmetry (for a detailed exposition see a recent paper ). At low energies, the minimal LNV model presented can be viewed as an alternative to another minimal model based on the MSSM, in which only the bilinear LNV terms $`\mu _iL_iH_u`$ in $`\mathrm{\Delta }W`$ (1) are included . Depending on the $`U(1)_A`$ charges assigned to the matter and Higgs superfields involved, one can generically obtain at low energies either the bilinear model or the trilinear one considered here. The bilinear model also leads to LNV-Yukawa coupling alignment, by virtue of which many predictions of both models concerning quark flavour conservation are very similar . However, there are principal differences as well. The point is that the influence of the SUSY soft breaking sector, being predominant for the bilinear model, is quite negligible for the present one. Therefore, the LNV-Yukawa alignment, while appearing in both models, leads in the latter case to distinctive flavour-dependent relations between various LNV processes arising from slepton and squark exchanges (which are basically conditioned by the quark and lepton mass hierarchy) . By contrast, in the bilinear model these processes appear to be essentially determined by $`W`$ and $`Z`$ bozon exchanges and, as a result, are largely flavour-independent. On the other hand, the bilinear model has a serious problem of extension to the GUT framework. Any such extension leads, together with a lepton mixing with a weak Higgs doublet, to a quark mixing with a colour Higgs triplet, thus inducing baryon number violation as well. The only handle one has to address this problem seems to be the use of electroweak scale masses $`\mu _i`$ in the GUT-symmetry invariant bilinear couplings. Their use would mean that new fine-tuning conditions, besides the ordinary gauge hierarchy one, should be satisfied in a very ad hoc way. An extended discussion of the properties of the $`SU(7)`$ GUT, including the solution to the doublet-triplet splitting problem, string scale unification, proton decay, hierarchy of baryon vs lepton number violation and neutrino masses, can be found in our recent paper . ## 4 Three-family unifying GUT: $`SU(8)`$ It is tempting to treat the extra gauge symmetry in a general $`SU(N)`$ GUT as a flavour symmetry. If so, according to the particular solution (5) for the weak-component missing VEV configuration, the numbers of fundamental colours and flavours must be equal ($`n_C=n_F=k/2`$) for any even-order $`SU(N)`$ group, among which the minimal one is $`SU(8)`$ ($`n_C=n_F=3`$). Thus, in the $`SU(8)`$ case, the missing VEV configuration requires an additional colour-flavour symmetry: $`SU(3)_C`$ $``$ $`SU(3)_F`$. Having considered the basic matter superfields (quarks and leptons and their superpartners), the question of whether the above flavour symmetry $`SU(3)_F`$ is really their family symmetry naturally arises. Needless to say, among many other possibilities, the special assignment treating the families as a fundamental triplet of $`SU(3)_F`$ is the most attractive. In such a case the anomaly-free set of $`SU(8)`$ antisymmetric multiplets (in a self-evident notation; $`A,B,C=1,2,\mathrm{}.,8`$ ) $$6\overline{8}^A+\overline{28}^{[AB]}+256_{[ABC]}+70_{[ABCD]}$$ (29) is singled out, if we require that after flavour symmetry breaking only three massless families of ordinary quarks and leptons (and their superpartners) are left as chiral triplets of $`SU(3)_F`$, stemming from the multiplets $$\overline{28}=(\overline{5},\overline{3})+\mathrm{},70=(10,\overline{3})+\mathrm{}$$ (30) The remaining $`SU(5)SU(3)_F`$ components, in these as well as in the other multiplets (29), acquire heavy masses of order $`M_FM_{GUT}`$ <sup>2</sup><sup>2</sup>2The special multiplet arrangement (29) was considered before by one of us as a possible basis for the family-unifying $`SU(8)`$ GUT. Remarkably, the multiplets (29) follow from the unique (”each multiplet - one time”) set of $`SU(11)`$ multiplets after the symmetry breaking $`SU(11)SU(8)`$ and the exclusion of all the conjugated (under $`SU(8)`$) multiplets except the self-conjugated one $`70_{[ABCD]}`$.. So, one arrives at the chiral $`SU(3)_F`$ family symmetry case , which leads to a natural conservation of flavour both in the particle and sparticle sectors. Furthermore, there is a universal see-saw mechanism in the $`SU(8)`$ model, with heavy intermediate states provided by the multiplets (29), which induces non-trivial fermion mass-matrices with many texture ansätze available. So the observed pattern of quark and lepton masses and mixings can appear once the electroweak $`SU(2)U(1)_Y`$ symmetry breaks . At the same time, by analogy with the $`SU(7)`$ case, see (25), the only $`RP`$-violating coupling allowed by $`SU(8)U(1)_A`$ symmetry is supposed to be given by the dimension-5 operator $$𝒪_{rpv}\frac{1}{M_P}(\overline{\mathrm{\Psi }}^{[AB]}\mathrm{\Psi }_{[ABCD]}\overline{\mathrm{\Psi }}^{[CD^{}]}\mathrm{\Sigma }_D^{}^D)$$ (31) Here the matter fields $`\overline{\mathrm{\Psi }}`$ and $`\mathrm{\Psi }`$ belong to the basic multiplets (30). One can see now that the weak-component missing VEV solution for $`\mathrm{\Sigma }`$ (5), when substituted into the operator $`𝒪_{rpv}`$, leaves only the LNV couplings and projects out the baryon number violating ones. At low energies the surviving effective couplings take the form $$\lambda ϵ_{\alpha \beta \gamma }(L^\alpha L^\beta \overline{E}^\gamma +rL^\alpha Q^\beta \overline{D}^\gamma )$$ (32) Here $`\alpha ,`$ $`\beta ,`$ $`\gamma =1,2,3`$ are the generation indices, belonging to the family $`SU(3)_F`$ symmetry, and $`r`$ is a factor giving the relative coupling constant renormalisation after evolution from the $`SU(8)`$ scale to low energies. So, as in the $`SU(7)`$ case, one has baryon number conservation at the same time as lepton number violation in the SUSY limit. Meanwhile, despite their common origin, there is a principal difference between the $`SU(7)`$ and $`SU(8)`$ cases. The point is that the basic adjoint $`\mathrm{\Sigma }`$ moduli mass ratio $`M_3/M_8`$ appears, according to the missing VEV vacua (5), to be 2 and 1/2 for $`SU(7)`$ and $`SU(8)`$ respectively. As was shown in recent papers , this ratio essentially determines the high-energy behavior of the MSSM gauge couplings. In fact it follows that the unification scale in $`SU(7)`$ is pushed to the string scale , while the unification scale in $`SU(8)`$ ranges, at best, close to the standard unification value . ## 5 Conclusions The absence of automatic global conservation laws in SUSY theories, in contrast to the Standard Model, is frequently considered as a drawback of supersymmetry. Meanwhile phenomenologically, whereas SUSY-inspired B number non-conservation must be highly suppressed, SUSY-inspired L-number violation could occur at a level large enough for the observation of its many spectacular manifestations . One of these manifestations may be the sizeable atmospheric neutrino oscillations recently reported , according to which one of the neutrino species is expected to have a mass at least of order $`0.1`$ eV. That means, in general, the particle content of the MSSM or the minimal $`SU(5)`$ SUSY GUT should be extended to include new states, that is fundamentally heavy right-handed neutrinos or even light sterile left-handed ones. Neutrino masses per se do not yet give any conclusive evidence in favour of SUSY theories. However sizeable LNV in the charged lepton sector and, of course, in the decays of the lightest supersuymmetric particle , if actually observed, could qualify as generic SUSY inspired phenomena. In such a situation the following question would arise, which should be addressed within the framework of Grand Unification rather than the MSSM: what could stand behind such a tremendous hierarchy of lepton vs baryon number violation? In this connection we suggested that the nature of the global conservation laws in SUSY theories is determined by the basic vacuum configuration which breaks the underlying GUT symmetry. Following this idea, we have argued that the GUTs with a natural missing VEV solution to the doublet-triplet splitting problem could, simultaneously, provide the reason for treating lepton and baryon number carrying matter fields differently. We have shown that missing VEV vacuum configurations, ensuring the survival of the MSSM gauge symmetry at low energies, only emerge in extended SU(N) GUTs with $`N7`$. Further, the one-family unifying $`SU(7)`$ and the three-family unifying $`SU(8)`$ GUTs have been constructed. In both cases the effective LNV couplings immediately evolve from the GUT scale, while the baryon number non-conserving ones are safely projected out by the missing VEV vacuum configuration, which breaks the starting GUT symmetry down to that of the MSSM. However, at the next stage when SUSY breaks, radiative corrections shift the missing VEV to some nonzero value of order $`M_{SUSY}`$, thus inducing the ordinary $`\mu `$-term of the MSSM, on the one hand, and BNV couplings with the hierarchically small constants $`\lambda _{ijk}^{\prime \prime }=O(M_{SUSY}/M_{GUT})`$, on the other. So, a missing VEV solution to the gauge hierarchy problem leads, in a literal sense, to the same hierarchy of baryon vs lepton number violation. ## Acknowledgments We would like to thank many of our colleagues, especially Riccardo Barbieri, Grahame Blair, Ilia Gogoladze, Mike Green, David Hutchcroft, Archil Kobakhidze, Gordon Moorhouse, Alexei Smirnov and David Sutherland for stimulating dicussions and useful remarks. Financial support by the INTAS grants No. RFBR 95-567, 96-155, PPARC grant No. PPA/V/S/1997/00644 and the Joint Project grant from the Royal Society are also gratefully acknowledged.
warning/0004/cond-mat0004161.html	ar5iv	text	# 1 Introduction ## 1 Introduction The $`q`$-state Potts antiferromagnet (AF) exhibits nonzero ground state entropy, $`S_0>0`$ (without frustration) for sufficiently large $`q`$ on a given lattice $`\mathrm{\Lambda }`$ or, more generally, on a graph $`G`$. This is equivalent to a ground state degeneracy per site $`W>1`$, since $`S_0=k_B\mathrm{ln}W`$. Such nonzero ground state entropy is important as an exception to the third law of thermodynamics . There is a close connection with graph theory here, since the zero-temperature partition function of the above-mentioned $`q`$-state Potts antiferromagnet on a graph $`G`$ satisfies $$Z(G,q,T=0)_{PAF}=P(G,q)$$ (1.1) where $`P(G,q)`$ is the chromatic polynomial expressing the number of ways of coloring the vertices of the graph $`G`$ with $`q`$ colors such that no two adjacent vertices have the same color (for reviews, see -). The minimum number of colors necessary for such a coloring of $`G`$ is called the chromatic number, $`\chi (G)`$. Thus $$W(\{G\},q)=\underset{n\mathrm{}}{lim}P(G,q)^{1/n}$$ (1.2) where $`n=v(G)`$ is the number of vertices of $`G`$ and $`\{G\}=lim_n\mathrm{}G`$. At certain special points $`q_s`$ (typically $`q_s=0,1,..,\chi (G)`$), one has the noncommutativity of limits $$\underset{qq_s}{lim}\underset{n\mathrm{}}{lim}P(G,q)^{1/n}\underset{n\mathrm{}}{lim}\underset{qq_s}{lim}P(G,q)^{1/n}$$ (1.3) and hence it is necessary to specify the order of the limits in the definition of $`W(\{G\},q_s)`$ . Denoting $`W_{qn}`$ and $`W_{nq}`$ as the functions defined by the different order of limits on the left and right-hand sides of (1.3), we take $`WW_{qn}`$ here; this has the advantage of removing certain isolated discontinuities that are present in $`W_{nq}`$. Using the expression for $`P(G,q)`$, one can generalize $`q`$ from $`_+`$ to $``$. The zeros of $`P(G,q)`$ in the complex $`q`$ plane are called chromatic zeros; a subset of these may form an accumulation set in the $`n\mathrm{}`$ limit, denoted $``$, which is the continuous locus of points where $`W(\{G\},q)`$ is nonanalytic. <sup>1</sup><sup>1</sup>1For some families of graphs $``$ may be null, and $`W`$ may also be nonanalytic at certain discrete points. The maximal region in the complex $`q`$ plane to which one can analytically continue the function $`W(\{G\},q)`$ from physical values where there is nonzero ground state entropy is denoted $`R_1`$. The maximal value of $`q`$ where $``$ intersects the (positive) real axis is labelled $`q_c(\{G\})`$. This point is important since it separates the interval $`q>q_c(\{G\})`$ on the positive real $`q`$ axis where the Potts model (with $`q`$ extended from $`_+`$ to $``$) exhibits nonzero ground state entropy (which increases with $`q`$, asymptotically approaching $`S_0=k_B\mathrm{ln}q`$ for large $`q`$, and which for a regular lattice $`\mathrm{\Lambda }`$ can be calculated approximately via large–$`q`$ series expansions) from the interval $`0qq_c(\{G\})`$ in which $`S_0`$ has a different analytic form. In the present work we report exact solutions for chromatic polynomials $`P(G,q)`$ for strips of the square lattice with arbitrarily great length $`L_x`$ vertices of the following types: (a) width $`L_y=4`$ vertices and $`(FBC_y,PBC_x)=`$ cyclic; (b) $`L_y=4`$ and $`(FBC_y,TPBC_x)=`$ Möbius, (c) $`L_y=5`$ and $`(PBC_y,FBC_x)=`$ cylindrical; and (d) $`L_y=5`$ and $`(FBC_y,FBC_x)=`$ open, where $`FBC`$, $`PBC`$, and $`TPBC`$ denote free, periodic, and twisted periodic (i.e. periodic with reversed orientation) boundary conditions, respectively. For each of these, taking the infinite-length limit, we calculate the degeneracy per site, $`W(\{G\},q)`$, and the continuous nonanalytic locus $``$. A comparative discussion is given of these results together with previous exact solutions for strips of smaller widths. These strips of regular lattices are examples of recursive families of graphs, where the latter are constructed by successive additions of subgraph units to an initial subgraph. There are several motivations for this work. We have mentioned the basic importance of nonzero ground state entropy in statistical mechanics . Physical examples are provided by ice - and certain other hydrogen-bonded molecular crystals . From the point of view of rigorous statistical mechanics, exact solutions are always valuable for the insight that they give into the behavior of the given system under study. Although infinite-length finite-width strips are quasi-one-dimensional systems and hence (for finite-range spin-spin interactions) do not have finite-temperature phase transitions, their zero-temperature critical points are of interest. Indeed, the presence of a zero-temperature critical point for the Ising antiferromagnet<sup>2</sup><sup>2</sup>2We recall that on bipartite graphs such as cyclic strips of the square lattice with even $`L_x`$, an elementary mapping shows the Ising ferromagnet and antiferromagnet to be equivalent; since the $`L_x\mathrm{}`$ limit can be taken with even $`L_x`$, this implies that the critical behavior of the Ising ferromagnet is equivalent to that of the Ising antiferromagnet on infinite-length limits cyclic strips of the square lattice. By similar elementary reasoning, one can show this equivalence for the infinite-length limit of Möbius strips of the square lattice. on infinite-length, finite-width strips of the square lattice has an interesting connection with the behavior of the singular locus $``$ for the strips that we have studied with global circuits<sup>3</sup><sup>3</sup>3A global circuit is a route following a lattice direction which has the topology of the circle, $`S^1`$, and a length $`\mathrm{}_{g.c.}`$ that goes to infinity as $`n\mathrm{}`$. For strip graphs, global circuits are equivalent to periodic or twisted periodic boundary conditions.: in these cases, this singular locus passes through the point $`q=2`$ in the complex $`q`$ plane. Our exact solutions for $`P`$ and, in the $`L_x\mathrm{}`$ limit, $`W`$, thus show quantitatively the relation between critical behavior as a function of temperature (at $`T=0`$) in the free energy and singularities as a function of $`q`$ in the per-site ground state degeneracy $`W`$. The present results also show many interesting connections with mathematical graph theory, as is clear from the identity (1.1), and algebraic geometry, as follows from the fact that for these strips, $``$ is an algebraic curve. Besides the works already cited, some related work on chromatic polynomials for recursive graphs includes -; further discussion of background and references may be found in . A generic form for chromatic polynomials for a strip graph of type $`G_s`$, width $`L_y`$, and length $`L_x`$ is $$P(G_s,L_y\times L_x,BC_y,BC_x,q)=\underset{j=1}{\overset{N_\lambda }{}}c_j(q)(\lambda _j(q))^{L_x}$$ (1.4) where $`c_j(q)`$ and the $`N_\lambda `$ terms $`\lambda _j(q)`$ depend on the type of strip graph $`G_s`$ including the boundary conditions but are independent of $`L_x`$. ## 2 $`L_y=4`$ Square-Lattice Strips with $`(FBC_y,(T)PBC_x)`$ In this section we give our solutions for the chromatic polynomials of the $`L_y\times L_x`$ strips of the square lattice with $`(FBC_y,PBC_x)`$ and $`(FBC_y,TPBC_x)`$, i.e. cyclic and Möbius, boundary conditions, respectively. For both the cyclic and Möbius strips, for $`L_x3`$ to avoid certain degenerate cases, the square lattice strips of width $`L_y`$ have $`n=L_xL_y`$ vertices and $`e=L_x(2L_y1)`$ edges. The cyclic square strips have $`\chi =2`$ for $`L_x`$ even and $`\chi =3`$ for $`L_x`$ odd, independent of $`L_y`$. For Möbius square strips, $`\chi =2`$ for $`(L_x,L_y)=(e,o)`$ or $`(o,e)`$, and $`\chi =3`$ for $`(L_x,L_y)=(e,e)`$ or $`(o,o)`$, where $`e`$ and $`o`$ denote even and odd. We calculate the chromatic polynomials by iterated use of the deletion-contraction theorem , together with coloring matrix methods . The calculation is considerably more involved than that for the $`L_y=3`$ cyclic strip given in , as is indicated by the number of $`\lambda _j`$ terms in eq. (1.4), namely, $`N_\lambda =26`$, as contrasted with the value $`N_\lambda =10`$ for the $`L_y=3`$ cyclic strip. Elsewhere we have given a general determination of $`N_\lambda `$ as a function of $`L_y`$ . As $`L_y`$ increases, the number of terms $`N_\lambda `$ in (1.4) grows rapidly; it is 70, 192, and 534 for $`L_y=5,6,`$ and 7. We obtain the exact solutions of the form (1.4) $$P(sq(4\times m,FBC_y,PBC_x),q)=\underset{j=1}{\overset{26}{}}c_{sq4,j}(\lambda _{sq4,j})^m$$ (2.1) and $$P(sq(4\times m,FBC_y,TPBC_x),q)=\underset{j=1}{\overset{26}{}}c_{sq4Mb,j}(\lambda _{sq4,j})^m$$ (2.2) where $`L_x=m`$. The fact that the $`\lambda _j`$’s for a Möbius strip must be the same as those for the cyclic strip of the same width and lattice type was proved in ; this also proves, a fortiori, that (i) the total number, $`N_\lambda `$, of $`\lambda _j`$’s, and (ii) the continuous nonanalytic locus $``$, including the point $`q_c`$, are the same for the cyclic and Möbius strips of a given type. For $``$ and $`q_c`$, we shall often indicate this by the notation $`(FBC_y,(T)PBC_x)`$. The explicit $`\lambda _{sq,j}`$’s that we calculate are as follows. The first six are $$\lambda _{sq4,1}=1$$ (2.3) $$\lambda _{sq4,2}=3q$$ (2.4) $$\lambda _{sq4,3}=1q$$ (2.5) $$\lambda _{sq4,(4,5)}=(3\pm \sqrt{2})q$$ (2.6) and $$\lambda _{sq4,6}=(q1)(q3).$$ (2.7) Three of the remaining $`\lambda _{sq4,j}`$’s, labelled $`j=7,8,9`$, including the one that is dominant in region $`R_1`$, are identical to the three that enter into the chromatic polynomial for the $`L_y=4`$ strip with $`(FBC_y,FBC_x)`$ boundary conditions, which was previously calculated in . This identity was shown in . The remaining $`\lambda _{sq4,j}`$’s for $`10j26`$ are roots of another cubic equation, for $`j=10,11,12`$; a quartic equation for $`13j16`$; and two 5th degree equations, for $`17j26`$. Since the equations defining these $`\lambda _j`$’s are somewhat lengthy, we give them in the appendix. In Table 1 we list various properties of our calculation and compare them with the properties that we have found for other related strips of the square (and triangular) lattice. The results for the $`L_y\mathrm{}`$ limit for the triangular lattice with $`(PBC_y,FBC_x)`$ are from . Comparisons for other lattices such as honeycomb and kagomé were given in . In particular, the fact that for this width, $`L_y=4`$ for the cyclic strip of the square lattice, we encounter equations of degree 5 for the $`\lambda _j`$’s means that it is not possible to solve for the corresponding $`\lambda _j`$’s as algebraic roots. Our experience with lattice strips of a given width $`L_y`$ (and arbitrary length) and a given set of boundary conditions is that the maximal degrees of the factors in the general equation for the $`\lambda _j`$’s are non-decreasing functions of $`L_y`$. Thus, assuming that this property of non-decreasing degrees of algebraic factors in the equation for the $`\lambda _j`$’s continues for higher $`L_y`$, our present results indicate that the exact solutions in for the $`\lambda _j`$’s for the width $`L_y=3`$ strip of the square lattice have completed the program of obtaining exact algebraic expressions for these terms for this type of lattice strip. Although no closed-form algebraic expression can be obtained for the $`\lambda _j`$’s, a theorem on symmetric polynomials of roots of algebraic equations, discussed in , enables one to calculate the chromatic polynomials exactly to arbitrary order. The key to this is the property that since the chromatic polynomial for a cyclic strip is a symmetric polynomial in the various roots, it can be expressed in terms of the coefficients of the algebraic equations that determine these $`\lambda _j`$’s. However, the fact that it is no longer possible to calculate the $`\lambda _j`$’s as algebraic roots when the width of the cyclic square strip is 4 means that the determination of the nonanalytic locus $``$ must be done in a somewhat more cumbersome manner than in our previous work where we had exact algebraic expressions for these $`\lambda _j`$’s. The coefficients $`c_j`$ that enter into the expressions for the chromatic polynomial (1.4) for the cyclic and Möbius strip of the square lattice of width $`L_y`$ are certain polynomials that we denote $`c^{(d)}`$, given by $$c^{(d)}=\underset{k=1}{\overset{d}{}}(qq_{d,k})$$ (2.8) where $$q_{d,k}=2+2\mathrm{cos}\left(\frac{2\pi k}{2d+1}\right)\mathrm{for}k=1,2,\mathrm{}d$$ (2.9) with $`0dL_y`$. We list below the specific $`c^{(d)}`$’s that appear in our results for the $`L_y=4`$ square lattice strip: $$c^{(0)}=1,c^{(1)}=q1,c^{(2)}=q^23q+1,$$ (2.10) $$c^{(3)}=q^35q^2+6q1,$$ (2.11) and $$c^{(4)}=(q1)(q^36q^2+9q1).$$ (2.12) In ascending order of degrees of $`c^{(d)}`$, we calculate $$c_{sq4,j}=c^{(0)}\mathrm{for}6j9$$ (2.13) $$c_{sq4,j}=c^{(1)}\mathrm{for}13j21$$ (2.14) $$c_{sq4,j}=c^{(2)}\mathrm{for}10j12\mathrm{and}22j26$$ (2.15) $$c_{sq4,j}=c^{(3)}\mathrm{for}2j5$$ (2.16) and $$c_{sq4,1}=c^{(4)}.$$ (2.17) In it was shown that the coefficient for the $`\lambda _j`$ that is leading in region $`R_1`$ must be 1. We define $$C(G)=\underset{j=1}{\overset{N_{\lambda _G}}{}}c_{G,j}.$$ (2.18) where the $`G`$-dependence in the coefficients is indicated explicitly. Note that for recursive graphs like the strip graphs considered here, the $`c_{G,j}`$ depend on $`L_y`$ and the boundary conditions, but not on $`L_x`$. Our results above give $$C(G)=q(q1)^3\mathrm{for}G=sq(L_y=4,FBC_y,PBC_x).$$ (2.19) in accord with the generalization $$C(G_s(L_y\times L_x,FBC_y,PBC_x),q)=P(T_{L_y},q)=q(q1)^{L_y1}$$ (2.20) for $`G_s`$ a strip of the square (or triangular) lattice, where $`P(T_n,q)`$ is the chromatic polynomial for the tree graph $`T_n`$. This is in accord with the coloring matrix approach . For the $`L_y=4`$ Möbius strip of the square lattice, we find $$c_{sq4Mb,j}=c^{(0)}\mathrm{for}7j12$$ (2.21) $$c_{sq4Mb,j}=c^{(0)}\mathrm{for}j=6\mathrm{and}22j26$$ (2.22) $$c_{sq4Mb,j}=c^{(1)}\mathrm{for}17j21$$ (2.23) $$c_{sq4Mb,j}=c^{(1)}\mathrm{for}j=1\mathrm{and}13j16$$ (2.24) $$c_{sq4Mb,j}=c^{(2)}\mathrm{for}j=4,5$$ (2.25) and $$c_{sq4Mb,j}=c^{(2)}\mathrm{for}j=2,3.$$ (2.26) Hence, the sum of the coefficients is $$C(G)=0\mathrm{for}G=sq(L_y=4,FBC_y,TPBC_x)$$ (2.27) in accord with the general result for the Möbius strip of the square (and triangular) lattice $$\underset{j=1}{\overset{N_{\lambda _G}}{}}c_{G(L_y,Mb),j}=\{\begin{array}{cc}P(T_{\frac{L_y+1}{2}},q)\hfill & \text{for odd }L_y\hfill \\ 0\hfill & \text{for even }L_y\hfill \end{array}.$$ (2.28) Chromatic zeros for the cyclic strip of the square lattice with $`L_y=4`$, $`L_x=m=20`$ and hence $`n=80`$ are shown in Fig. 1; with this value of $`m`$, the complex chromatic zeros lie close to the boundary $``$ and give an approximate indication of its position. Note that there is a zero very close to $`q=2`$, but $`P(sq(L_y\times L_x,FBC_y,PBC_x),q)`$ is nonzero for $`q=2`$ for the case shown, where $`L_x=m`$ is even, as is clear from the fact that $`\chi =2`$ in this case. The maximal point at which $``$ crosses the real axis, $`q_c`$, is determined as a solution of the equation of degeneracy of leading terms $`\|\lambda _{eq79,max}\|=\|\lambda _{eq2226,max}\|`$, where $`\lambda _{eq79,max}`$ and $`\lambda _{eq2226,max}`$ are the roots of eqs. (8.1.7) and (8.1.39) with the largest magnitudes, respectively. Since only two $`\lambda _j`$’s are degenerate in magnitude at this point, it is a regular point on the algebraic curve $``$ in the terminology of algebraic geometry. This is also the case for the $`L_y=3`$ (and $`L_y=1`$) strip of the square lattice , whereas, in contrast, $`q_c`$ is a multiple point on $``$ for $`L_y=2`$. We find $$q_c2.4928456\mathrm{for}sq(4\times \mathrm{},FBC_y,(T)PBC_x)$$ (2.29) This may be compared with the values $`q_c=2`$ for the $`L_y\times \mathrm{}`$ strip of the square lattice with $`L_y=1,2`$ and the same $`(FBC_y,(T)PBC_x)`$ boundary conditions , and the value $`q_c2.33654`$ for $`L_y=3`$ . We calculate that $`W(sq,4\times \mathrm{},FBC_y,BC_x)=1.2697336..`$ at the value $`q=q_c`$ in eq. (2.29). The locus $``$ also crosses the real $`q`$ axis at $`q=2`$ and at $`q=0`$. In addition to region $`R_1`$ which extends outward from the envelope of $``$ and includes the real axis for $`q>q_c`$ and $`q<0`$, there are two other regions that contain segments of the real axis: $`R_2`$, including the interval $`2<q<q_c`$ and $`R_3`$, including the interval $`0<q<2`$. In region $`R_1`$, the dominant $`\lambda _j`$ is the root of the cubic equation (8.1.7) with the largest magnitude, which we label $`\lambda _{79,max}`$. In region $`R_2`$, the dominant $`\lambda _j`$ is the root of the fifth-degree equation (8.1.39) with the largest magnitude, which we label $`\lambda _{2226,max}`$. In region $`R_3`$, the dominant $`\lambda _j`$ is the root of the fifth-degree equation (8.1.30) with the largest magnitude, which we label $`\lambda _{1721,max}`$. We have $$W=(\lambda _{79,max})^{1/4},\mathrm{for}qR_1$$ (2.30) $$\|W\|=\|\lambda _{2226,max}\|^{1/4},\mathrm{for}qR_2$$ (2.31) $$\|W\|=\|\lambda _{1721,max}\|^{1/4},\mathrm{for}qR_3$$ (2.32) (In regions other than $`R_1`$, only the magnitude $`\|W\|`$ can be determined unambiguously .) The locus $``$ has support for $`Re(q)<0`$ as well as $`Re(q)0`$. It separates the $`q`$ plane into several regions, including the three described above and two complex-conjugate ones which we denote $`R_4`$ and $`R_4^{}`$, centered approximately at $`q2.6\pm 0.8i`$. In the regions $`R_4`$ and $`R_4^{}`$, we have $$\|W\|=\|\lambda _{1316,max}\|^{1/4},\mathrm{for}qR_4,R_4^{}$$ (2.33) Just as complex-conjugate pairs of tiny sliver regions were found for the cyclic $`L_y=3`$ square and triangular strips, so also these may be present here; we have not carried out a search for such regions (but have ruled out the possibility of tiny regions on the real axis). ## 3 $`L_y=5,6`$ Square-Lattice Strips with $`(PBC_y,FBC_x)`$ Here we report our exact solutions for the chromatic polynomials for the width $`L_y=5,6`$ strips of the square lattice of arbitrary length and with $`(PBC_y,FBC_x)`$, i.e., cylindrical, boundary conditions. Results for the cases $`L_y=3`$ and $`L_y=4`$ were given previously in . We recall that $`N_\lambda =1`$ for $`L_y=3`$ and $`N_\lambda =2`$ for $`L_y=4`$. For $`L_y=5`$ and $`L_y=6`$ we calculate $`N_\lambda =2`$ and $`N_\lambda =5`$, respectively. As before, it is convenient to present the results in terms of a generating function, denoted $`\mathrm{\Gamma }(G_s,q,x)`$. The chromatic polynomial $`P((G_s)_m,q)`$ is determined as the coefficient in a Taylor series expansion of this generating function in an auxiliary variable $`x`$ about $`x=0`$: $$\mathrm{\Gamma }(G_s,q,x)=\underset{m=0}{\overset{\mathrm{}}{}}P((G_s)_m,q)x^m.$$ (3.1) The generating function $`\mathrm{\Gamma }(G_s,q,x)`$ is a rational function of the form $$\mathrm{\Gamma }(G_s,q,x)=\frac{𝒩(G_s,q,x)}{𝒟(G_s,q,x)}$$ (3.2) with $$𝒩(G_s,q,x)=\underset{j=0}{\overset{d_𝒩}{}}A_{G_s,j}(q)x^j$$ (3.3) and $$𝒟(G_s,q,x)=1+\underset{j=1}{\overset{d_𝒟}{}}b_{G_s,j}(q)x^j$$ (3.4) where the $`A_{G_s,i}`$ and $`b_{G_s,i}`$ are polynomials in $`q`$, and $`d_𝒩deg_x(𝒩)`$, $`d_𝒟deg_x(𝒟)`$, In factorized form $$𝒟(G_s,q,x)=\underset{j=1}{\overset{d_𝒟}{}}(1\lambda _{G_s,j}(q)x).$$ (3.5) Equivalently, the $`\lambda _{G_s,j}`$ are roots of the equation $$\xi ^{d_𝒟}𝒟(G_s,q,1/\xi )=\xi ^{d_𝒟}+\underset{j=1}{\overset{d_𝒟}{}}b_{G_s,j}\xi ^{d_𝒟j}.$$ (3.6) The general formula expressing $`P(G_m,q)`$ in terms of these quantities is $$P(G_m,q)=\underset{j=1}{\overset{d_𝒟}{}}\left[\underset{s=0}{\overset{d_𝒩}{}}A_s\lambda _j^{d_𝒟s1}\right]\left[\underset{1id_𝒟;ij}{}\frac{1}{(\lambda _j\lambda _i)}\right]\lambda _j^m.$$ (3.7) For $`L_y=5`$ we find $$\lambda _{sq5PF,j}=\frac{1}{2}\left[T_{sq5PF}\pm \sqrt{R_{sq5PF}}\right],j=1,2$$ (3.8) where $$T_{sq5PF}=q^510q^4+46q^3124q^2+198q148$$ (3.9) and $`R_{sq5PF}=q^{10}20q^9+188q^81092q^7+4356q^612596q^5+27196q^4`$ (3.10) (3.11) $`44212q^3+52708q^241760q+16456.`$ (3.12) The coefficients $`c_{sq5PF,j}`$ can be computed using eq. (3.7) in terms of the generating function, which is given in the appendix. In the $`L_x\mathrm{}`$ limit, the locus $``$ includes five arcs, consisting of two complex-conjugate pairs and a fifth, self-conjugate, arc. The endpoints of these arcs are located at the five complex-conjugate pairs of roots of $`R_{sq5PF}`$. The self-conjugate arc crosses the real axis at the real zero of $`T_{sq5PF}`$, namely at $$q_c2.691684\mathrm{for}sq(5\times \mathrm{},PBC_y,FBC_x)$$ (3.13) In Fig. 3 we show a plot of chromatic zeros for the $`L_y=6`$ strip of the square lattice with $`(PBC_y,FBC_x)`$ and length $`L_x=m+2=16`$ vertices, so that the strip has $`n=96`$ vertices in all. With this large a value of $`m`$, the complex chromatic zeros lie close to the boundary $``$ and give an approximate indication of its position. (We have not searched for very minute features in $``$.) From our exact analytic results, we calculate indication of its position. From our exact analytic results, we calculate $$q_c2.6089\mathrm{for}sq(6\times \mathrm{},PBC_y,FBC_x)$$ (3.14) The morphology of chromatic zeros for this long $`6\times 16`$ cylindrical strip is similar to that found for a $`8\times 8`$ patch of the square lattice, again with cylindrical boundary conditions, in . In both cases, the chromatic zeros have support for $`Re(q)<0`$ and prongs extending to the right; further, our exact calculation shows that in the limit $`L_x\mathrm{}`$ with $`L_y=6`$, the locus $``$ has support for $`Re(q)<0`$. In Fig. 3, one of the chromatic zeros is very close to $`q=2`$, but for $`q=2`$ exactly, the chromatic polynomial is nonzero, equal to 2, in accord with the fact that this strip is bipartite for any value of $`L_x`$. For the $`L_x\mathrm{}`$ limit of these respective strips we have $$W(sq(5\times \mathrm{},PBC_y,FBC_x),q)=(\lambda _{sq5PF,j,max})^{1/5}$$ (3.15) and $$W(sq(6\times \mathrm{},PBC_y,FBC_x),q)=(\lambda _{sq6PF,j,max})^{1/6}$$ (3.16) where $`\lambda _{sq5PF,j,max}`$ and $`\lambda _{sq6PF,j,max}`$ denote the solutions to the respective equations (3.6) with maximal magnitude in region $`R_1`$. It is of interest to use this exact result to study further the approach of $`W`$ to the limit for the full infinite 2D square lattice. This extends our previous study in . In Table 2 we list various values of $`W(sq(L_y\times \mathrm{},PBC_y,BC_x),q)`$ (which, for this range of $`q`$, are independent of $`BC_x`$), denoted as $`W(sq(L_y),P,q)`$, together with Monte Carlo measurements of $`W`$ for the full 2D square lattice, $`W(sq,q)`$ from and the $`q=3`$ value $`W(sq,3)=(4/3)^{3/2}`$ from . We also list the ratio $$R_W(\mathrm{\Lambda }(L_y),BC_y,q)=\frac{W(\mathrm{\Lambda }(L_y),BC_y,q)}{W(\mathrm{\Lambda },q)}$$ (3.17) for the present square lattice $`\mathrm{\Lambda }=sq`$. One sees that for $`L_y=5`$ and moderate values of $`q`$, say 5 or 6, the agreement of $`W(sq(L_y),q)`$ for the infinite-length, finite-width strips with the respective values $`W(sq,q)`$ for the infinite square lattice is excellent; the differences are of order $`10^3`$ to $`10^4`$. As noted before , for $`PBC_y`$ (and any $`BC_x`$) this approach is not monotonic. ## 4 $`L_y=5`$ Square-Lattice Strips with $`(FBC_y,FBC_x)`$ We have also gone beyond the previous studies in to calculate the chromatic polynomial for the strip of the square lattice with width $`L_y=5`$ and $`(FBC_y,FBC_x)`$, i.e., open, boundary conditions. A related study on wide strips is in . In , a given strip $`(G_s)_m`$ was constructed by $`m`$ successive additions of a subgraph $`H`$ to an endgraph $`I`$; here, $`I=H`$, so that, following the notation of , the total length of the strip graph $`(G_s)_m`$ is $`L_x=m+2`$ vertices, or equivalently, $`m+1`$ edges in the longitudinal direction. The results are conveniently expressed in terms of the coefficient functions in the generating function, as discussed above. For the width $`L_y=5`$ strip of the square lattice we find $`deg_x(𝒟)=N_\lambda =7`$. The coefficient functions $`b_{sq5FF,j}`$ in eq. (3.4) that determine the $`\lambda _{sq5FF,j}`$’s via eq. (3.6) are listed in the appendix. Because the $`A_{sq5FF,j}`$’s (cf. eq. (3.3)) are quite lengthy, we do not give them here<sup>4</sup><sup>4</sup>4The $`A_{sq5FF,j}`$ are listed in the copy of this paper in the cond-mat archive.. In Table 1 this result is compared with the findings from the previous calculations in for narrower open strips of the square lattice, and with strips of the triangular lattice . One observes that the equation (3.6) defining the $`\lambda _j`$’s increases in degree as $`L_y`$ increases for the open strips. In particular, because we now encounter an equation (3.6) of degree higher than 4 (specifically, degree 7), it is not possible to solve for the $`\lambda _j`$’s as algebraic roots. Furthermore, assuming that this increase in degree of (3.6) continues for greater widths $`L_y`$ of open strips, our present results show that the previous calculations of the $`\lambda _j`$’s in up to $`L_y=4`$ have completed the program of calculating these terms exactly as algebraic roots for open strips of the square lattice. As noted above, because of the theorem on symmetric polynomial functions of roots an algebraic equations , one can still calculate the chromatic polynomial in terms of the coefficients of the algebraic equation for the $`\lambda _j`$’s. In Fig. 4 we show a plot of chromatic zeros for the open strip of the square lattice with $`L_y=5`$ and length $`L_x=m+2=16`$ vertices, so that the strip has $`n=80`$ vertices in all. With this large a value of $`m`$, the complex chromatic zeros lie close to the boundary $``$ and give an approximate indication of its position. From an analysis of the degeneracy of leading $`\lambda _j`$’s, we find that (in the $`L_x\mathrm{}`$ limit where $``$ is defined) $$q_c2.42843\mathrm{for}sq(5\times \mathrm{},FBC_y,FBC_x)$$ (4.1) This is in agreement with the chromatic zeros shown in Fig. 4. Comparing Fig. 4 with the corresponding plots for $`L_y=2`$ and $`L_y=3`$ (Fig. 3(a,b) of ), we see that the arcs forming $``$ are elongating and that the arc endpoints nearest to the origin are approaching more closely to the origin. This agrees with the behavior that we had observed earlier from narrower strips and with the conclusions that were drawn from that behavior , in particular, the statement that these results are consistent with, and provide further support for, the hypothesis that in the limit as $`L_y\mathrm{}`$, the locus $``$ will extend all the way through the origin of the $`q`$ plane and will separate this plane into different regions containing the real axis. For $`q>q_c`$, we have, for the physical ground state degeneracy per site of the $`q`$-state Potts antiferromagnet, $$W(sq(5\times \mathrm{},FBC_y,FBC_x),q)=(\lambda _{sq5FF,j,max})^{1/5}$$ (4.2) where $`\lambda _{sq5FF,j,max}`$ denotes the solution of eq. (3.6) with the coefficients (8.3.1)-(8.3.29) that has the maximal magnitude in region $`R_1`$. As with the cylindrical strips, we can use our new exact solution for the $`L_y=5`$ open square strip to study the approach of $`W`$ to the limit for the infinite 2D square lattice, extending . In Table 3 we list various values of $`W(sq(L_y\times \mathrm{},FBC_y,BC_x),q)`$ (which, for this range of $`q`$, are independent of $`BC_x`$), denoted as $`W(sq(L_y),F,q)`$, together with Monte Carlo measurements of $`W`$ for the full 2D square lattice, $`W(sq,q)`$ from and the $`q=3`$ value $`W(sq,3)=(4/3)^{3/2}`$ from . We also list the ratio $`R_W(sq(L_y),FBC_y,q)`$ defined in (3.17). In it was proved that for $`FBC_y`$ the approach of $`W`$ to the $`L_y=\mathrm{}`$ limit is monotonic. One sees from Table 3 that for $`L_y=5`$ and moderate values of $`q`$, say 5 or 6, the agreement of $`W(sq(L_y),q)`$ for the infinite-length, finite-width strips with the respective values $`W(sq,q)`$ for the infinite square lattice is very good, accurate to a few per cent, although the approach is somewhat slower for open strips than for cylindrical strips. This is understandable since the condition of periodic boundary conditions in the transverse direction minimizes finite-size effects in this direction. ## 5 Comparative Discussion on $`P`$ and $``$ In this section we give a general discussion of some properties of (i) the chromatic polynomials for cyclic lattice strips with both arbitrarily great length and arbitrarily great width, and (ii) the loci $``$ for the infinite-length limit of strips of the square and triangular lattice with various boundary conditions. This discussion incorporates the exact solutions given in the present work and also in our previous papers. 1. From our exact solutions for cyclic and Möbius strips of the square and triangular lattices, we draw the following inference: for these lattice strips, with arbitrary $`L_y`$ (independent of $`L_x`$), the $`\lambda _{G_s,j}`$ in eq. (1.4) with the highest-degree $`c^{(d)}`$, namely $`c^{(L_y)}`$ (see eq. (2.8)), is $$\lambda _{G_s,1}=(1)^{L_y}$$ (5.1) 2. A second inference concerns the set of terms $`\lambda _{G_s,j}`$ for the cyclic strips of the square and triangular lattice with coefficients $`c_{G_s,j}=c^{(L_y1)}`$. There are $`L_y`$ of these terms $`\lambda _{G_s,j}`$ . Let $`\overline{\lambda }_{G_x,L_y,j}=(1)^{L_y}\lambda _{G_s,L_y,j}`$ for $`G_s=sq,tri`$. Then the $`\overline{\lambda }_{sq,L_y,j}`$’s and hence the $`\lambda _{sq,L_y,j}`$’s with coefficients $`c_{sq,L_yj}=c^{(L_y1)}`$ can be calculated as follows. Denote the equation whose solution is $`\overline{\lambda }_{sq,L_y,j}`$ as $`f(sq,L_y,\xi )`$. Thus, $`f(sq,1,\xi )=\xi +(q1)`$ and $$f(sq,2,\xi )=f(sq,1,\xi )(\xi +q3).$$ (5.2) The $`\lambda _{sq,L_y,j}`$’s for higher values of $`L_y`$ are then given by $$f(sq,L_y,\xi )=f(sq,L_y1,\xi )(\xi +q3)f(sq,L_y2,\xi )\mathrm{for}L_y3.$$ (5.3) We find that in the chromatic polynomial for the cyclic strip of the square lattice, of the $`\lambda _{sq,j}`$’s with coefficient $`c_{sq,j}=c^{(L_y1)}`$, (i) one is $`\lambda _{sq,j}=(1)^{L_y}(1q)`$; (ii) if $`L_y`$ is even, then another is $`3q`$; (iii) if $`L_y=0`$ mod 3, two others are $`(1)^{L_y}(2q)`$ and $`(1)^{L_y}(4q)`$; (iv) if $`L_y=0`$ mod 6, then two others are $`3\pm \sqrt{3}q`$. (This is not an exhaustive list of special factors.) For cyclic strips of the triangular lattice, denote the equation whose solution is $`\overline{\lambda }_{tri,L_y,j}`$ as $`f(tri,L_y,\xi )`$. Thus, $`f(tri,1,\xi )=\xi +(q1)`$ and $$f(tri,2,\xi )=f(tri,1,\xi )(\xi +q3)(\xi 1)$$ (5.4) The $`\lambda _{tri,L_y,j}`$’s for higher values of $`L_y`$ are then given by $$f(tri,L_y,\xi )=f(tri,L_y1,\xi )(\xi +q3)\xi f(tri,L_y2,\xi )\mathrm{for}L_y3$$ (5.5) The equations defining the $`\lambda _{tri,L_y,j}`$’s involve progressively higher degrees in $`\xi `$. It was shown in that the $`\lambda _{G_s,j}`$’s are the same for the cyclic and Möbius strips of a given lattice with width $`L_y`$. Therefore, for each width $`L_y`$, the $`\lambda _{G_s,j}`$’s identified above also occur in the respective Möbius strips of the square and triangular lattices, although they do not, in general, have the same coefficients $`c_{G_s,j}`$. These inferences are important because they show how one can reduce the problem of calculating the $`\lambda _{G_s,j}`$’s for larger-width strips graphs of cyclic and Möbius type from those for lower widths without recourse to the usual iterative application of the deletion-contraction or coloring matrix methods. That is, after having used these latter methods to obtain the chromatic polynomials for the first few values of $`L_y`$, the rest can be obtained purely algebraically, without further direct analysis of the graphs involved. Work on constructing the recursive formulas for the other $`\lambda _{G_s,j}`$’s is currently in progress. 3. For the infinite-length limit of a given strip graph $`G_s`$, the dominant $`\lambda _{G_s,j}`$ in region $`R_1`$ is independent of the longitudinal boundary condition, and its coefficient is $`c^{(0)}=1`$ . In particular, this $`\lambda _{G_s,j}`$ is the same for $`(FBC_y,(T)PBC_x)`$ and $`(FBC_y,FBC_x)`$ boundary conditions. When this $`\lambda _{G_s,j}`$ is the root of an algebraic equation of degree higher than linear, then, for the theorem on symmetric functions of roots of algebraic equations to apply and guarantee the polynomial nature of $`P(G_s,q)`$, it is necessary and sufficient that all of the other roots of this equation enter with the same coefficient . Hence, the analysis given in for cyclic strips of type $`G_s`$ that determines the number of $`\lambda _{G_s,j}`$’s with a specified $`c_{G_s,j}=c^{(d)}`$ places, for $`d=0`$, an upper bound on the number $`N_\lambda `$ of $`\lambda _{G_s,j}`$’s that occur in the strip of type $`G_s`$ with $`(FBC_x,FBC_y)`$ boundary conditions. In particular, for $`L_y`$ from 1 through 8, this number $`n_P(L_y,0)`$ takes the values 1,1,2,4,9,21,51,127. For the square lattice, one finds, for $`L_y`$ from 1 through the current results presented here for $`L_y=5`$, the values $`N_\lambda =1,1,2,3,7`$, as listed in Table 1. For the strips of the triangular lattice, this upper bound is realized as an equality: for $`L_y`$ from 1 to 5, the open strips have $`N_\lambda =1,1,2,4,9`$. The reason that the inequality is realized as an equality for the strips of the triangular lattice is a consequence of the different behavior of the coefficients of the square and triangular lattice strips in the Möbius case . 4. For all of the strips of the square lattice containing global circuits that we have studied, the locus $``$ encloses regions of the $`q`$ plane including certain intervals on the real axis and passes through $`q=0`$ and $`q=2`$ as well as other possible points, depending on the family. Note that the presence of global circuits is a sufficient, but not necessary, condition for $``$ to enclose regions, as was shown in (see Fig. 4 of that work). Our present results for the square lattice are in accord with, and strengthen the evidence for, the inference (conjecture) that $`\{q=0,2\}`$ $`\mathrm{for}sq(L_y,FBC_y,(T)PBC_x)L_y1`$ (5.8) $`\mathrm{and}sq(L_y,PBC_y,(T)PBC_x)L_y3.`$ (For the upper line of this equation, note that the $`L_y=1`$ graphs with $`(FBC_y,TPBC_x)`$ and $`(FBC_y,PBC_x)`$ boundary conditions are identical.) 5. The crossing of $``$ at the point $`q=2`$ for the (infinite-length limit of) strips with global circuits nicely signals the existence of a zero-temperature critical point in the Ising antiferromagnet (equivalent to the Ising ferromagnet on bipartite graphs). This has been discussed in in the context of exact solutions for finite-temperature Potts model partition functions on the $`L_y=2`$ cyclic and Möbius strips of the square lattice. In contrast, this connection is not, in general, present for strips with free longitudinal boundary conditions since $``$ does not, in general, pass through $`q=2`$. (Of the strips with $`FBC_x`$ that we have studied so far, such a crossing at $`q=2`$ was only found for the $`L_y=3`$ $`(FBC_y,FBC_x)`$ case, as one can see from Table 1.) Furthermore, for the strips without global circuits, there is no indication of any motion of the respective loci $``$ toward $`q=2`$ as $`L_y`$ increases. 6. Our exact solutions show that in the limit as $`L_x\mathrm{}`$, the respective loci $``$ for the $`W`$ functions for the infinite strips of the square lattice with the fixed values of $`L_y`$ considered and with (i) periodic or twisted periodic longitudinal boundary conditions and (ii) free longitudinal boundary conditions differ; in particular, the loci $``$ for cases with (i) pass through $`q=2`$, whereas the loci for cases with (ii) do not. This dependence of $``$ on the boundary conditions means that an $`n\mathrm{}`$ limit does not exist in a manner independent of these boundary conditions. If one fixes $`q=2`$ at the outset, i.e. considers the Ising antiferromagnet on the square-lattice strips (or if one fixes $`q`$ to the trivial value $`q=1`$ at the outset) and then calculates $`W`$, there are no pathologies; these arise when one considers nonintegral real $`q`$ in the range $`0<q<3`$. This was already discussed in the more general context of the full temperature-dependent free energy for the Potts antiferromagnet in , together with other pathologies such as a negative partition function (lack of Gibbs measure), noted earlier in , and negative specific heat. In general, the the conclusion is that a nonpathological $`n\mathrm{}`$ limit of the antiferromagnetic Potts model fails to exist at sufficiently low temperature and sufficiently small real nonintegral positive $`q`$ on strips of the square lattice. Since these strips are of fixed width, the $`L_x\mathrm{}`$ limit may be considered to be effectively quasi-one-dimensional; in contrast, a true two-dimensional thermodynamic limit would be $`L_x\mathrm{}`$, $`L_y\mathrm{}`$, with the ratio $`L_y/L_x`$ a finite nonzero constant in this limit. However, as is clear from the random cluster representation of the Potts model, the problem of a negative partition function (lack of Gibbs measure) for sufficiently small positive real nonintegral $`q`$ is present for both infinite-length, finite width strips and for the above two- or higher-dimensional infinite volume limit . Our exact results for infinite-length strips of various widths and our inference above that in the $`L_y\mathrm{}`$ limit, the loci $``$ and $`W`$ functions obtained with periodic (or twisted periodic) versus free longitudinal boundary conditions would differ is in connected with the other pathologies noted above. From the analysis in , we also conclude that a nonpathological $`L_x\mathrm{}`$ limit for the antiferromagnetic Potts model fails to exist at sufficiently low temperature and sufficiently small positive nonintegral $`q`$ on strips of the triangular lattice and a nonpathological thermodynamic limit fails to exist at sufficiently low temperature for nonintegral $`0<q<4`$ for the full triangular lattice. One could infer a generalization of this for other lattices also: a thermodynamic limit would fail to exist for the Potts antiferromagnet at sufficiently low temperature for positive nonintegral $`q`$ in the range from 0 to $`q_c`$ for the given 2D lattice, e.g., $`q_c=3`$ for the kagomé lattice. Our exact solutions are consistent with the understanding that the point $`q_c`$ for the infinite 2D (or higher-dimensional) lattice is independent of the boundary conditions used to define this infinite lattice. 7. For cyclic strips, we note a correlation between the coefficient $`c_{G_s,j}`$ of the respective dominant $`\lambda _{G_s,j}`$’s in regions that include intervals of the real axis. Before, it was shown that the $`c_{G_s,j}`$ of the dominant $`\lambda _{G_s,j}`$ in region $`R_1`$ including the real intervals $`q>q_c(\{G\})`$ and $`q<0`$ is $`c^{(0)}=1`$, where the $`c^{(d)}`$ were given in eqs. (2.8), (2.9). We observe further that the $`c_{G_s,j}`$ that multiplies the dominant $`\lambda _{G_s,j}`$ in the region containing the intervals $`0<q<2`$ is $`c^{(1)}`$. For the cyclic $`L_y=3`$ and $`L_y=4`$ strips, there is also another region containing an interval $`2qq_c`$ on the real axis, where $`q_c2.34`$ and 2.49 for $`L_y=3,4`$; in this region, we find that the $`c_{G_s,j}`$ multiplying the dominant $`\lambda _{G_s,j}`$ is $`c^{(2)}`$. 8. Our new results on cylindrical and open strips with $`L_y=5`$ confirm and extend various features that had been discussed earlier : for these values of $`L_y`$, $``$ forms arcs, and as $`L_y`$ increases, these arcs elongate and move closer together, with the arc endpoints nearest to the origin moving toward this point. This is consistent with the inference that in the $`L_y\mathrm{}`$ limit, the arcs would close to form a closed boundary that contained $`q=0`$ and $`q=q_c(sq)=3`$. One sees this general trend in the $`L_y=6`$ cylindrical strip (Fig. 3). However, in contrast with the strip graphs containing global circuits, for which the loci $``$ contained a region-enclosing boundary passing through $`q=0`$ for any $`L_y`$, this feature is evidently only approached in the limit as $`L_y\mathrm{}`$ for the strips that do not contain global circuits. The earlier calculations of cylindrical strips of the triangular lattice showed an example of a strip, namely the $`L_y=4`$ case, where $``$ contains arcs and an self-conjugate oval on the real axis , but for the cylindrical strips of the square lattice that we have investigated so far, we have not yet encountered such an oval. 9. For the $`L_x\mathrm{}`$ limit of all of the strips of the square lattice containing global circuits, a $`q_c`$ is defined, and our results for the cyclic and Möbius strips with widths from $`L_y=1`$ through $`L_y=4`$ indicate that $`q_c`$ is a non-decreasing function of $`L_y`$ in these cases. The same behavior was found for the strips of the triangular lattice with $`L_y=2`$ (and subsequently also $`L_y=3,4`$ ). This motivated the inference (conjecture) that $`q_c`$ is a non-decreasing function of $`L_y`$ for strips of regular lattices with $`(FBC_y,(T)PBC_x)`$ boundary conditions , and our present results strengthen the support for this inference. Given that, as $`L_y\mathrm{}`$, $`q_c`$ reaches a limit, which is the $`q_c`$ for the 2D lattice of the specified type (square, triangular, etc.), this inference leads to the following inequality: $$q_c(\mathrm{\Lambda },L_y\times \mathrm{},BC_y,(T)PBC_x)q_c(\mathrm{\Lambda }).$$ (5.9) Our exact solutions show that this inequality can be saturated. For example, $`q_c=3`$ for the $`L_y=3`$ torus and Klein bottle strip of the square lattice , which is equal to the $`q_c`$ value for the infinite 2D square lattice . In contrast, for (the $`L_x\mathrm{}`$ limit of) strips without global circuits, the locus $``$ does not necessarily cross the real axis, and hence there is not necessarily any $`q_c`$ defined, as was shown in . Furthermore, in these cases, even if $``$ does cross the real axis, so that a $`q_c`$ is defined, the value of $`q_c`$ is not a non-decreasing function of $`L_y`$. This is shown by our calculations of $``$ for the $`L_y=4`$, $`L_y=5`$, and $`L_y=6`$ strips of the triangular lattice with cylindrical boundary conditions in ; for these we get $`q_c=4`$ for $`L_y=4`$ but $`q_c=3.28`$ for $`L_y=5`$ and $`q_c=3.25`$ for $`L_y=6`$. Similarly, for the $`L_y=5`$ and $`L_y=6`$ cylindrical strips of the square lattice we get $`q_c=2.69`$ and $`q=2.61`$, respectively. 10. A generalized conjecture would be to consider a slab of a $`d`$-dimensional lattice $`\mathrm{\Lambda }`$ of size $`L_1\times L_2\times \mathrm{}\times L_d`$, and let $`d1`$ of the lengths of this slab go to infinity, holding one length, which can be chosen without loss of generality to be $`L_d`$, fixed and finite, and to define $`W`$ via (1.2) as $$W(\mathrm{\Lambda },L_d\times \mathrm{}^{d1},BC_1,\mathrm{},BC_d,q)=\underset{L_j\mathrm{},j=1,\mathrm{},d1}{lim}P(\mathrm{\Lambda },L_1\times \mathrm{}\times L_d,BC_1,\mathrm{},BC_d,q)^{1/n}$$ (5.10) For each of these $`W`$ functions, one would consider the corresponding continuous singular locus $``$ and its $`q_c`$, for choices of the $`BC_j`$ and $`L_d`$ where this point exists. We display the dependence of $`q_c`$ on these inputs by writing it as $`q_c(\mathrm{\Lambda }_d,L_d\times \mathrm{}^{d1},BC_1,\mathrm{},BC_d)`$. Next, we define a $`W`$ function for the $`d`$-dimensional lattice as $$W(\mathrm{\Lambda }_d,BC_1,\mathrm{},BC_d,q)=\underset{L_j\mathrm{},j=1,\mathrm{},d}{lim}P(\mathrm{\Lambda },L_1\times \mathrm{}\times L_d,BC_1,\mathrm{},BC_d,q)^{1/n}$$ (5.11) and a corresponding singular locus and $`q_c(\mathrm{\Lambda }_d)`$. As indicated in the notation, one expects that this $`q_c`$ would be independent of the $`BC_j`$, $`j=1,\mathrm{},d`$ just as is the case for the exactly known $`q_c`$ values for certain 2D lattices. Then we conjecture the inequality $$q_c(\mathrm{\Lambda }_d,L_d\times \mathrm{}^{d1},BC_1,\mathrm{},BC_d)q_c(\mathrm{\Lambda }_d)$$ (5.12) Similarly, a generalization of our inference that $`q_c`$ is a non-decreasing function of $`L_y`$ for the strips with $`(FBC_y,(T)PBC_x)`$ would be the conjecture that $`q_c(\mathrm{\Lambda }_d,L_d\times \mathrm{}^{d1},(T)PBC_1,\mathrm{},(T)PBC_{d1},FBC_d)`$ is a non-decreasing function of $`L_d`$. Our exact solutions for strips with $`(PBC_y,FBC_x)`$ boundary conditions show that if one uses periodic rather than free boundary conditions in the direction in which the slab is finite, then the resultant $`q_c`$ is not, in general, a non-decreasing function of $`L_d`$. 11. Our exact solutions for the $`L_y=4`$ cyclic and Möbius strips of the square lattice yield a singular locus $``$ that has support for $`Re(q)<0`$. In comparison (see Table 1), this was also true for the same type of strip with $`L_y=3`$, while for $`L_y=1,2`$, $``$ only had support for $`Re(q)0`$, and the only point on $``$ with $`Re(q)=0`$ was $`q=0`$ itself. This shows that for a given type of strip, increasing $`L_y`$ can shift the left-most chromatic zeros and, in the $`L_x\mathrm{}`$ limit, the left-most portion of the locus $``$,into the $`Re(q)<0`$ half-plane. The same type of behavior was found for the cyclic and Möbius strips of the triangular lattice; for $`L_y=2`$ and $`L_y=3`$, $``$ and chromatic zeros had support only for $`Re(q)0`$, while for $`L_y=4`$, this support extended into the $`Re(q)<0`$ region. In it was conjectured that global circuits were a necessary condition for lattice strips to have chromatic zeros and, in the limit $`L_x\mathrm{}`$, a locus $``$ with support for $`Re(q)<0`$. However, this conjecture was ruled out by our exact solutions for chromatic polynomials, $`W`$, and $``$ for homeomorphic expansions<sup>5</sup><sup>5</sup>5We recall two definitions from graph theory: (i) a homeomorphic expansion of a graph is obtained by inserting one or more degree-2 vertices on edge(s) of the graph; (ii) the girth of a graph is the number of edges or vertices in a minimum-distance circuit. of lattice strips with $`(FBC_y,FBC_x)`$ boundary conditions in , as also by the results for lattice strips with $`(PBC_y,FBC_x)`$ in . The homeomorphic expansions in have the effect of increasing the girth of these strip graphs, and it was found that for a given type of open strip graph, increasing the degree of homeomorphic expansion and hence the girth shifts the left-most chromatic zeros and, in the limit $`L_x\mathrm{}`$, the left-most portion of $``$, farther to the left. This is thus a different way of getting chromatic zeros and part of $``$ to have support for $`Re(q)<0`$ than in the present case of cyclic strips, where this result is obtained as a consequence of increasing the width of the strip while the girth remains constant. We remark that for all of these families of graphs, the magnitudes of the chromatic zeros and points $`q`$ on $``$ are bounded. Yet another way to get chromatic zeros and $``$ with negative real parts involves families with unbounded chromatic zeros and loci $``$ ; indeed, in we constructed families where these zeros and loci $``$ had arbitrarily large negative $`Re(q)`$ 12. There have been a number of theorems proved concerning real chromatic zeros. An elementary result is that no chromatic zeros can lie on the negative real axis $`q<0`$, since a chromatic polynomial has alternating coefficients. It has also been proved that there are no chromatic zeros in the intervals $`0<q<1`$, and $`1<q<32/27`$ . The bound of 32/27 in has been shown to be sharp; i.e., for any $`ϵ>0`$, there exists a graph with a chromatic zero at $`q=32/27+ϵ`$ . Based on our studies of strips of the square (and triangular) lattices with all of the various boundary conditions considered, we make the following observation: for such strips, we have not found any chromatic zeros, except for the zero at $`q=1`$, in the interior of the disk $`\|q1\|=1`$. This motivates the conjecture that for these strips, there are no chromatic zeros with $`\|q1\|<1`$ except for the zero at $`q=1`$. Assuming that this conjecture is valid, the bound would be a sharp bound, since the circuit graph with $`n`$ vertices, $`C_n`$, has chromatic zeros lying precisely on the circle $`\|q1\|=1`$ and at $`q=1`$ . ## 6 Values of $`W`$ for Low Integral Values of $`q`$ In previous works and sections of the present paper we have discussed values of $`W`$ for various infinite-length, finite-width lattice strips. For infinite-length limits of strips with global circuits, where the region(s) of the positive real axis in the interval $`0<q<q_c`$ are not analytically connected with the region $`R_1`$ including $`q>q_c`$ (and $`q<0`$), the ground state degeneracy per site, $`W`$, has a qualitatively different behavior than for integer or real $`qq_c`$. A comparative discussion of this was given in with the results available at that time, and it is worthwhile to use our new exact solutions to study this behavior further. In particular, it is of interest to inquire what the values of the $`W`$ functions are for the infinite-length limits of various strips of the square lattice at the points $`q=0`$, 1, and 2. Our exact analytic expressions yield the numerical values listed in Table 4. The notation follows that in Table 1. As was noted in , in general, for regions other than $`R_1`$, it is only possible to determine $`\|W\|`$ unambiguously. Hence, for uniformity, we list $`\|W\|`$ for all of the strips, including those with only a region $`R_1`$. For comparison, we also include values of $`\|W\|`$ at $`q=0,1,2`$ and 3 for infinite-length strips of the triangular lattice in Table 5. In addition, for families where, in the $`L_x\mathrm{}`$ limit, there exists a $`q_c`$, we include the respective values of $`W`$ at $`q_c`$. For the smallest widths, the $`\|W\|`$ values are relatively simple analytic expressions, e.g., for the square strips, $`\|W\|=3^{1/2}`$ for $`(FBC_y,BC_x)`$, $`L_y=2`$, $`q=0`$; $`\|W\|=13^{1/3}`$ for $`(PBC_y,BC_x)`$, $`L_y=3`$, $`q=0`$, and so forth. In the case of the triangular lattice, $`L_y=\mathrm{}`$, $`(PBC_y,FBC_x)`$, the values of $`\|W\|`$ for $`q=0`$ and $`q=4`$ are from ; the exact value for $`q=3`$ is our analytic evaluation, and the numerical values for $`q=1,2`$ are our numerical evaluations, of an integral representation in . Although we list the values in the tables only to three significant figures, we note that the $`q=1`$ value, $`\|W(tri)\|3.1716`$, is different from the $`q=0`$ value $`\|W(tri,3\times \mathrm{},PBC_y,BC_x)\|3.1748`$. For these values of $`q`$, the noncommutativity of eq. (1.3) occurs . Thus, for any connected graph $`G`$, and in particular, the lattice strips considered here, the chromatic polynomial $`P(G,q)`$ vanishes at $`q=0`$ and $`q=1`$ and hence also the function $`W_{nq}`$ defined via the order of limits on the right-hand side of eq. (1.3) vanishes. In contrast, in general, $`W_{qn}`$ defined by the limits on the left-hand side of (1.3) is nonzero. For cyclic strips of the square lattice of length $`L_x`$, at $`q=2`$, the chromatic polynomial $`P`$ is equal to 2 if $`L_x`$ is even but 0 if $`L_x`$ is odd, so that at $`q=2`$, no $`W_{nq}`$ is defined, since the limit on the right-hand side of (1.3) does not exist; however, $`W_{qn}`$ is well-defined and, in general, nonzero. Analogous comments apply for strips of the triangular lattice: at $`q=2`$, the chromatic polynomial $`P`$ vanishes identically, so $`W_{nq}=0`$, but $`W_{qn}`$ is, in general, nonzero. For cyclic strips of the triangular lattice, at $`q=3`$, then $`P=3!`$ if $`L_x=0`$ mod 3, and $`P=0`$ if $`L_x=1`$ or 2 mod 3; hence, no $`W_{nq}`$ is defined, since the limit on the right-hand side of (1.3) does not exist, but $`W_{qn}`$ is well-defined and, in general, nonzero. As with the other results given in this paper, the values of $`W`$ given in Tables 4 and 5 follow the definition $`WW_{qn}`$. Some general comments follow: 1. As is evident in Tables 4 and 5, for values of $`q`$ that are positive but sufficiently small, for a given lattice, boundary conditions, and value of $`L_y`$ studied, $`\|W\|`$ is a non-increasing function of $`q`$. In contrast, for sufficiently large $`q`$, $`\|W\|`$ increases with $`q`$. For families of graphs that involve global circuits, these two different types of behavior occur, respectively, for $`0<q<q_c`$ and $`q>q_c`$. The latter behavior is the one expected for the $`q`$-state Potts antiferromagnet, since increasing $`q`$ increases the physical ground state entropy. As examples, for the ($`L_x\mathrm{}`$ limit of) circuit graph, $`W`$ is constant for $`0q2`$, while for the cyclic or Möbius strip of the square lattice $`L_y=2`$, it decreases as $`\|W\|=\|3q\|`$ in this interval; in both of these cases, $`q_c=2`$ and $`W`$ is real and increasing for $`q>q_c`$. For the cyclic or Möbius $`L_y=3`$ and $`L_y=4`$ strips of the square lattice, $`\|W\|`$ has a different analytic form in the interval $`0q2`$ and $`2qq_c`$ but is everywhere decreasing for $`0<q<q_c`$, for the respective values of $`q_c`$. As an example of a strip with no $`q_c`$, for the open line, $`L_y=1`$, $`\|W\|`$ decreases from 1 to 0 as $`q`$ increases from 0 to 1 and increases for larger $`q`$. As another example of a strip with no $`q_c`$, for the $`L_y=3`$ strip of the square lattice with $`(PBC_y,FBC_x)`$ boundary conditions, $`\|W\|`$ decreases monotonically as $`q`$ increases from 0 and vanishes at $`q2.453`$; for larger values of $`q`$, $`W`$ is real and positive and increases with $`q`$. 2. For the $`L_x\mathrm{}`$ limit of strips with free transverse boundary conditions, $`FBC_y`$ and any longitudinal boundary conditions $`BC_x`$, it was proved that for a fixed physical $`qq_c`$, $`W`$ is a monotonically decreasing function of $`L_y`$ . However, as is evident from Tables 4 and 5 for sufficiently small positive values of $`q`$ (smaller than $`q_c`$ for strips with a $`q_c`$), $`\|W\|`$ is a non-decreasing function of $`L_y`$. 3. For the strips that we have studied whose $`L_x\mathrm{}`$ limit yields a locus $``$ with a $`q_c`$, $`\|W(q)\|`$ for fixed $`q[0,q_c]`$ is a non-decreasing function of $`L_y`$. 4. It has been shown that for physical values of $`q`$ in the $`q`$-state Potts antiferromagnet, in the $`L_x\mathrm{}`$ limit of a strip of a given type of lattice $`\mathrm{\Lambda }`$, $`W(\mathrm{\Lambda },L_y\times \mathrm{},BC_y,BC_x,q)`$ is independent of the longitudinal boundary condition $`BC_x`$ . However, for small positive values of $`q`$, $`\|W\|`$ does depend on both $`BC_y`$ and $`BC_x`$, as is evident from Tables 4 and 5. One observes that for the small integral values of $`q`$ shown in these tables, $`\|W(\mathrm{\Lambda },L_y\times \mathrm{},FBC_y,FBC_x,q)\|\|W(\mathrm{\Lambda },L_y\times \mathrm{},FBC_y,(T)PBC_x,q)\|`$ and $`\|W(\mathrm{\Lambda },L_y\times \mathrm{},PBC_y,FBC_x,q)\|\|W(\mathrm{\Lambda },L_y\times \mathrm{},PBC_y,(T)PBC_x,q)\|`$. 5. It was observed in and proved in (section 7 of) that for integer, and, by analytic continuation, real, values of $`q>max(q_c)`$ for the square and triangular lattice strips, i.e., $`q4`$, $`W(tri,q)<W(sq,q)`$. Most of the values of $`\|W\|`$ shown in Tables 4 and 5 show the opposite inequality. Together with various other properties noted above, this shows that $`\|W\|`$ behaves qualitatively differently for sufficiently small positive values of $`q`$ than for larger values. ## 7 Conclusions In conclusion, we have presented exact solutions of the zero-temperature partition function (chromatic polynomial $`P`$) and the ground state degeneracy per site $`W`$ (= exponent of the ground-state entropy) for the $`q`$-state Potts antiferromagnet on strips of the square lattice of width $`L_y`$ vertices and arbitrarily great length $`L_x`$ vertices. The specific solutions were for (a) $`L_y=4`$, $`(FBC_y,PBC_x)`$ (cyclic); (b) $`L_y=4`$, $`(FBC_y,TPBC_x)`$ (Möbius); (c) $`L_y=5,6`$, $`(PBC_y,FBC_x)`$ (cylindrical); and (d) $`L_y=5`$, $`(FBC_y,FBC_x)`$ (open), where $`FBC`$, $`PBC`$, and $`TPBC`$ denote free, periodic, and twisted periodic boundary conditions, respectively. Some inferences were given for certain terms $`\lambda _{G_s,j}`$ for cyclic and Möbius strip graphs of the square and triangular lattice that allow one to calculate them for arbitrarily wide strips (of any length). These are important because they show how one can reduce the problem of calculating the $`\lambda _{G_s,j}`$’s for these strips of arbitrarily large width from those for lower widths without recourse to the usual iterative application of the deletion-contraction or coloring matrix methods. A comparative discussion was given of the continuous nonanalytic locus $``$ for these strips and numerical results of $`W`$ were given for a range of values of $`q`$. In general, our exact solutions give further insight into the properties of the Potts antiferromagnet in the setting of infinite-length, finite width systems. Acknowledgment: The research of R. S. was supported in part by the NSF grant PHY-9722101 and at Brookhaven by the DOE contract DE-AC02-98CH10886.<sup>6</sup><sup>6</sup>6Accordingly, the U.S. government retains a non-exclusive royalty-free license to publish or reproduce the published form of this contribution or to allow others to do so for U.S. government purposes. ## 8 Appendix ### 8.1 Terms $`\lambda _j`$ for the Cyclic $`4\times m`$ Strip of the Square Lattice In this appendix we give the equations for the terms $`\lambda _j`$ for $`7j26`$. The $`\lambda _{sq4,j}`$, $`j=7,8,9`$, are roots of the cubic equation $`\xi ^3+(q^4+7q^323q^2+41q33)\xi ^2`$ (8.1.1) (8.1.2) $`+(2q^623q^5+116q^4329q^3+553q^2517q+207)\xi `$ (8.1.3) (8.1.4) $`+(q^8+16q^7112q^6+449q^51130q^4+1829q^31858q^2+1084q279)=0.`$ (8.1.5) (8.1.6) (8.1.7) The $`\lambda _j`$ for $`j=10,11,12`$ are the roots of another cubic equation $`\xi ^32(q2)(q3)\xi ^2+(q^411q^3+42q^268q+38)\xi `$ (8.1.8) (8.1.9) $`(q1)(q3)(q^3+7q^215q+11)=0.`$ (8.1.10) The $`\lambda _j`$ for $`13j16`$ are the roots of the quartic equation $`\xi ^4+(2q^312q^2+28q23)\xi ^3`$ (8.1.11) (8.1.12) $`+(q^613q^5+73q^4224q^3+396q^2381q+152)\xi ^2`$ (8.1.13) (8.1.14) $`(q1)(q3)(q^612q^5+62q^4179q^3+304q^2288q+119)\xi `$ (8.1.15) (8.1.16) $`q^9+17q^8127q^7+549q^61518q^5+2790q^43411q^3+2673q^21215q+243=0.`$ (8.1.17) (8.1.18) (8.1.19) Finally, there are two sets of roots of two degree-5 equations. The set $`\lambda _j`$ for $`17j21`$ are the roots of the equation $`\xi ^5+(q3)(2q^29q+14)\xi ^4`$ (8.1.20) (8.1.21) $`+(q^617q^5+119q^4446q^3+947q^21080q+515)\xi ^3`$ (8.1.22) (8.1.23) $`(q3)(2q^732q^6+224q^5883q^4+2106q^33028q^2+2428q835)\xi ^2`$ (8.1.24) (8.1.25) $`+(q2)(q^922q^8+213q^71193q^6+4267q^510120q^4+15922q^316013q^2+9329q2394)\xi `$ (8.1.26) (8.1.27) $`+(q1)(q^817q^7+125q^6520q^5+1342q^42206q^3+2261q^21325q+341)(q3)^2=0.`$ (8.1.28) (8.1.29) (8.1.30) The set $`\lambda _j`$ for $`22j26`$ are the roots of the equation $`\xi ^5+(4q^2+19q26)\xi ^4+(6q^458q^3+214q^2354q+219)\xi ^3`$ (8.1.31) (8.1.32) $`+(4q^6+60q^5370q^4+1198q^32144q^2+2013q773)\xi ^2`$ (8.1.33) (8.1.34) $`+(q2)(q^720q^6+162q^5693q^4+1697q^32391q^2+1805q565)\xi `$ (8.1.35) (8.1.36) $`+(q1)(q3)(q^716q^6+106q^5378q^4+788q^3967q^2+653q189)=0.`$ (8.1.37) (8.1.38) (8.1.39) ### 8.2 Generating Functions for the $`L_y=5,6`$ Strips of the Square Lattice with $`(PBC_y,FBC_x)`$ For the $`L_y=5`$ strip we calculate a generating function of the form (3.2) with $`d_𝒟=2`$, $`d_𝒩=1`$ and, in the notation of eqs. (3.4) and (3.3), we find $$b_{sq5PF,1}=q^5+10q^446q^3+124q^2198q+148$$ (8.2.1) $$b_{sq5PF,2}=q^819q^7+159q^6767q^5+2339q^44627q^3+5800q^24212q+1362$$ (8.2.2) $$A_{sq5PF,0}=q(q1)(q2)(q^712q^6+67q^5225q^4+494q^3719q^2+650q282)$$ (8.2.3) $`A_{sq5PF,1}=q(q1)(q2)(q^22q+2)(q^819q^7+159q^6767q^5`$ (8.2.4) (8.2.5) $`+2339q^44627q^3+5800q^24212q+1362)`$ (8.2.6) For the $`L_y=6`$ strip we calculate a generating function of the form (3.2) with $`d_𝒟=5`$, $`d_𝒩=4`$, with $$b_{sq6PF,1}=q^6+12q^568q^4+234q^3524q^2+727q483$$ (8.2.7) $`b_{sq6PF,2}=2q^{10}44q^9+456q^82917q^7+12710q^639322q^5+87323q^4137193q^3`$ (8.2.8) (8.2.9) $`+145624q^294100q+28114`$ (8.2.10) $`b_{sq6PF,3}=q^{14}+33q^{13}509q^{12}+4872q^{11}32374q^{10}+158152q^9586234q^8`$ (8.2.11) (8.2.12) $`+1676100q^73715937q^6+6358772q^58268225q^4+7921161q^3`$ (8.2.13) (8.2.14) $`5284418q^2+2197026q429510`$ (8.2.15) $`b_{sq6PF,4}=q^{17}+38q^{16}681q^{15}+7649q^{14}60357q^{13}+355400q^{12}1618550q^{11}`$ (8.2.16) (8.2.17) $`+5828269q^{10}16812727q^9+39098146q^873327191q^7+110295876q^6131415610q^5`$ (8.2.18) (8.2.19) $`+121386275q^483893487q^3+40850378q^212502528q+1809361`$ (8.2.20) $`b_{sq6PF,5}=q^{19}41q^{18}+794q^{17}9658q^{16}+82760q^{15}531052q^{14}`$ (8.2.21) (8.2.22) $`+2647330q^{13}10495556q^{12}+33592560q^{11}87588439q^{10}+186851845q^9326185418q^8`$ (8.2.23) (8.2.24) $`+464098186q^7533530852q^6+488389118q^5347889815q^4+185960167q^370211630q^2`$ (8.2.25) (8.2.26) $`+16703951q1884267`$ (8.2.27) $`A_{sq6PF,0}=q(q1)(q^{10}17q^9+136q^8674q^7+2296q^65640q^5+10183q^413457q^3`$ (8.2.28) (8.2.29) $`+12563q^27517q+2183)`$ (8.2.30) $`A_{sq6PF,1}=q(q1)(2q^{14}54q^{13}+695q^{12}5631q^{11}+31999q^{10}134668q^9`$ (8.2.31) (8.2.32) $`+432404q^81075802q^7+2085064q^63137110q^5+3615627q^43106751q^3+1890461q^2`$ (8.2.33) (8.2.34) $`733250q+137516)`$ (8.2.35) $`A_{sq6PF,2}=q(q1)(q^{18}38q^{17}+684q^{16}7756q^{15}+62128q^{14}373554q^{13}`$ (8.2.36) (8.2.37) $`+1748131q^{12}6513823q^{11}+19602672q^{10}48032023q^9+96128905q^8156920332q^7`$ (8.2.38) (8.2.39) $`+207640116q^6220043849q^5+183010634q^4115543495q^3+52290297q^2`$ (8.2.40) (8.2.41) $`15182726q+2135038)`$ (8.2.42) $`A_{sq6PF,3}=q(q1)(q^{21}43q^{20}+881q^{19}11444q^{18}+105796q^{17}740641q^{16}`$ (8.2.43) (8.2.44) $`+4078480q^{15}18111664q^{14}+65961019q^{13}199240735q^{12}+502713558q^{11}`$ (8.2.45) (8.2.46) $`1063474616q^{10}+1887470282q^92803761470q^8+3465937164q^73530769703q^6+2919336052q^5`$ (8.2.47) (8.2.48) $`1914246633q^4+960052617q^3346744827q^2+80479446q9033772)`$ (8.2.49) $`A_{sq6PF,4}=q(q1)(q^45q^3+10q^210q+5)(q^{19}41q^{18}+794q^{17}9658q^{16}`$ (8.2.50) (8.2.51) $`+82760q^{15}531052q^{14}+2647330q^{13}10495556q^{12}+33592560q^{11}87588439q^{10}`$ (8.2.52) (8.2.53) $`+186851845q^9326185418q^8+464098186q^7533530852q^6+488389118q^5`$ (8.2.54) (8.2.55) $`347889815q^4+185960167q^370211630q^2+16703951q1884267)`$ (8.2.56) ### 8.3 Generating Function for the $`L_y=5`$ Open Strip of the Square Lattice For this strip we calculate a generating function of the form (3.2) with $`d_𝒟=7`$ and $`d_𝒩=6`$. In the notation of eq. (3.4) we find $$b_{sq(5),1}=q^5+9q^440q^3+107q^2167q+118$$ (8.3.1) $`b_{sq5FF,2}=4q^863q^7+458q^62011q^5+5840q^411477q^3+14844q^211466q+4003`$ (8.3.2) (8.3.3) (8.3.4) $`b_{sq5FF,3}=6q^{11}+136q^{10}1432q^9+9250q^840749q^7+128594q^6`$ (8.3.5) (8.3.6) $`296624q^5+499762q^4601803q^3+492117q^2245164q+56113`$ (8.3.7) $`b_{sq5FF,4}=4q^{14}120q^{13}+1685q^{12}14681q^{11}+88695q^{10}393187q^9+1319323q^8`$ (8.3.8) (8.3.9) $`3404712q^7+6790667q^610414582q^5+12084263q^410278730q^3`$ (8.3.10) (8.3.11) $`+6051725q^22204111q+373840`$ (8.3.12) $`b_{sq5FF,5}=(q1)(q^{16}38q^{15}+674q^{14}7419q^{13}+56807q^{12}321258q^{11}`$ (8.3.13) (8.3.14) $`+1389731q^{10}4696189q^9+12540817q^826576855q^7+44582788q^658613690q^5`$ (8.3.15) (8.3.16) $`+59234653q^444497390q^3+23436736q^27733009q+1204091)`$ (8.3.17) $`b_{sq5FF,6}=(q1)^2(q^{17}38q^{16}+682q^{15}7680q^{14}+60795q^{13}359135q^{12}+1639962q^{11}`$ (8.3.18) (8.3.19) $`5915021q^{10}+17065698q^939623309q^8+74056302q^7110813572q^6`$ (8.3.20) (8.3.21) $`+131155616q^5120231650q^4+82455281q^339872376q^2`$ (8.3.22) (8.3.23) $`+12141916q1753922)`$ (8.3.24) $`b_{sq5FF,7}=(q1)^3(q2)^2(q^{15}34q^{14}+538q^{13}5259q^{12}+35541q^{11}176036q^{10}`$ (8.3.25) (8.3.26) $`+660682q^91914798q^8+4324155q^77615130q^6+10381339q^510768339q^4`$ (8.3.27) (8.3.28) $`+8235159q^34388527q^2+1459163q228580).`$ (8.3.29) Since the $`A_{sq5FF,j}`$ are rather lengthy, they are given in the copy of this paper in the cond-mat archive. With the definition $`A_{sq5FF,j}=q(q1)\overline{A}_{sq5FF,j}`$, we have $$\overline{A}_{sq5FF,0}=(D_4)^4$$ (8.3.30) where $`D_4=q^23q+3`$; $`\overline{A}_{sq5FF,1}=4q^{11}+75q^{10}653q^9+3478q^812572q^7+32346q^6`$ (8.3.31) (8.3.32) $`60381q^5+81687q^478370q^3+50664q^219788q+3517`$ (8.3.33) $`\overline{A}_{sq5FF,2}=6q^{14}154q^{13}+1854q^{12}13864q^{11}+71883q^{10}`$ (8.3.34) (8.3.35) $`273164q^9+784036q^81725384q^7+2923023q^63789945q^5+3697547q^4`$ (8.3.36) (8.3.37) $`2627998q^3+1283656q^2384667q+53170`$ (8.3.38) $`\overline{A}_{sq5FF,3}=4q^{17}+132q^{16}2056q^{15}+20067q^{14}137414q^{13}`$ (8.3.39) (8.3.40) $`+700413q^{12}2750993q^{11}+8502143q^{10}20926266q^9+41238149q^8`$ (8.3.41) (8.3.42) $`65036748q^7+81574624q^680321984q^5+60731068q^434014621q^3`$ (8.3.43) (8.3.44) $`+13280602q^23222200q+365089`$ (8.3.45) $`\overline{A}_{sq5FF,4}=(q1)^2(q^{18}40q^{17}+751q^{16}8804q^{15}+72287q^{14}441816q^{13}`$ (8.3.46) (8.3.47) $`+2084720q^{12}7769759q^{11}+23199424q^{10}55934061q^9+109187879q^8`$ (8.3.48) (8.3.49) $`172199452q^7+217795442q^6217905554q^5+168626444q^497343549q^3`$ (8.3.50) (8.3.51) $`+39444397q^210000178q+1191823)`$ (8.3.52) $`\overline{A}_{sq5FF,5}=(q1)^3(q^{19}40q^{18}+759q^{17}9082q^{16}+76836q^{15}488373q^{14}`$ (8.3.53) (8.3.54) $`+2418556q^{13}9549855q^{12}+30509903q^{11}79556288q^{10}+169998158q^9`$ (8.3.55) (8.3.56) $`297636195q^8+425129474q^7490934225q^6+451509111q^5323033449q^4`$ (8.3.57) (8.3.58) $`+173271795q^365533774q^2+15574281q1747588)`$ (8.3.59) $`\overline{A}_{sq5FF,6}=(q1)^6(q2)^2(q^{15}34q^{14}+538q^{13}5259q^{12}+35541q^{11}`$ (8.3.60) (8.3.61) $`176036q^{10}+660682q^91914798q^8+4324155q^77615130q^6+10381339q^5`$ (8.3.62) (8.3.63) $`10768339q^4+8235159q^34388527q^2+1459163q228580)`$ (8.3.64)
warning/0004/hep-ph0004103.html	ar5iv	text	# 1 STRUCTURE OF THE STANDARD MODEL ## 1 STRUCTURE OF THE STANDARD MODEL The Standard Model is our theory for the quantitative descriptions of all interactions of fundamental particles except quantum gravity effects. It is highly successful: all measurements are in agreement with the Standard Model predictions. The Standard Model is a renormalizable relativistic quantum field theory based on non-Abelian gauge symmetry of the gauge group $`SU(3)_C\times SU(2)_L\times U(1)_Y`$. It has two sectors: Quantum Chromodynamics (QCD) and the Electroweak Theory (EW). QCD is a vector gauge theory which describes the $`SU(3)_C`$ color interactions of quarks and gluons. It has rich dynamical structure such as chiral symmetry breaking, asymptotic freedom, quark confinement, topologically non-trivial configurations (monopoles, instantons). The Electroweak Theory (EW) describes the electromagnetic and weak interactions of the quarks and leptons as a chiral non-Abelian isospin and an Abelian hypercharge gauge symmetry $`SU(2)_L\times U(1)_Y`$. As a result of the Higgs mechanism, the gauge bosons $`W^\pm ,Z`$ become massive while the photon remains massless. The true dynamics behind the Higgs mechanism is not yet known. The simple one doublet Higgs sector predicts the existence of a single Higgs boson with well defined properties and its experimental search has first priority. Quarks carry both color and electroweak charges. Quarks and leptons cooperate to cancel the weak gauge anomalies. The Lagrangian of the Standard Model has important accidental global symmetries leading to baryon number and individual lepton number conservation in all orders of perturbation theory without implying absolute conservation of these quantum numbers. ### 1.1 Basic QCD In the sixties and early seventies an exciting series of beautiful experiments with many puzzling and unexpected results have lead to the discovery of QCD. #### 1.1.1 Quarks, flavor, color Spin 1/2 quarks as elementary constituents of strongly interacting hadrons have been invented by Gell-Mann and Zweig in 1964 to explain the approximate SU(3) spectral symmetry of baryons and mesons. They come in three flavor (up, down and strange) and form the fundamental triplet representation of the approximate SU(3) symmetry. The bound state wave functions of the spin 3/2 baryons decouplet composed from such objects, however, did not follow the Fermi-Dirac statistics. Color was invented by Greenberg in 1964 to restore the validity of the correct spin-statistics, requiring that every quark with a given flavor comes in three colors (red, blue, yellow). The measured normalization of the decay rate $`\mathrm{\Gamma }(\pi ^02\gamma )`$ and of the cross-section $`\sigma (e^+e^{}\mathrm{hadrons})`$ dramatically confirmed this assumption. The low lying hadron spectrum also had a more delicate chiral $`SU(3)_L\times SU(3)_R`$ symmetry that was broken spontaneously and explicitly and was described in terms of algebra of currents. The nature of the color interactions was not clear. For example, motivated by the success of current algebra, Gell-Mann suggested that the underlying field theory of strong interactions is a quark-gluon theory with one Abelian colorless gluon. Nambu instead assumed that the gluons form the octet representation of the color group $`SU(3)_C`$. These qualitative physical concepts, however, could not yet be summarized into a consistent quantitative theory. As next development, the quark constituent picture got confirmed by deep inelastic electron and neutrino experiments. The results have naturally been interpreted as the backward scattering of electrons and neutrinos on free pointlike constituents of the proton (parton model). Their quantum numbers could be extracted from the data and it turned out that the partons are quarks invented to explain hadron spectroscopy. The final theory could not be formulated since the parton model was based on the assumption of having approximately free point like constituents at short distances within the bound state wave function. This was “not consistent with the known class of renormalizable field theories” (Feynman) . #### 1.1.2 Breakthrough by ’t Hooft The breakthrough was achieved by ’t Hooft in 1971 by proving that non-Abelian gauge theories are renormalized: a new class of field theories have been discovered with strikingly new properties. The quantization of non-Abelian gauge theories is far more complex than the well-known case of Quantum Electrodynamics (QED) because of the self-interaction of the gauge field. The algebraic complexity of the Feynman rules as well as the Ward-identities of the exact gauge symmetry made the study very difficult. These theories, however, were considered by many people as irrelevant (non-physical) because the gauge bosons are necessarily massless. ’t Hooft’s proofs of the renormalizability of massless non-Abelian gauge theories and of massive gauge theories with Higgs-mechanism opened the possibility to find the fundamental theory of strong interactions as well as the electroweak interactions. #### 1.1.3 Towards QCD The discovery of the renormalizability of Yang-Mills theory by ’t Hooft helped the model builders to put together the concepts of quarks, color and flavor as the basic ingredients of a non-Abelian field theory called Quantum Chromodynamics. Gell-Mann et al. pointed out that if the previously suggested quark-gluon model of strong interactions is modified by replacing the Abelian colorless gluon with a non-Abelian colored gluon sector with exact $`SU(3)_C`$ gauge symmetry, one obtains better agreement with the experimentally established qualitative features of strong interactions. They have also speculated that in these theories quarks and gluons might be permanently confined (quark confinement), chiral symmetry breaking could take place and the so called $`U(1)_A`$ problem may be solved. In addition, a completely new fundamental property of Yang-Mills theories was discovered: at shorter and shorter distances the physics looks the same but the interaction of the particles are reduced . In contrary to the old type of field theories, Yang-Mills theories are well defined at short distances. A field theory is asymptotically free if and only if it is a non-Abelian gauge theory. In addition, the importance of asymptotic freedom in connection with Bjorken scaling was also realized and the possibility to use perturbative methods to calculate strong interaction effects has been pointed out. Gross and Wilczek using Wilson’s operator product expansion and renormalization group method have shown that the short and long distance contributions can be factorized and the short distance part can be consistently described using perturbation theory . They have derived the parton picture and interpreted Bjorken-scaling of deep inelastic scattering as leading order absorption of the virtual photon by free quarks inside the proton. As they could evaluate the corrections in first order, the predictions got spectacularly confirmed by a long experimental effort. By reformulating these results in terms of Feynman diagrams, the so called QCD improved parton model has been established. It provides us a well-defined algorithm for calculating cross-sections of hard scattering processes involving hadrons precisely. In particular, jet, W, Z and heavy quark production have been predicted and properties of bound states involving heavy quarks could be calculated. The discovery of heavy particles offered new experimental possibilities to test the QCD improved parton model predictions. #### 1.1.4 Confinement Asymptotic freedom implies that at long distances the color interactions are strong. The system condenses in some way. Quarks and gluons may get permanently confined within the hadronic bound states such that massless gauge-bosons do not appear in the particle spectrum. Wilson has pointed out that quark-confinement is a direct consequence of local gauge symmetry in the (non-physical) strong coupling limit when QCD is formulated on a four dimensional Euclidean lattice . It has been suggested that color neutralization is energetically favored in comparison with color separation. It is generally accepted by now that the mechanism of color confinement is due to the condensation of magnetic monopoles<sup>1</sup><sup>1</sup>1 Magnetic monopoles appear when the gauge is completely fixed such that the so called Gribov ambiguity is avoided. suggested by ’t Hooft . The vacuum of color dynamics is a dual superconductor where instead of condensate of Cooper-pairs one has the condensate of magnetic monopoles. The best formulation of the non-perturbative domain of gluon dynamics is the lattice gauge theory and ’t Hooft’s mechanism of quark confinement has been supported by the results of a number of numerical simulation work . #### 1.1.5 Extended objects, topologically non-trivial gauge configurations By studying extended objects we can get information on the non-perturbative aspects of the theory. They are interpreted as particles with masses inversely proportional to the coupling constant (’t Hooft-Polyakov magnetic monopoles ). Extended solutions of the classical field equations in Euclidean space (instantons) can produce tunneling effects with amplitude depending exponentially on the inverse of the coupling constant . In QCD, instantons are important for breaking the global flavor $`U(1)_A`$ symmetry and providing strong CP-violation effects ($`\theta `$-term). #### 1.1.6 The Lagrangian of QCD The transformation matrix of the fundamental representation of the local $`SU(3)`$ gauge group is $$\mathrm{\Omega }(x)_{ab}=\left(e^{iT^A\xi ^A(x)}\right)_{ab},T^A=\frac{1}{2}\lambda ^A$$ (1) where $`\lambda ^A`$ denotes the $`SU(3)`$ Gell-Mann matrices and $`\xi ^A(x)`$ is the group parameter $`A=1,\mathrm{}8`$ and $`a,b=1,2,3`$. The Dirac spinor of the quarks transforms like $$q_a^{}(x)=\mathrm{\Omega }(x)_{ab}q_b(x)$$ (2) while the gauge (gluon) fields transform inhomogenously $$T^CG^C(x)=\mathrm{\Omega }(x)T^CG^C(x)\mathrm{\Omega }^1+\frac{i}{g_s}\left(_\mu \mathrm{\Omega }(x)\right)\mathrm{\Omega }^1(x)$$ (3) After the choice of the matter field, renormalizability and gauge symmetry dictates uniquely the form of the classical Lagrangian. Introducing the non-Abelian field strength $`G_{\mu \nu }^A`$ $`=`$ $`_\mu G_\nu ^A_\nu G_\mu ^Ag_sf^{ABC}G_\mu ^BG_\nu ^C`$ (4) $`D_{ab}^\mu `$ $`=`$ $`\delta _{ab}^\mu +ig_ST_{ab}^CG^{\mu C}`$ (5) where $`f^{ABC}`$ is the structure constant of the $`SU(3)`$ Lie algebra, one gets $`_{\mathrm{classical}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}G_{\mu \nu }^AG_A^{\mu \nu }+{\displaystyle \underset{f}{}}\overline{q}^{fa}(x)i\gamma _\mu D_{ab}^\mu q(x)_b^f`$ $`{\displaystyle \underset{f}{}}m_f\overline{q}^{fa}(x)q(x)_{fa},f=u,d,\mathrm{},t`$ where the spinor label are suppressed. Mass terms are allowed since QCD is a vector theory: the color properties of the left and right handed quarks are the same. Renormalizability and gauge invariance, however, allows to add an additional term in the Lagrangian $$_\theta =\frac{\theta g_s^2}{32\pi ^2}G_{\mu \nu }^A\stackrel{~}{G}_A^{\mu \nu },\stackrel{~}{G}_A^{\mu \nu }=\frac{1}{2}ϵ^{\mu \nu \lambda \rho }G^{A,\lambda \rho }.$$ (7) This term can be written as a total derivative of a non-gauge invariant vector field $`K_\mu `$, composed from the gauge field $`G_A^\mu `$ and, therefore it can be dropped in the context of the perturbation theory. $`_\theta `$, however, can not be neglected in general. The vector field $`K_\mu `$ is not gauge invariant, it can have singular behavior at infinity and those non-trivial topological field configurations of the QCD vacuum can lead to physical CP-violating effects. But strong CP-violation is severely constrained by the data. One such effect would be the observation of electric dipole moment of the neutron. Using the experimental upper limit one gets $`\theta <10^9`$. It is puzzling why this term is so small. This question is referred to in literature as the strong CP-problem and its resolution leads to the suggestion of the existence of axions . The classical Lagrangian upon quantization gets modified: the perturbative treatment requires that gauge fixing terms, Fadeev-Popov ghost terms are added to the Lagrangian and renormalization requires counter terms $$=_{\mathrm{classical}}+_{\mathrm{gauge}\mathrm{fixing}}+_{\mathrm{ghost}}+_{\mathrm{counter}\mathrm{terms}}.$$ (8) This Lagrangian then uniquely defines the algorithm (Feynman rules) for calculating finite physical amplitudes in perturbative expansion . The quadratic terms give the propagators, the trilinear and quartic terms give the vertices. The form of the counter terms depends on the choice of regularization and renormalization scheme and are obtained by calculating a few ultraviolet divergent self-energy and vertex contribution. The algorithm is particularly simple in the case of dimensional regularization, supplemented by the mass independent $`\overline{\mathrm{MS}}`$ renormalization scheme . The Feynman rules and the one loop counter terms can be found in many textbooks . #### 1.1.7 Running coupling $`\alpha _S(\mu )`$ and the $`\mathrm{\Lambda }`$ parameter From the explicit form of the one loop counter terms one can easily derive the leading term of the beta function. It gives the measure of the change of the coupling constant with the change of the renormalization scale. In next-to-leading order of perturbation theory one obtains $`\mu ^2{\displaystyle \frac{d\alpha _S}{d\mu ^2}}`$ $`=`$ $`\beta (\alpha _S)=b_0\alpha _{S}^{}{}_{}{}^{2}b_1\alpha _{S}^{}{}_{}{}^{3}+\mathrm{}`$ (9) $`b_0`$ $`=`$ $`{\displaystyle \frac{11C_A2n_f}{12\pi }},b_1={\displaystyle \frac{17C_A^25C_An_f3C_Fn_f}{24\pi ^2}}`$ (10) where $`\alpha _S=g_s^2/(4\pi )`$, $`n_f`$ denotes the number of the quark flavors and $`C_A`$ is the color charge of the gluons $`f_{ACD}f^{BCD}=C_A\delta _{AB}`$, for $`SU(N_C)`$ gauge symmetry $`C_A=N_C`$. $`b_0`$ and $`b_1`$ are independent from regularization and renormalization schemes. Equation (9) can be easily integrated $$\alpha _S(\mu ^2)=\frac{1}{b_0\mathrm{ln}\frac{\mu ^2}{\mathrm{\Lambda }^2}}\left[1\frac{b_1}{b_0^2}\frac{\mathrm{ln}\mathrm{ln}\frac{Q^2}{\mathrm{\Lambda }^2}}{\mathrm{ln}\frac{Q^2}{\mathrm{\Lambda }^2}}\right]$$ (11) $$\alpha _S(\mu ^2)=\frac{1}{b_0\mathrm{ln}\frac{\mu ^2}{\mathrm{\Lambda }^2}}$$ where $`\mathrm{\Lambda }`$ is an integration constant. The value of the $`\mathrm{\Lambda }`$ parameter can be extracted for example from data on the scaling violation of deep inelastic scattering (see section 2.1). Its actual value $`300\mathrm{MeV}`$ gives the measure of the strength of the gluon interaction and the energy scale at which the coupling constant becomes strong. By now a large number of competing methods for extracting the value of the coupling is available. The results are conveniently normalized to the scale $`\mu =M_Z`$ and the world average is $`\alpha _S(M_Z)=0.119\pm 0.003`$. This value is relatively large, therefore, the proper scale choice of the running coupling is far more important issue in QCD than in QED. With appropriate choice of the renormalization scale one can avoid the occurrence of large logarithms in the higher order perturbative corrections. #### 1.1.8 Classical versus Quantum Symmetries, Approximate Symmetries Local consistent relativistic quantum field theories form a rather limited class. Full consistency requires renormalizability and gauge symmetry. One should note in this respect two important features. First, with requiring only local gauge symmetry, the classical Lagrangian, in many cases, possess additional global or discrete symmetries called accidental symmetries. Secondly, symmetries of the classical theory can be lost at the quantum level. We shall consider two examples. If the quarks are massless, the QCD Lagrangian, is (accidentally) scale invariant. The scale invariance of the classical theory, however, can not be maintained in the quantum theory. Quantum fluctuations of the hard modes are eliminated by the procedure of renormalization defined with the help of some scale parameter providing a hard source of scale symmetry violation. The coupling is renormalized at a given scale and it makes clear distinction between the behaviors at different mass scales. The hidden scale is introduced via the running coupling. The $`\mathrm{\Lambda }`$ parameter gives the characteristic scale of the quantum theory and by definition it is independent from the scale choice in the running coupling constant. In terms of $`\alpha _S`$ and $`\mu `$ in leading order it is given as $$\mathrm{\Lambda }_{\overline{MS}}=\mu e^{\frac{1}{2b_0\alpha _S(\mu )}}.$$ (12) QCD can be defined non-perturbatively by formulating it on a four dimensional Euclidean space-time lattice . The continuum limit is obtained in the zero lattice site and large volume limit. All lattice studies rely on the fundamental assumption that the bare coupling constant goes to zero as given by the perturbative renormalization group (see equation (11)). The validity of this assumption is by far not trivial but if we accept it then the $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ gives the physical scale also for non-perturbative quantities. In particular, if we calculate hadron masses, we have to get $$m_H=\mathrm{\Lambda }_{\mathrm{QCD}}c_H$$ (13) where $`c_H`$ is some pure number of order one. A quantum theory must have an intrinsic scale. This phenomenon is referred to in literature as dimensional transmutation. The second example is the spontaneously broken chiral symmetry and the $`U(1)_A`$ problem. For many purposes the quark masses are negligible if $`m_q<<\mathrm{\Lambda }_{\mathrm{QCD}}`$. This is fulfilled for $`m_u`$ and $`m_d`$ and with less accuracy for $`m_s`$, while it is badly broken for $`m_c,m_b`$ and $`m_t`$. Since color interactions are flavor blind rotating the quark field in the flavor space will leave the Lagrangian invariant. With three light quarks QCD has accidental global (approximate) $`U(3)_L\times U(3)_R`$ symmetry $$_{m_f=0}=\underset{f=1}{\overset{3}{}}\overline{q}^f\widehat{D}q^f=\underset{f=1}{\overset{3}{}}\left(\overline{q}_L^f\widehat{D}q_L^f+\overline{q}_R^f\widehat{D}q_R^f\right).$$ (14) The approximate $`SU(3)_L\times SU(3)_R\times U(1)_V\times U(1)_A`$ symmetry, however, has to be broken spontaneously to $`SU(3)_V\times U(1)_V`$ since the hadron spectrum has only approximate $`SU(3)_V`$ symmetry (there are no parity doublets) and baryon number conservation. The pion could successfully be interpreted as the Goldstone boson of the corresponding broken generators of the axial $`SU(2)_A`$ symmetry. The isospin singlet pseudoscalar mesons $`\eta `$ and $`\eta ^{}`$, however, are too heavy to be considered as Goldstone-bosons of the broken $`U(1)_A`$ symmetry. This difficulty is referred to in literature as the $`U(1)_A`$ problem. It has been pointed out, however, by ’t Hooft , that the $`U(1)_A`$ classical symmetry is lost at the quantum level. The conservation of the singlet axial current is formally violated by the Adler-Bell-Jackiw anomaly. It was not clear how entire units of axial U(1) charge can be created from or annihilated into the vacuum. ’t Hooft pointed out that instantons provide the necessary mechanism: they break the chiral U(1) symmetry explicitly and $`\eta `$ and $`\eta ^{}`$ get masses due to instanton contributions. All this gives a consistent picture of the strong interactions of the pseudo-scalar triplet pions (or with less accuracy of the pseudo-scalar octet mesons). Using chiral perturbation theory quark mass corrections can be taken into account and the spectroscopy of the pseudoscalar hadrons can be used to extract the light quark masses. #### 1.1.9 Heavy quark symmetries An approximate symmetry is also obtained in the infinite quark mass limit. In the bound states involving heavy quarks, the heavy quark acts only as a static source of color charge, therefore the physics does not depend on the flavor of the heavy quark and its spin orientation. In this limit the corresponding bound states must exhibit an $`SU(2N_{hq})`$ spectrum symmetry where $`N_{hq}`$ is the number of the heavy quarks . In practice this symmetry is only useful for heavy hadrons containing charm and bottom quarks. The top quark decays too fast via weak interaction to form bound states. The relations obtained in the exact heavy quark limit can be corrected systematically by calculating small perturbative symmetry breaking $`1/m_Q`$ corrections. ### 1.2 Basic Electroweak Theory #### 1.2.1 Fermi theory The theory of weak interactions started in 1933 with Fermi’s theory of beta decay. He also suggested the name neutrino for the hypothetical particle invented by Pauli in 1930 to explain the continuous energy spectrum of the electrons (or apparent non-conservation of energy). The Fermi theory of weak interactions provides an exploitation of quantum field theory outside the realm of electromagnetism by describing processes when electrons, neutrinos and atomic nuclei are created and annihilated. Fermi’s original interaction involves two vector currents in analogy with the electromagnetic interaction describing electron-electron scattering. The correct form of the interaction, however, became clear only after the discovery of parity violation in 1957 and its theoretical interpretation by Lee and Yang which lead to the proposal that the Lagrangian of weak interactions is given by the products of (V-A) currents $$_F(x)=\frac{G_F}{\sqrt{2}}\overline{p}(x)\gamma ^\alpha (g_Vg_A\gamma _5)n(x)\overline{e}(x)\gamma _\alpha (1\gamma _5)\nu (x)+\mathrm{h}.\mathrm{c}.$$ (15) where the vector coupling of the nucleon is slightly smaller than one and is given by the Cabibbo angle $`g_V=\mathrm{cos}\theta _C0.97`$. The ratio of the axial to vector couplings of the nucleon is known from the study of beta-decay with total angular momentum transitions $`\mathrm{\Delta }J=0,1`$ giving $`g_A/g_V=1.2573\pm 0.0028`$ . This Lagrangian can be used to calculate the neutron lifetime in leading order of perturbation theory in terms of the Fermi coupling $`G_F`$. From the experimental value of the neutron lifetime $`\tau =887.0\pm 2.0`$ sec one obtains a first estimate of the value of the Fermi constant $`G_F(250\mathrm{GeV})^2=1.6\times 10^5`$ $`\mathrm{GeV}^2`$. The theory is not renormalizable and the interaction is weak at low energies. With the discoveries of the pion, the muon and strange hadrons the V-A structure of weak interactions has been established in a variety of experiments. Further progress has been made with the discovery that the electron and muon number are separately conserved and that the neutrinos associated with the muons are new particles. The data have indicated that the strength and form of the four fermion interactions between fermionic doublets $`(\mathrm{p},\mathrm{n})`$, $`(\mathrm{e},\nu _\mathrm{e})`$, $`(\mu ,\nu _\mu )`$ is universal in particular the muon decay is described by the Lagrangian as $$_\mu (x)=\frac{G_F}{\sqrt{2}}\overline{\nu }_\mu (x)\gamma ^\alpha (1\gamma _5)\mu (x)\overline{e}(x)\gamma _\alpha (1\gamma _5)\nu _e(x)+\mathrm{h}.\mathrm{c}.$$ (16) This interaction allows for an important non-renormalization theorem: the photonic corrections to this transition are finite in all orders of perturbation theory . The leading corrections have been calculated 20 years ago , the $`𝒪(\frac{\alpha }{\pi })^2`$ term has been obtained by van Ritbergen and Stuart only very recently . The muon lifetime is then given by the theoretical expression $$\frac{1}{\tau _\mu }=\frac{G_F^2m_\mu ^5}{192\pi ^3}\left(1\frac{8m_e^2}{m_\mu ^2}\right)\left[1+1.810\left(\frac{\alpha }{\pi }\right)+(6.701\pm 0.002)\left(\frac{\alpha }{\pi }\right)^2+\mathrm{}\right].$$ This equation offers a convenient definition of the Fermi-coupling $`G_F`$ by assuming that the non-photonic corrections are all lumped into $`G_F`$ in a way that it can be considered a physical quantity. Using the measured value of $`\tau _\mu `$ we get $$G_F=(1.16637\pm 0.00001)\times 10^5\mathrm{GeV}^2.$$ (17) #### 1.2.2 Weak isospin and hypercharge In seeking analogy between electromagnetism and week interaction, the four-fermion interactions can be considered as the effective low energy theory of a charged massive vector boson interacting with the charged chiral current $`_I`$ $`=`$ $`{\displaystyle \frac{g}{2\sqrt{2}}}W_\alpha ^{}J^{+\alpha }+\mathrm{h}.\mathrm{c}.`$ (18) $`J_\alpha ^+`$ $`=`$ $`\left[\overline{\nu }_e(x)\gamma _\alpha (1\gamma _5)e(x)+\mathrm{}\right]`$ $`G_F`$ $`=`$ $`{\displaystyle \frac{g^2}{8M_\mathrm{W}^2}},M_\mathrm{W}110\mathrm{GeV},\mathrm{if}g<1`$ where $`W^\pm =W^1\pm iW^2`$. It is natural to consider the charged current as the charged component of the weak isospin $`SU(2)_L`$ current $`{\displaystyle \frac{1}{2}}J_\alpha ^i`$ $`=`$ $`\overline{N}_L(x)\gamma _\alpha T^iN_L(x)+\overline{L}_L(x)\gamma _\alpha T^iL_L(x)`$ $`N_L(x)`$ $`=`$ $`𝒫_L\left(\begin{array}{c}p(x)\\ n(x)\end{array}\right),L_L(x)=𝒫_L\left(\begin{array}{c}\nu (x)\\ e(x)\end{array}\right)𝒫_L={\displaystyle \frac{1}{2}}(1\gamma _5)`$ (23) where $`T^i=\tau ^i/2`$ is the $`SU(2)`$ generator in the fundamental representation. This assumption, however, implies necessarily that in addition to electromagnetism weak neutral current must exist since $$[T^+,T^{}]=2T^3Q.$$ where $`T^\pm =T^1\pm iT^2`$. Actually, one assumes that the $`SU(2)_L`$ doublets and singlets carry a diagonal hypercharge quantum numbers such that $$Q=T^3+Y$$ (24) is fulfilled. Furthermore, since hadrons are composite state of quarks, the weak hadronic currents have to be given not in terms of nuclei but quark doublets. The $`SU(2)_L\times U(1)_Y`$ quantum numbers of left and right handed quarks and leptons are listed in Table 1. The spin half matter field form three identical quark-lepton families. It is convenient to classify the matter fields in terms of left handed Weyl spinors. This is possible since under CP conjugation a right-handed spin half fermion is transformed into a left-handed antifermion. We can group the fundamental spin half particles into a reducible multiplet of doublet fermions and singlet antifermions. One quark-lepton family is composed from 15 left-handed Weyl fermions grouped in 5 irreducible components $$\psi _L^f=[Q_L^f(3,2,1/6),U_{cL}^f(3,1,2/3),D_{cL}^f(3,1,1/3),L_L^f(1,2,1/2),E_{cL}^f(1,1,1)],f=1,2,3$$ (25) where the first two numbers in the ordinary parenthesis are the dimensions of the SU(3) and SU(2) representations, respectively the third number is the value of the hypercharge and $`f`$ is the family label. The corresponding Dirac spinors will be labeled as $`\psi _{f_\chi }`$ where $`\chi `$ runs over the values $`\chi =U,D,E,N`$ and $$f_U=u,c,t,f_D=d,s,b,f_E=e,\mu ,\tau \mathrm{and}f_N=\nu _\mathrm{e},\nu _\mu ,\nu _\tau $$ (26) where $`f_\chi `$ is again the family label but for a given component of the families. #### 1.2.3 Towards Yang-Mills theories The universality of the interactions, the weak isospin structure and the analogy with QED pointed to the Yang-Mills theory with gauge group of $`SU(2)_L\times U(1)_Y`$ . The symmetric part of the Lagrangian density is given in terms of the two gauge coupling constant $`g`$ and $`g^{^{}}`$ $$_{\mathrm{ew}}=\frac{1}{4}W^{i,\mu \nu }W_{\mu \nu }^i\frac{1}{4}B^{\mu \nu }B_{\mu \nu }+2\underset{f=1}{\overset{3}{}}\overline{\psi }_L^f\gamma _\mu D^\mu \psi _L^f.$$ (27) where $`D_\mu `$ is the covariant derivative $$D^\mu =^\mu +igt^iW^{i,\mu }+ig^{}YB^\mu $$ (28) All terms containing the gluon fields are dropped and $`t^i`$ is the $`SU(2)_L`$ matrix of the reducible fermionic representation $`\psi _L^f`$. The photon field is the linear combination of $`W_3`$ and $`B`$ coupled to the electromagnetic current $$A_\mu =\mathrm{sin}\theta _WW_\mu ^3+\mathrm{cos}\theta _WB_\mu $$ (29) with $$\mathrm{tan}\theta _W=\frac{g^{}}{g},e=g\mathrm{sin}\theta _W.$$ (30) The Z-boson field is the orthogonal combination $$Z_\mu =\mathrm{cos}\theta _WW_\mu ^3+\mathrm{sin}\theta _WB_\mu $$ (31) coupled to the weak neutral current. The interaction terms of the fermions are $`_{I_f}`$ $`=`$ $`\left({\displaystyle \frac{g}{2\sqrt{2}}}J_\mu ^+W^{,\mu }+{\displaystyle \frac{g}{2\sqrt{2}}}J_\mu ^{}W^{+,\mu }+{\displaystyle \frac{g}{2\mathrm{cos}\theta _W}}J_\mu ^{\mathrm{NC}}Z_\mu +eJ_\mu ^{\mathrm{elm}}A^\mu \right)`$ with currents defined in terms of Dirac spinors $`J_\mu ^+`$ $`=`$ $`{\displaystyle \underset{f_U,f_D}{}}\overline{\psi }_{f_U}\gamma _\mu (1\gamma _5)V_{CKM}^{f_Uf_D}\psi _{f_D}+{\displaystyle \underset{f_E}{}}\overline{\psi }_{f_N}\gamma _\mu (1\gamma _5)\psi _{f_E}`$ $`J_\mu ^{\mathrm{NC}}`$ $`=`$ $`{\displaystyle \underset{f_\chi }{}}\overline{\psi }_{f_\chi }\gamma _\mu (v_\chi a_\chi \gamma _5)\psi _{f_\chi }`$ $`J_\mu ^{\mathrm{em}}`$ $`=`$ $`{\displaystyle \underset{f_\chi }{}}\overline{\psi }_{f_\chi }\gamma _\mu Q_\chi \psi _{f_\chi }`$ (32) where $`f_\chi `$ are the labels defined in equation (26), the color labels and spinor labels are suppressed and $`V_{CKM}^{ff^{}}`$ is the CKM-matrix (see subsection 1.2.5). The requirement of non-Abelian gauge symmetry leads to the universality of the gauge boson interactions and predicts the neutral current couplings $$v_\chi =T_{\chi ,L}^32Q_\chi \mathrm{sin}^2\theta _W,a=T_{\chi ,L}^3\chi =U,D,E,N.$$ (33) The chiral gauge symmetry of the Lagrangian (27) forbids mass terms both for the gauge bosons and the fermions. Adding mass terms by hand is disastrous since it destroys gauge invariance. Because of this in the early sixties these theories have not been taken seriously and the successful predictions for the neutral currents were considered to be very vague. The fact that the low energy effective theory was rather successful in explaining the charged current data implied small correction terms and gave an experimental hint that somehow massive renormalizable Yang-Mills theories must exist . The breakthrough came with the solid theoretical understanding of the renormalizability of the Yang-Mills theories and the mechanism of mass generation. #### 1.2.4 Higgs mechanism The difficulty with the mass terms and its resolution can be understood already in the case of abelian theories. Massless spin one particles have only two spin degrees of freedom, the longitudinal component does not contribute to the kinetic energy and the free theory is gauge invariant. Keeping the interactive theory gauge invariant, the longitudinal components remain decoupled and one gets renormalizable theories. Adding even an infinitesimal mass term is disastrous: the longitudinal component of the gauge bosons becomes physical and it destroys unitarity. The trouble is related to the number of degrees of freedom: the massive gauge bosons have three spin states, therefore, the massless theory can not be obtained simply as the massless limit of a massive theory. At high energies, however, the longitudinal component behaves like a scalar particle suggesting that perhaps the gauge symmetry may be maintained if we add scalar particles to the theory. The gauge transformation rules then will also involve the scalar field. This was the crucial observation of Higgs Brout and Englert leading to the discovery of the Higgs mechanism. Even if the energetically preferred value of the scalar field is not equal to zero, the Ward-Takahashi identities required by local gauge invariance can be maintained. If the ground state $`<\varphi >`$ is non-vanishing, without the requirement of local symmetry, we get spontaneous symmetry breaking with massless Goldstone bosons associated with each broken generators. In gauge theory at high energies when masses are negligible we have massless Goldstone bosons and massless gauge bosons. At low energies, however, the Goldstone bosons disappear from the theory: they provide the longitudinal component of the massive gauge bosons since the number of degrees of freedom of the theory has to be preserved. This feature of gauge theories coupled to Goldstone bosons is called the Higgs mechanism. One can obtain massive gauge bosons by supplementing the Lagrangian with some new sector providing us with the appropriate Goldstone bosons. The Standard Model is defined with the simplest realization of the Higgs mechanism : one adds to the theory one scalar doublet with appropriate hypercharge $`Y(\mathrm{\Phi })=1/2`$ $`\mathrm{\Phi }=\left(\begin{array}{c}\varphi ^+\\ \varphi ^0\end{array}\right)`$ (36) with gauge kinetic energy term and self-interactions $$_\mathrm{\Phi }(x)=\left(D_\mu \mathrm{\Phi }\right)^{}D^\mu \mathrm{\Phi }+\mu ^2\mathrm{\Phi }^{}\mathrm{\Phi }\lambda \left(\mathrm{\Phi }^{}\mathrm{\Phi }\right)^2$$ (37) and Yukawa couplings $`_{\mathrm{Yukawa}}(x)`$ $`=`$ $`{\displaystyle \underset{ff^{^{}}}{}}\lambda _{ff^{^{}}}^U\left(\overline{Q}_{Lf}\stackrel{~}{\mathrm{\Phi }}\right)u_{Rf^{^{}}}+\lambda _{ff^{^{}}}^D\left(\overline{Q}_{Lf}\mathrm{\Phi }\right)d_{Rf^{^{}}}`$ (38) $`+\lambda _{ff^{^{}}}^E\left(\overline{L}_{Lf}\mathrm{\Phi }\right)e_{Rf^{^{}}}+\mathrm{h}.\mathrm{c}.`$ where $`Q_{Lf}`$ and $`L_{Lf}`$ denote the quark and lepton doublet Weyl spinors for family $`f`$ and $`\lambda _{ff^{^{}}}^u`$, $`\lambda _{ff^{^{}}}^d`$, $`\lambda _{ff^{^{}}}^e`$ denote complex coupling matrices in the family space. There is no Yukawa coupling for neutrinos, since it is assumed that in nature only left-handed neutrinos exist. Assuming $`\mu ^2>0`$ there is a circle of degenerate minima at $$\left\|\mathrm{\Phi }\right\|^2=\frac{\mu ^2}{2\lambda }\frac{v^2}{2}.$$ (39) The excitations along the circle correspond to the Goldstone bosons. The local gauge transformations, however, also rotate $`\left\|\mathrm{\Phi }\right\|`$ along the circle. One can choose gauge condition in a way that the scalar field points (at least in leading order) to a fixed direction (unitary gauge). If we had dealt with global symmetry we would have gotten spontaneous symmetry breaking as the vacuum is not symmetric with non-vanishing scalar field. The vacuum, however, does not break the local gauge invariance. Any state in the Hilbert space that is not invariant under local gauge transformations is unphysical. With choosing unitary gauge $`\mathrm{\Phi }(x)={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}0\\ h(x)+v\end{array}\right)`$ (42) local gauge invariance is not broken but one rotates away the three unphysical components of the scalar doublet field $`\mathrm{\Phi }(x)`$. The $`h(x)`$ field describes the neutral Higgs boson remaining in the physical spectrum. Rewriting the Lagrangian in terms of $`h(x)`$ and $`v`$ we obtain mass terms for the gauge bosons $$_m^{VB}=\left\|D_\mu <\mathrm{\Phi }>\right\|^2=\frac{v^2}{4}\left[g^2W_\mu ^+W^\mu +\frac{1}{2}(W_\mu ^3,B_\mu )\left(\begin{array}{cc}g^2& gg^{^{}}\\ gg^{^{}}& g^{}_{}{}^{}2\end{array}\right)\left(\begin{array}{c}W^{3\mu }\\ B^\mu \end{array}\right)\right].$$ (43) In terms of the $`A_\mu `$ and $`Z_\mu `$ fields the mass matrix of the neutral gauge bosons becomes diagonal and one gets $$M_\mathrm{W}=\frac{1}{2}gv,M_\mathrm{Z}=\frac{1}{2}\sqrt{g^2+g^{}_{}{}^{}2},m_\gamma =0.$$ (44) For the mass of the Higgs boson we obtain $$M_\mathrm{H}^2=2\lambda v^2$$ (45) therefore the strength of the self-interaction of the Higgs boson can be expressed in terms of the Higgs and gauge boson masses and the gauge coupling $$\lambda =\frac{M_\mathrm{H}^2}{8M_\mathrm{W}^2}g^2.$$ (46) The gauge symmetry uniquely defines the coupling of the gauge bosons to the Higgs boson allowing to predict for example the value of the half-width of the Higgs boson $$\mathrm{\Gamma }(h\mathrm{W}^+\mathrm{W}^{})=\frac{g^2M_\mathrm{H}^3}{64\pi M_\mathrm{W}^2}\sqrt{14x_h}(14x_h+12x_h^2),x_h=\frac{M_\mathrm{W}^2}{M_\mathrm{H}^2}.$$ (47) With increasing $`M_\mathrm{H}`$ it grows as $`M_\mathrm{H}^3`$. In particular for $`M_H1\mathrm{TeV}`$ we get $`\mathrm{\Gamma }(h)M_\mathrm{H}`$ indicating the difficulty with the validity of perturbative unitarity in case of a heavy Higgs boson. With substituting the shifted field (42) in the Yukawa coupling of $`\mathrm{\Phi }`$ to fermions we get the mass matrices of the fermions and their couplings to the Higgs boson. The physical fermion states are obtained by diagonalizing the mass matrices (with biunitary rotations) $$\frac{v}{\sqrt{2}}U(\chi )_L^{}\lambda ^\chi U(\chi )_R,=_{\mathrm{diag}}^\chi ,\chi =U,D,E.$$ (48) $`\chi =N`$ does not occur since it is assumed that right handed neturinos do not exist. The diagonal element of $`_{\mathrm{diag},\mathrm{ff}^{}}^\chi =m^\chi \delta _{ff^{^{}}}`$ gives the mass values and $`f`$ runs over the three families. This diagonalization produces three important physical results. First, the couplings of the Higgs boson to fermions are flavor diagonal and proportional to the fermion mass $$_h^Y(x)=\underset{\chi ,f}{}\frac{gm_{f_\chi }}{2M_\mathrm{W}}\overline{\psi }_{f_\chi }(x)\psi _{f_\chi }(x)h(x),$$ (49) therefore, the coupling of the Higgs bosons to light fermions is very weak. This makes its experimental search extremely difficult. Secondly, the charged current of quarks is not flavor diagonal $$J_\mu ^{}(x)=\underset{f_U,f_D}{}\overline{\psi }f_U(x)V_{\mathrm{CKM}}^{f_Uf_D}\gamma _\mu (1\gamma _5)\psi _{f_D}(x),f_U=u,c,t,f_D=d,s,b$$ (50) where $`V_{\mathrm{CKM}}=U^{}(U)_LU(D)_L`$ denotes the Cabibbo-Kobayashi-Maskawa matrix. In the charged current of quarks six fields with 5 physically irrelevant independent phases are involved (with one relevant phase for $`U(1)_Y`$), therefore four phases of unitary CKM matrix can be rotated away and we end up with $`3^25=4`$ physically relevant parameters. If the neutrinos are massless, one can arbitrarily rotate the neutrinos so that the charged lepton current remains diagonal. Recently, experimental evidence has been obtained for neutrino oscillations and so for neutrino masses . The pattern for massive neutrinos is more complicated than those for massive quarks as the electrically neutral neutrinos can have Dirac and/or Majorana mass terms. In the minimal Standard Model with one Higgs doublet, however, only the Dirac mass term is possible. Thirdly, after the Higgs mechanism the neutral current remains obviously flavor diagonal (GIM mechanism ), therefore in leading order there are no flavor changing neutral current transitions. Before the discovery of the charm quark this feature was less obvious. In higher order, as a result of virtual flavor changing charged current exchanges, flavor changing neutral current transitions are allowed but suppressed strongly by the smallness of the Fermi coupling. #### 1.2.5 CKM matrix The elements of the CKM matrix are denoted conveniently as $`V_{\mathrm{CKM}}^{ff^{}}`$ with $`f=u,c,t`$ and $`f^{^{}}=d,s,b`$. As we noted above, they can be described in terms of four independent parameters. The values of these have to be extracted from the data. Although the independent parameters of the CKM matrix are fundamental parameters of the theory, there is considerable freedom in their definite choice. According to the data the matrix elements are large in the diagonal and they get smaller and smaller as we move away from it. Wolfenstein has suggested a convenient parametrization which takes this hierarchical behavior into account $`V^{\mathrm{CKM}}=\left(\begin{array}{ccc}1\lambda ^2/2& \lambda & A\lambda ^3(\rho i\eta )\\ \lambda & 1\lambda ^2/2& A\lambda ^2\\ A\lambda ^3(1\rho i\eta )& A\lambda ^2& 1\end{array}\right)+𝒪(\lambda ^4).`$ (54) Experimentally $`\lambda 0.22`$, $`A0.82`$ $`\sqrt{\rho ^2+\eta ^2}0.4`$ and $`\eta 0.3`$ . Since $`\eta `$ is non-vanishing CP-violation occurs both in the neutral kaon mass matrix and in the transition amplitudes, therefore, the CKM model of CP-violation is milliweak. Recent experiments also confirmed the presence of CP-violation in the decay amplitudes of kaons into pions in agreement with the prediction of the Standard Model. Unfortunately, the theoretical predictions have large theoretical error due to non-perturbative QCD effects , therefore the agreement gives only a qualitative confirmation. The CKM predictions of CP-violating effects in terms of a single CP-violating parameter depend sensitively on the assumptions that the Higgs sector is minimal and that we have only three fermion families. If the Higgs sector is more complicated or there are additional heavy fermions the predictions of the Standard Model for CP violation effects will fail. This is why the precision test of the CKM predictions for CP-violation is so important. The CKM matrix is unitary in the Standard Model but is not necessarily unitary in its extensions. It is of great interest to test the relations among the CKM matrix elements required by unitarity. A particularly interesting relation is the prediction $$V_{td}V_{ud}^{}+V_{ts}V_{us}^{}+V_{tb}V_{ub}^{}=0.$$ (55) All the three terms at the left-hand side are proportional to $`A\lambda ^3`$. Up to this overall factor the three terms are equal to $`1\rho +i\eta `$$`\rho i\eta `$$`1`$ and they form a triangle (Bjorken) in the complex plane. The surface of the triangle is equal to the Jarlskog invariant $$JA\lambda ^3\eta .$$ (56) If the CKM matrix conserves CP, the unitarity triangle shrinks to a line. The observation of CP-violation in b-decays in the near future will provide a decisive test on the validity of the CKM model of CP-violation and the fermion-mass generation mechanism of the Standard Model. #### 1.2.6 Custodial symmetry We have already noted in section 1.18 that in the case of a simple matter field content the requirement of local gauge symmetry and renormalizability may lead to a Lagrangian density with accidental global symmetries. The Higgs sector of the Standard Model has an accidentally global $`SU(2)_L\times SU(2)_RSO(4)`$ symmetry. The non-vanishing vacuum expectation value of the Higgs field breaks it down spontaneously to the diagonal $`SU(2)_V`$ global symmetry (called custodial symmetry). In the limit $`g^{^{}}=0`$ the gauge interaction preserves this symmetry and in this limit the massive gauge bosons must form a degenerate triplet representation of $`SU(2)_V`$. Non-vanishing $`g^{^{}}`$ coupling leads to the mass splitting $`M_\mathrm{W}=\mathrm{cos}\theta _WM_\mathrm{Z}`$ (see equation (43)). This relation therefore remains valid for any Higgs mechanism which respects custodial symmetry. The Yukawa couplings also violate the custodial symmetry if the mass values of the up and down components of a fermion doublets are not degenerate. Since the top-bottom mass splitting is large, virtual top and bottom quark contributions give large corrections to the leading order value of the $`\rho `$ parameter $`\rho =M_\mathrm{W}/(\mathrm{cos}\theta _WM_\mathrm{Z}).`$ #### 1.2.7 Cancellation of chiral gauge anomalies Classical chiral gauge symmetries may be broken in the quantum theory by triangle anomalies. Chiral fermions are massless but renormalization requires the regularization of the theory which necessarily introduces masses for the fermions. It may happen that the chiral symmetry of the classical theory will not survive in the quantum theory. This is disastrous for the chiral gauge theories because the gauge symmetry is broken and the theory makes no sense. Fortunately, all terms which break chiral symmetry have very simple origins as they come from simple fermionic triangle diagrams coupled to vector and axial vector currents. Therefore, the anomaly is proportional to the difference of the trace of coupling matrices of the left-handed and right-handed fermions (masses are negligible in the ultraviolet limit). A chiral gauge theory is only meaningful if the chiral anomalies cancel each other. One can easily see that the condition of anomaly cancellation in the Standard Model is that the sum of the charges of the fermions vanishes. In the Standard Model this is fulfilled individually for each family $$TrQ=3(Q_u+Q_d)+Q_e=0.$$ (57) This simple result follows because the group $`SU(2)`$ is anomaly free. The condition (57) forms an important bridge between the electroweak sector and the strong sector: without quarks the lepton sector is anomalous and the quarks must come in three colors. We note the problem of charge quantization. It is a phenomenological fact that the charges of the proton and the positron are equal to each other within very high experimental accuracy. The relation $`Q=T_3+Y`$, however, would allow arbitrary relation since the $`U(1)_Y`$ charge is not quantized. By choosing the value of $`Y`$ consistently with value of the charges of the positron and the proton we get the anomalies cancelled. This indicates that charge quantization and anomaly cancellation may be connected. In the case of global symmetries the anomalies do not destroy renormalizability, but the quantum theory will not be symmetric. For example the strong chiral-isospin symmetry forbids the decay $`\pi ^0\gamma \gamma `$, but in the quantum theory this symmetry is violated by the triangle anomaly and the decay is allowed. The anomaly is proportional to $`Tr(Q^2T_3)=N_c(Q_u^2Q_d^2)`$. This result played a crucial role in the discovery of color. #### 1.2.8 Accidental continuous global symmetries It is a success of the Standard Model that the requirement of local gauge invariance and renormalizability leads to accidental global symmetries. The origin of these symmetries is the large $`U(45)`$ symmetry of the Lagrangian of the 45 Weyl fermions of the Standard Model. The gauge interaction breaks down this symmetry to three copies of 5 irreducible $`SU(3)\times SU(2)\times U(1)`$ representations (see equation (25)) which still has $`U(3)^5\times U(1)`$ global symmetry. This is broken by the Yukawa coupling to $`U(1)^4`$. Because of the CKM matrix in the quark sector, only one phase rotation survives: each quark is rotated with the same phase and each antiquark field is rotated with the opposite phase. Since we do not have a CKM matrix in the lepton sector we can have individual phase rotations of leptons for each family. These symmetries lead to the $`\mathrm{𝑒𝑥𝑎𝑐𝑡}`$ baryon number $`B`$ and to individual lepton number conservations for each family $`L_e`$, $`L_\mu `$ and $`L_\tau `$. According to the data these conservation laws are valid to very high precision . I have to note that we do not consider these symmetries to be absolute. They are violated even within the Standard Model by instanton contributions . The baryon current coupled to two gauge currents via the anomalous triangle diagrams is necessarily anomalous. The condition (57) requires that $`BL`$ is conserved. The instantons can absorb baryon and lepton numbers. Also, considering the Standard Model as an effective low energy field theory with all possible higher dimensional non-renormalizable gauge invariant operators baryon and lepton number can be violated. The baryon and lepton number violation is strongly suppressed by power corrections. The search for such effects is an important tool for testing the range of validity of the Standard Model. ## 2 PRECISION CALCULATIONS ### 2.1 Testing QCD QCD is asymptotically free, therefore the physical phenomena at short distances and at finite time intervals may be in principle subject to perturbative treatment in terms of weakly interacting quarks and gluons. It is crucially important for collider physics that the perturbative description is valid for large momentum transfer reactions since perturbation theory is the only systematic method for calculating scattering cross sections directly from the QCD Lagrangians. The application of perturbative QCD to scattering phenomena with large momentum transfer, however, is not straightforward. It is not obvious a priori that the short and long distance properties can be meaningfully separated. In general the cross sections of scattering processes in perturbative QCD with massless quarks and gluons are singular due to the presence of soft and collinear contributions. Fortunately, in simple cases when only one hard scale is involved like the total cross section of $`e^+e^{}`$ annihilation to hadrons, deep inelastic scattering and the Drell-Yan processes, one can prove in all order of the perturbation theory that all the infrared sensitive contributions given by soft and collinear parton configurations are cancelled except some remaining collinear singularities . They are, however, universal and can be factored into the parton distribution functions of the incoming hadrons or into the fragmentation functions of final hadrons. The fundamental assumption of the QCD improved parton model is that this theorem remains valid after including non-perturbative effects up to some power corrections of $`𝒪(\mathrm{\Lambda }_{QCD}/Q)`$ where $`Q`$ denotes the hard scale of the process. Furthermore, it is assumed that the theorem remains valid for any infrared safe quantities. A physical observable is called infrared safe if its value calculated in the perturbation theory is not sensitive to the emission of additional soft gluons or the splitting of a hard parton into two collinear partons. In the QCD improved parton model the incoming hadrons are considered as wide band beams of hadrons with well defined momentum distributions. Infrared safe hard scattering cross sections of hadrons are calculated in the perturbation theory in terms of partons. #### 2.1.1 Hadron production in $`e^+e^{}`$ annihilation In simple inclusive reactions, such as the total cross section of $`e^+e^{}`$ annihilation into quarks and gluons, the soft and collinear contributions cancel (KLN theorem). Therefore, the cross section is free from infrared singularities and can be calculated in power series of the effective coupling $$R=\frac{\sigma (e^+e^{}\mathrm{hadrons})}{\sigma (e^+e^{}\mu ^+\mu ^{})}=\stackrel{\overline{R}}{\stackrel{}{(1+\frac{\alpha _s}{\pi }+\mathrm{})}}3\underset{2}{}e_q^2$$ (58) where $$\overline{R}=1+\frac{\alpha _s(\mu )}{\pi }+\left(\frac{\alpha _s(\mu )}{\pi }\right)^2\left[\pi b_0\mathrm{ln}\frac{\mu ^2}{s}+B_2\right]+\mathrm{}$$ (59) $`\alpha _S`$ is the running coupling constant, $`\mu `$ is the renormalization scale, $`B_2`$ is a known constant given by the NNLO calculation and $`b_0`$ is the first coefficient in the beta function given in equation (9). The explicit $`\mu `$ dependence in (59) is cancelled by the $`\mu `$ dependence of the running coupling constant to $`𝒪(\alpha _{S}^{}{}_{}{}^{3})`$. In general, the truncated series is $`\mu `$-dependent but the $`\mu `$ dependence is order of $`𝒪(\alpha _{S}^{}{}_{}{}^{(n+1)})`$ if the cross section is calculated to $`𝒪(\alpha _{S}^{}{}_{}{}^{n})`$. The result obtained for partons can be applied to hadrons assuming that the unitarity sum over all hadronic final state can be replaced with the unitarity sum over all final state quark and gluons $$\underset{h}{}\|h><h\|=\underset{q,g}{}\|\mathrm{gluons},\mathrm{quarks}><\mathrm{gluons},\mathrm{quarks}\|.$$ (60) This assumption is only valid if the annihilation energy is much larger than the quark and hadron mass values and if the annihilation energy $`Q`$ is in a region far from resonances and thresholds (or we smear over the threshold and resonance regions). The KLN theorem remains valid also for integrating over final states in a limited phase space region, as in the case of jet production. The Sterman-Weinberg two-jet cross section is defined by requiring that all the final state partons are within a back-to-back cone of size $`\delta `$, provided their energy is less than $`ϵ\sqrt{s}`$. At NLO $`\sigma _{2\mathrm{j}\mathrm{e}\mathrm{t}}`$ $`=`$ $`\sigma _{\mathrm{SW}}(s,\epsilon ,\delta )`$ (61) $`=`$ $`\sigma _{\mathrm{tot}}\sigma _{q\overline{q}g}^{(1)}(\mathrm{all}E>\epsilon \sqrt{s},\mathrm{all}\theta _{ij}>\delta )`$ $`=`$ $`\sigma _0\left[1{\displaystyle \frac{4\alpha s}{3\pi }}\left(4\mathrm{ln}2\epsilon \mathrm{ln}\delta +3\mathrm{ln}\delta 5/2+\pi ^2/3\right)\right]`$ where $`\sigma _0=4\pi \alpha ^2/3s`$. We can easily see that the jet definition is infrared safe and therefore the cancellation theorem remains valid. The dependence of the cross section on the jet defining parameters $`ϵ,\delta `$ is physical since the same parameters have to be used in the measurements of jet cross sections when jets are defined in terms of hadrons. #### 2.1.2 Hard scattering with hadrons in the initial state In infrared safe quantities for processes with partons in the initial state, the initial state collinear singularities are not cancelled but they are universal ( process independent) in all orders in the perturbation theory . Therefore, they can be removed by collinear counter terms generated by the ‘renormalization’of the incoming parton densities. The choice of the finite part of the collinear counter terms is arbitrary and allows the definition of the different factorization schemes. Of course the physics under a change of the factorization scheme remains the same since the change in the parton cross-sections is compensated with the change in the parton densities. The collinear subtraction terms define the kernels of the scale evolution of the parton number densities (Altarelli-Parisi equation ). In the parton model, the differential cross section for hadron collisions has the form $$d\sigma _{AB}(p_A,p_B)=\underset{ab}{}𝑑x_1𝑑x_2f_{a/A}(x_A,\mu )f_{b/B}(x_B,\mu )𝑑\widehat{\sigma }_{ab}(x_Ap_A,x_Bp_B,\mu ),$$ (62) where $`A`$ and $`B`$ are the incoming hadrons, $`p_A`$ and $`p_B`$ their momentum, and $`a,b`$ run over all the parton flavors which can contribute. $`d\widehat{\sigma }_{ab}`$ denotes the finite partonic cross section, in which the singularities due to collinear emission of incoming partons are subtracted and the scale $`\mu `$ is the factorization scale. Equation (62) applies also when the incoming hadrons are formally substituted for partons. In NLO their densities are singular $$f_{a/d}(x)=\delta _{ad}\delta (1x)\frac{\alpha _S}{2\pi }\left(\frac{1}{\overline{ϵ}}P_{a/d}(x,0)K_{a/d}(x)\right)+𝒪\left(\alpha _{S}^{}{}_{}{}^{2}\right),$$ (63) where $`P_{a/d}(x,0)`$ are the Altarelli-Parisi kernels in four dimensions ( we use $`42ϵ`$ dimensions and the $`0`$ in the argument of $`P_{a/d}`$ stands for $`ϵ=0`$). The finite functions $`K_{a/d}`$ are arbitrary, the $`\overline{\mathrm{MS}}`$ subtraction scheme is obtained by choosing $`K_{a/d}0`$. Expanding the unsubtracted and subtracted partonic cross sections to next-to-leading order $$d\sigma _{ab}=d\sigma _{ab}^{(0)}+d\sigma _{ab}^{(1)},d\widehat{\sigma }_{ab}=d\widehat{\sigma }_{ab}^{(0)}+d\widehat{\sigma }_{ab}^{(1)},$$ (64) we obtain $`d\widehat{\sigma }_{ab}^{(0)}(p_1,p_2)`$ $`=`$ $`d\sigma _{ab}^{(0)}(p_1,p_2)`$ (65) $`d\widehat{\sigma }_{ab}^{(1)}(p_1,p_2)`$ $`=`$ $`d\sigma _{ab}^{(1)}(p_1,p_2)+d\sigma _{ab}^{\mathrm{count}}(p_1,p_2)`$ (66) where $`d\sigma _{ab}^{\mathrm{count}}(p_1,p_2)`$ $`=`$ $`{\displaystyle \frac{\alpha _S}{2\pi }}{\displaystyle \underset{d}{}}{\displaystyle 𝑑x\left(\frac{1}{\overline{ϵ}}P_{d/a}(x,0)K_{d/a}(x)\right)𝑑\sigma _{db}^{(0)}(xp_1,p_2)}`$ $`+{\displaystyle \frac{\alpha _S}{2\pi }}{\displaystyle \underset{d}{}}{\displaystyle 𝑑x\left(\frac{1}{\overline{ϵ}}P_{d/b}(x,0)K_{d/b}(x)\right)𝑑\sigma _{ad}^{(0)}(p_1,xp_2)}.`$ Equation (LABEL:counterterms2) defines the collinear counter terms for any finite hard scattering cross section. The parton number densities at a given scale have to be extracted from the data. One can systematically improve the accuracy of the predictions by calculating also higher order corrections. The accuracy of the recent experimental results requires the inclusion of higher order radiative corrections for a large number of measured quantities . We have a large number of quantitative well tested prediction for deep inelastic scattering, $`W/Z`$ production, jet-production and heavy quark production. The method of calculation follows of the procedure described above. Both the theoretical and experimental error could be significantly reduced during the last decade and bt now they are about $`15\%30\%`$. As illustration of this in Fig. 1 the inclusive transverse energy $`E_T`$ distribution for jet production at Tevatron is shown. The data points obtained by CDF are compared by the absolute NLO QCD prediction . These data are interesting since i) they test the absolute prediction of QCD up to NLO accuracy ii) they give information on the parton number densities at large $`Q^2`$ and iii) they can be used to constrain possible new physics. #### 2.1.3 Measurements of $`\alpha _S`$ QCD has been discovered by its qualitative properties. Over the last 20 years, however, significant progress has been achieved in the field theoretical techniques for deriving its consequences and by now we can test QCD $`\mathrm{𝑞𝑢𝑎𝑛𝑡𝑖𝑡𝑎𝑡𝑖𝑣𝑒𝑙𝑦}`$. The free parameters of QCD are the quark masses and the strong coupling constant. Once we fitted these values to some measurable infrared safe quantities, hard scattering processes can be predicted from first principles for $`e^+e^{}`$ annihilation. For processes with hadrons in the initial state the parton number densities should also be fitted. We do not want to go into this detail (see ref. ). The status of QCD tests is well characterized with the agreement between the various values of $`\alpha _S`$ extracted from different measurements. We illustrate the present situation with two figures . Fig. 3 gives a test for the running of $`\alpha _S(Q)`$ with plotting the value of $`\alpha _S(Q)`$ as obtained by various experiments. The data confirms the energy dependence of $`\alpha _S(Q)`$ as predicted by QCD. In Fig. 3 we show the summary of the results of the same experiments but evolving the values of $`\alpha _S(Q)`$ to a common energy scale, $`Q=M_\mathrm{Z}`$ using the QCD $`\beta `$ function in $`𝒪(\alpha _{S}^{}{}_{}{}^{4})`$ with 3-loop matching at the heavy quark pole masses $`M_b=4.7\mathrm{GeV}`$ and $`M_c=1.5\mathrm{GeV}`$. The corresponding world average is also given as well as the values of $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ for $`n_f=4,5`$. ### 2.2 Testing the electroweak theory In the Standard Model at tree level the gauge bosons $`\gamma ,W,Z`$ and their interactions are described in terms of three parameters: the two gauge coupling constants $`g,g^{}`$ and the vacuum expectation value of the Higgs field $`v`$. We need to know their values as precisely as possible. They have to be fitted to the three best measured physical quantities of smallest experimental error: $`G_\mu ,M_\mathrm{Z}`$ and $`\alpha `$. In leading order we have the simple relations $$G_\mu =\frac{1}{\sqrt{2}v^2},M_\mathrm{Z}=\frac{1}{2\mathrm{cos}\theta _W}gv,\alpha =\frac{g^2}{4\pi }\mathrm{sin}^2\theta _W,\mathrm{tan}\theta _W=\frac{g^2}{g^2}.$$ (68) The muon coupling $`G_\mu `$ is extracted from the precise measurement of the muon lifetime using the theoretical expression given by equation (1.2.1) $$G_F=(1.16637\pm 0.00001)\times 10^5\mathrm{GeV}^2.$$ (69) The value of $`M_\mathrm{Z}`$ is extracted from the line shape measurement at the $`Z`$-pole. There are subtleties in the theoretical definition of the mass and the width at higher order associated with the truncation of the perturbative series and gauge invariance. The latest best value is $$M_\mathrm{Z}=(91.1871\pm 0.0021)\mathrm{GeV}.$$ (70) The best value of $`\alpha `$ is extracted from the precise measurement of the electron anomalous magnetic moment $`(g_\mathrm{e}2)`$ $$1/\alpha =137.03599959\pm 0.00000038.$$ (71) Additional physical quantities like the mass of the W-boson $`M_\mathrm{W}`$, the lepton asymmetries at the Z-pole, the leptonic width of the Z-boson $`\mathrm{\Gamma }_\mathrm{l}`$ etc. are derived quantities. At the level of the per mil accuracy the predictions obtained in Born approximations for derived quantities, however, disagree with the measured values significantly. #### 2.2.1 Quantum corrections At LEP, SLC and Tevatron an enormous amount of data has been collected on the Z and W bosons and their interactions . This allows for an unprecedented precision test of the Standard Model at the level of the per mil accuracy. At this precision one and two-loop quantum fluctuations give measurable contributions and data data show sensitivity also to the Higgs mass and the top mass. Since the Standard Model is a renormalizable quantum field theory the theoretical predictions of the theory can be improved systematically by calculating higher order corrections. In particular, the recent precision of the data requires the study of the complete next-to-leading order corrections, resummation of large logarithmic contributions and a number of two loop corrections. At higher order the derived quantities show sensitivity also to the values of the mass parameters $`m_t`$, $`M_\mathrm{H}`$, $`m_b`$ and the QCD coupling constant $`\alpha _s`$. From direct measurements one obtains $`\alpha _s=0.119\pm 0.002`$ (see Fig. 3), $`m_t=173.8\pm 5.0\mathrm{GeV}`$ and $`\overline{m}_b(\overline{m}_b)=4.25\pm 0.08\mathrm{GeV}`$, $`M_\mathrm{H}102\mathrm{G}\mathrm{e}\mathrm{V}`$ where $`\overline{m}_b(\overline{m}_b)`$ denotes the $`\overline{\mathrm{MS}}`$ mass . The error bars give parametric uncertainties in the predictions and limit our ability to extract a precise value of the Higgs mass. The calculation of the higher orders requires a choice of the renormalization scheme.<sup>2</sup><sup>2</sup>2 For a complete discussion of this technical detail with references see . In the perturbation theory the higher order effects can be given in bare and renormalized parameters. Let us consider the basic observables as the basic bare parameters $`a_0^i(G_0,\alpha _0,M_{Z0})`$. Calculating the radiative corrections in the regularized theory the radiatively corrected renormalized values of the basic parameters can be written as $$a^i(a_0^i)=a_0^i+\delta a^i(a_0^i).$$ (72) This relations can be inverted $$a_0^ia_0^i(a^i).$$ (73) Similarly for derived observables we can write $$O(a_0^i)O_0(a_0^i)+\delta O(a_0^i).$$ (74) At one loop we can express the radiatively corrected derived quantities as functions of the radiatively corrected basic quantities $`O(a_0^i(a_i))`$ $``$ $`O_0(a_0^i)+\delta O(a_0^i)`$ (75) $``$ $`O_0(a^i)+\delta ^{(1)}O(a^i){\displaystyle \underset{i}{}}{\displaystyle \frac{O_0}{a^i}}\delta ^{(1)}a^i`$ $`=`$ $`O_0(a_i)+\mathrm{\Delta }^{(1)}(a^i).`$ The coefficients of this expansion are finite since we expand renormalized finite quantities in terms of the renormalized finite basic observables. Equation (75), however, tells us that the finite one-loop corrections have a direct part $`\delta ^{(1)}O`$ and an indirect one $`_i\frac{O_0}{a^i}\delta ^{(1)}a^i`$ coming from corrections to the basic parameters. The two contributions can be separately divergent, only the sum is finite. In the on shell scheme the mixing angle $`s^2=\mathrm{sin}^2\theta `$ is defined by the tree level relation $`s^2=1M_\mathrm{W}^2/M_\mathrm{Z}^2`$ in all order, therefore, it is a physical quantity. It is customary to define auxiliary dimensionless parameters $`r_W`$ by the relation $$s^2c^2\frac{\pi \alpha }{\sqrt{2}G_\mu M_\mathrm{Z}^2(1r_W)}$$ (76) where $`c^2=1s^2`$. Obviously $`r_W`$ is also a finite derived quantity and it gives the radiative corrections to the $`M_\mathrm{W}`$. In the $`\overline{\mathrm{MS}}`$-scheme the measured values of $`\alpha `$, $`G_\mu `$, $`M_\mathrm{Z}`$, $`m_f`$, $`\alpha _S`$ are used to fix the input parameters of the theory with $`M_\mathrm{H}`$ as free parameter. The $`\overline{\mathrm{MS}}`$ gauge couplings evaluated at the scale of $`M_\mathrm{Z}`$ are denoted as $`\widehat{e}`$ and $`\widehat{s}^2=\mathrm{sin}^2\widehat{\theta }_W(M_\mathrm{Z})`$. The renormalized parameters $`\widehat{s}^2`$, $`\widehat{e}^2`$ can be completely calculated in terms of $`G_\mu `$, $`\alpha `$ and $`M_\mathrm{Z}`$. Another useful auxiliary quantity is the effective mixing angle $$\mathrm{sin}^2\theta _W^{eff}=\frac{1}{4}\left(1\frac{\overline{v}_l}{\overline{a}_l}\right)=s^2(1+\mathrm{\Delta }k^{^{}}).$$ (77) This is uniquely related to the ratio of the effective neutral current vector and axialvector couplings $`\overline{v}_l`$ and $`\overline{a}_l`$ of leptons (see equation (33)) and it gives the leptonic forward-backward asymmetries $`A_{FB}^l`$ at the Z-pole in all order of the perturbation theory as well as the tau polarization $`A_\tau ^{\mathrm{pol}}`$ and left-right polarization asymmetries $`A_{\mathrm{pol}}`$ $$A_{\mathrm{FB}}^l=\frac{3}{4}A_eA_f,A^{\mathrm{pol}}=A_\tau ,A_{\mathrm{LR}}=A_e$$ (78) where $$A_f=\frac{2\overline{v}_v\overline{a}_f}{\overline{v}_f^2+\overline{a}_f^2}.$$ (79) Measurements of the asymmetries hence are measurements of the effective mixing angle and therefore it is a physical quantity as well the dimensionless parameter $`\mathrm{\Delta }k^{^{}}`$ defined by equation (77). The leptonic width depends on the vector and axial vector coupling and on the corrections to the $`Z`$-propagator. This requires the introduction of the so called $`\rho `$-parameter $$\mathrm{\Gamma }_l=\frac{G_\mu M_\mathrm{Z}^3}{6\pi \sqrt{2}}\left(\overline{g}_{Vl}^2+\overline{g}_{Al}^2\right)\rho .$$ (80) The corrections $`\mathrm{\Delta }r_W`$, $`\mathrm{\Delta }k^{}`$, $`\mathrm{\Delta }\rho `$ are known in various schemes and play an important role in the analysis of electroweak physics, because they give the precise predictions of the theory for simple observables as $`M_\mathrm{W}`$, the leptonic asymmetries etc. in terms of $`\alpha ,G_\mu `$ and $`M_\mathrm{Z}`$. It is very useful to have the results in different schemes since it allows for cross-checking the correctness of the result and to estimate the remaining theoretical errors given by the missing higher order contributions. In the precision tests assuming that the analysis is not restricted to the Standard Model the radiative corrections it is convenient to use the $`ϵ`$ parameters defined as $$ϵ_1=\mathrm{\Delta }\rho ,ϵ_2=c^2\mathrm{\Delta }\rho +\frac{s^2\mathrm{\Delta }r_W}{(c^2s^2)},ϵ_3=c^2\mathrm{\Delta }\rho +(s^2c^2)\mathrm{\Delta }k^{^{}}.$$ (81) The electroweak radiative corrections are dominated by two leading contributions: the running of the electromagnetic coupling and large $`m_t`$ effects to $`\mathrm{\Delta }\rho `$ $`(\mathrm{\Delta }\rho _t3G_\mu m_t^2/(8\pi ^2\sqrt{2}))`$. These corrections can be absorbed into the parameters of the Born cross section when we get improved Born approximation. #### 2.2.2 Running electromagnetic coupling The running of $`\alpha `$ is completely given by the photon self energy contributions $$\alpha (M_\mathrm{Z})=\frac{\alpha }{1\mathrm{\Delta }\alpha }$$ (82) where $$\mathrm{\Delta }\alpha =\mathrm{Re}\left(\widehat{\mathrm{\Pi }}^\gamma (M_\mathrm{Z}^2)\right)=\mathrm{Re}\left(\mathrm{\Pi }^\gamma (M_\mathrm{Z}^2)\right)+\mathrm{Re}\left(\mathrm{\Pi }^\gamma (0)\right).$$ (83) The self energy contribution is large ($`6\%`$). It can be split into leptonic and hadronic contributions $$\mathrm{\Delta }\alpha =\mathrm{\Delta }\alpha _{\mathrm{lept}}+\mathrm{\Delta }\alpha _{\mathrm{had}}$$ (84) The leptonic part is known up to three loop $$\mathrm{\Delta }\alpha _{\mathrm{lept}}=314.97687(16)\times 10^4$$ (85) and the remaining theoretical error is completely negligible. The hadronic contribution is more problematic as it can not be calculated theoretically with the required precision since the light quark loop contributions have non-perturbative QCD effects. One can extract it, however, from the data using the relation $`\mathrm{\Delta }\alpha _{\mathrm{had}}`$ $`=`$ $`{\displaystyle \frac{\alpha }{3\pi }}M_\mathrm{Z}^2\mathrm{Re}{\displaystyle _{4m_\pi ^2}^{\mathrm{}}}𝑑s{\displaystyle \frac{R_{e^+e^{}}(s^)}{s^{}(s^{}M_\mathrm{Z}^2iϵ)}}`$ $`R_{e^+e^{}}(s)`$ $`=`$ $`{\displaystyle \frac{\sigma (e^+e^{}\gamma ^{}\mathrm{hadrons})}{\sigma (e^+e^{}\mu ^+\mu ^{})}}.`$ (86) Conservatively, one calculates the high energy $`\sqrt{s}40\mathrm{GeV}`$ contribution using perturbative QCD. The low energy contribution $`\sqrt{s}40\mathrm{GeV}`$ is estimated using data . Unfortunately, the precision of the low energy data is not good enough and the error from this source dominates the error of the theoretical predictions $$\mathrm{\Delta }\alpha _{\mathrm{had}}=0.02804\pm 0.00064,\alpha ^1(M_\mathrm{Z})=128.89\pm 0.09.$$ (87) One can, however, achieve a factor of three reduction of the estimated error assuming that the theory can be used down to $`\sqrt{s}=m_\tau `$ when quark mass effects can be included up to three loops. Such an analysis is quite well motivated by the successful results on the tau lifetime. In the hadronic vacuum polarization the non-perturbative power corrections appear to be suppressed and the unknown higher order perturbative contributions are relatively small. In this theory driven approach the error is reduced to an acceptable $`0.25\%`$ value $$\alpha ^1(M_\mathrm{Z})=128.905\pm 0.036.$$ (88) It is unlikely that the low energy hadronic total cross section will be measured in the foreseeable future with a precision leading to essential improvement. #### 2.2.3 Calculation of $`\mathrm{\Delta }\rho _t`$ We have noted in section 1.26 that in the limit of custodial symmetry in leading order $`\rho =1`$ and the dominant radiative correction comes from the virtual effects of the top quark since the top-bottom mass splitting gives the largest violation of custodial symmetry. The importance of this correction was first pointed out by Veltman . It is an elegant technical trick to calculate this correction using the effective field theory obtained in the $`m_t\mathrm{}`$ limit . In this limit we need to keep only the third generation and the gauge bosons can be treated as external classical currents without kinetic terms. The Standard Model Lagrangian can be reduced to the terms $$_{\mathrm{eff}}=i\overline{\psi }_L^Q\gamma _\mu D^\mu \psi _L^Q+i\overline{\psi }_R^t\gamma _\mu D^\mu \psi _R^t+i\overline{\psi }_R^b\gamma _\mu D^\mu \psi _R^b+\lambda _t\overline{\psi }_L^Q\mathrm{\Phi }\psi _r^tV(\mathrm{\Phi })$$ (89) with the Higgs doublet $`\mathrm{\Phi }=\left(\begin{array}{c}\varphi ^+\\ \frac{1}{\sqrt{2}}(v+h+i\chi )\end{array}\right)`$ (92) where $`\chi `$ and $`\varphi ^+`$ are the Goldstone bosons. The renormalized Lagrangian then has the form $$_{\mathrm{eff}}=Z_2^\varphi \left\|_\mu \varphi ^+i\frac{gv}{2}W_\mu ^+\right\|^2+\frac{Z_2}{2}^\chi \left\|_\mu \chi i\frac{gv}{2c}Z_\mu \right\|^2+\frac{Z_h}{2}\left(_\mu h\right)^2+\mathrm{}$$ (93) where we dropped the top and bottom kinetic energy terms, the top mass terms and the gauge boson fermion couplings. In this limit the gauge boson couplings do not get corrections but the $`Z`$ and $`W`$ mass terms are modified by the self energy corrections of the Goldstone bosons $$M_\mathrm{W}^2=Z_2^\varphi \frac{g^2v^2}{4},M_\mathrm{Z}^2=Z_2^\chi \frac{g^2v^2}{4c^2}$$ (94) and therefore the correction to the $`\rho `$ parameter is $$\mathrm{\Delta }\rho =\frac{Z_2^\varphi }{Z_2^\chi }1$$ (95) that is we can get the correction to the $`\rho `$ parameter simply by calculating the the contributions of top and top-bottom fermion loops to difference of the self energies of the neutral and charged Goldstone bosons. Carrying out this simple calculation we can easily check that the answer is $$\mathrm{\Delta }\rho _t3\frac{G_Fm_t^2}{8\pi \sqrt{2}}.$$ (96) The method can be extended to two loop order and the corresponding two loop calculation has been carried out in confirming previous result . #### 2.2.4 Higher order corrections to $`M_\mathrm{W}`$ and the mixing angle As we noted above, the simplest physical observables for precise tests of the Standard Model are $`M_\mathrm{W}`$ and the $`\mathrm{sin}^2\theta _W^{eff}`$. It is convenient to consider the radiative corrections in the $`\overline{\mathrm{MS}}`$ scheme where with good accuracy $`\mathrm{sin}^2\theta _W^{eff}\widehat{s}^2`$. It is given in terms of the input parameters via the relation $$\widehat{s}^2\widehat{c}^2=\frac{\pi \alpha (M_\mathrm{Z})}{\sqrt{2}G_\mu M_\mathrm{Z}^2(1\widehat{r}_W)}$$ (97) where $`\widehat{r}_W=0`$ in leading order. Using the measured value of $`\mathrm{sin}^2\theta _W^{eff},M_\mathrm{Z}`$ and $`G_\mu `$ we obtain a value $`\widehat{r}_W=0.0058\pm 0.000480`$ different from zero at the $`12\sigma `$ level. If one carries out a similar analysis for $`M_\mathrm{W}`$ the evidence for the presence of subleading corrections is even better. The radiative correction $`\widehat{r}_W`$ does not contain the large effect from the running $`\alpha `$ but it receives large custodial symmetry violating corrections because of the large top-bottom mass splitting $$\mathrm{\Delta }\widehat{r}_W\|_{\mathrm{top}}=c^2/s^2\mathrm{\Delta }\rho 0.0096\pm 0.00095.$$ (98) Subtracting this value we get about $`6\sigma `$ difference coming from the loops involving the bosonic sector (W,Z,H) and subleading fermionic contributions. At this level of accuracy many other corrections start to become important and the size of errors coming from the errors in the input parameters leads to effects of the same order. In particular, we get some sensitivity to the value of the Higgs mass. Beyond the complete one loop corrections it was possible to evaluate all important two loop corrections: $`𝒪(\alpha ^2\mathrm{ln}(M_\mathrm{Z}/m_f)`$ corrections with light fermions, mixed electroweak QCD corrections of $`𝒪(\alpha \alpha _s)`$, two loop electroweak corrections enhanced by top mass effects of $`𝒪(\alpha ^2(m_t^2/M_\mathrm{W}^2)^2)`$ together with the subleading parts of $`𝒪(\alpha \alpha _s^2m_t^2/M_\mathrm{W}^2)`$ and the very difficult subleading correction of $`𝒪(\alpha ^2m_t^2/M_\mathrm{W}^2)`$. It is remarkable that this last contribution proved to be important in several respect . Its inclusion reduced significantly the scheme dependence of the results and lead to a significant reduction of the upper limit on the Higgs mass. #### 2.2.5 Global fits This summer the LEP experiments and SLD could finalize their results on the electroweak precision data. The most important development is that the final value of SLD on the leptonic polarization asymmetry which implies $`\mathrm{sin}^2\theta _W^{eff}=0.23119\pm 0.00020`$. A nice summary of the results is given in Fig. 4. According to a recent analysis of the EWWW working group , the new world average is $$\mathrm{sin}^2\theta _W^{eff}=0.23151\pm 0.00017\mathrm{with}\chi ^2/\mathrm{d}.\mathrm{o}.\mathrm{f}=13.3/7.$$ (99) This gives only rather low confidence level of $`6.4\%`$. The origin of this unsatisfactory result is the $`2.9\sigma `$ discrepancy between the values $`\mathrm{sin}^2\theta _W^{eff}`$ deriving from the SLAC leptonic polarization asymmetry data and from the forward backward asymmetry in the b-b channel at LEP and SLC. The results obtained from a global fit to all data give somewhat better result but there we are hampered with the problem that the polarization asymmetry parameters disagree with each other with $`2.7\sigma `$, therefore the $`\chi ^2`$ is relatively large. #### 2.2.6 Upper limit on $`M_\mathrm{H}`$ The final results of the electroweak radiation corrections for $`M_\mathrm{W}`$ and $`\mathrm{sin}^2\theta _W^{eff}`$ can be parameterized in terms of the input parameters including their errors in simple approximate analytic form . For example in the $`\overline{MS}`$-scheme one obtains for the W-mass $`M_\mathrm{W}`$ $`=`$ $`80.38270.0579\mathrm{ln}({\displaystyle \frac{M_\mathrm{H}}{100}})0.008\mathrm{ln}^2({\displaystyle \frac{M_\mathrm{H}}{100}})`$ (100) $`0.517\left({\displaystyle \frac{\mathrm{\Delta }\alpha _h^{(5)}}{0.0280}}1\right)+0.543\left[\left({\displaystyle \frac{m_t^2}{175}}1\right)\right]`$ $`0.085\left({\displaystyle \frac{\alpha _s(M_\mathrm{Z})}{0.118}}1\right)`$ where $`m_t`$, $`M_\mathrm{H}`$ and $`M_\mathrm{W}`$ are in $`\mathrm{GeV}`$ units. This formula accurately reproduces the result obtained with numerical evaluation of all corrections in the range $`75\mathrm{GeV}M_\mathrm{H}350\mathrm{GeV}`$ with maximum deviation of less than $`1\mathrm{MeV}`$. In Fig. 5 the measured values of $`M_\mathrm{W}`$ are summarized . Using the world average $`M_\mathrm{W}=80.394\pm 0.042\mathrm{GeV}`$ with input parameters $`\alpha _s=0.119\pm 0.003`$, $`m_t=174.3\pm 5.1\mathrm{GeV}`$, $`\delta \alpha ^{(5)}=0.02804\pm 0.00065`$ one obtains at $`95\%`$ confidence level an allowed range for the Higgs mass of $`73\mathrm{GeV}M_\mathrm{H}294\mathrm{GeV}`$. Similar result exists also for $`\mathrm{sin}^2\theta _W^{eff}`$ extracted from the asymmetry measurements at the Z-pole with somewhat better (95% confidence) limits of $`95\mathrm{GeV}M_\mathrm{H}260\mathrm{GeV}`$. Without global fits we got a semi-analytic insight on the sensitivity of the precision tests to the Higgs mass. We also see that the precise measurements of $`M_\mathrm{W}`$ have already provided us with competitive values in comparison with those obtained from the measurement of $`\mathrm{sin}^2\theta _W^{eff}`$. It is interesting that the values of the Higgs mass obtained in a recent global fit are in good agreement with the simple analysis based on the value of $`M_\mathrm{W}`$ or $`\mathrm{sin}^2\theta _W^{eff}`$ as described above. One obtains an expected value for the Higgs boson of $`160170\mathrm{GeV}`$ with error of $`\pm 5060\mathrm{GeV}`$. The 95% confidence level upper limit is about $`260290\mathrm{GeV}`$. In Fig. 6 $`\mathrm{\Delta }\chi ^2`$ is plotted as the function of $`M_\mathrm{H}`$. #### 2.2.7 Can the Higgs boson be heavy? The precision data can not yet rule out dynamical symmetry breaking with some heavy Higgs like scalar and vector resonances. The minimal model to describe this alternative is obtained by assuming that the new particles are heavy (more than 0.5 TeV) and the linear $`\sigma `$-model Higgs sector of the Standard Model is replaced by the non-renormalizable non-linear $`\sigma `$-model. It can be derived also as an effective chiral vector-boson Lagrangian with non-linear realization of the gauge symmetry . How can we reconcile this more phenomenological approach with the precision data? Removing the Higgs boson from the Standard Model while keeping the gauge invariance is a relatively mild change. Although the model becomes non-renormalizable, at the one-loop level the radiative effects grow only logarithmically with the cut-off at which new interactions should appear. In equation (100) the Higgs mass is replaced by this cut-off. The logarithmic terms are universal, therefore their coefficients must remain the same. The constant terms, however, can be different from those of the Standard Model. The one loop corrections of the effective theory require the introduction of new free parameters which influence the value of the constant terms. The data, unfortunately, do not have sufficient precision to significantly constrain the constant terms appearing in $`M_W`$, $`\mathrm{sin}^2\theta _W^{eff}`$ and $`\mathrm{\Gamma }_l`$ (or alternatively in the parameters $`ϵ_1,ϵ_2,ϵ_3`$ or $`S,T,U`$ ). In a recent analysis it has been found that due to the screening of the symmetry breaking sector , alternative theories with dynamical symmetry breaking and heavy scalar and vector bosons still can be in agreement with the precision data up to a cut-off scale of $`3\mathrm{T}\mathrm{e}\mathrm{V}`$. #### 2.2.8 W-pair production At LEP the precise measurement of the production of $`W^+W^{}`$ is also an important physics goal. The production of gauge boson pairs provide us with the best test of the non-Abelian gauge symmetry of the Standard Model. Deviation from the Standard Model predictions may come either from the presence of anomalous couplings or the production of new heavy particles and their decays into vector-boson pairs. If the particle spectrum of the Standard Model has to be enlarged with new particles (as in the Minimal Supersymmetric Standard Model) with mass values of $`0.51\mathrm{T}\mathrm{e}\mathrm{V}`$, small anomalous couplings are generated at low energy. In Fig. 7 we show the recent measurement of the W-pair cross-section. ## 3 HIGGS SECTOR, HIGGS SEARCH ### 3.1 Difficulties with the Higgs sector The Standard Model is defined only in perturbation theory, but the perturbative treatment of the Higgs sector can not be valid up to arbitrary high energies. Its range of validity depends strongly on the value of the Higgs mass. #### 3.1.1 Theoretical upper limits on $`M_\mathrm{H}`$ For high Higgs mass values we get conflict with perturbative unitarity since according to equation (46) if $`M_\mathrm{H}>>M_\mathrm{W}`$ the scalar self interaction becomes strong. Unitarity requires that in a given angular momentum channel the scattering matrix element fulfills the relations $$\left\|M_J\right\|^2\left\|\mathrm{Im}(M_J)\right\|,\left\|\mathrm{Re}(M_J)\right\|<\frac{1}{2}$$ (101) Applying these constraints for the Born amplitude of $`W_LW_L`$ scattering we get $$M_\mathrm{H}^2\frac{2\pi \sqrt{2}}{G_F}(850\mathrm{GeV})^2$$ (102) A more precise coupling channel analysis leads to somewhat better limit $$M_\mathrm{H}^2\frac{2\pi \sqrt{2}}{G_F}(700\mathrm{GeV})^2.$$ (103) If the Higgs self coupling is large, the gauge interactions are negligible. The scalar interaction, however, is not asymptotically free, therefore in the perturbation theory the running coupling constant has a Landau pole. Actually it has been proven that the scalar theory is trivial . If we require to have finite scalar coupling at very short distances we get vanishing coupling at large distances: the theory becomes free. The scalar sector mathematically can not be rigorous and it should be considered as an effective low energy theory. The scalar self coupling $`\lambda `$, therefore, has to be smaller than its value at the Landau pole. This condition also gives an upper limit on $`M_\mathrm{H}`$. The one loop running coupling is $$\lambda (\mu )=\frac{\lambda (M_\mathrm{H})}{112\frac{\lambda (M_\mathrm{H})}{16\pi ^2}\mathrm{ln}\frac{\mu ^2}{M_\mathrm{H}^2}}$$ (104) where $`\lambda (M_\mathrm{H})=g^2M_\mathrm{H}^2/(8M_\mathrm{W}^2)`$. The position of the Landau pole is at the scale $$\mu _c=M_\mathrm{H}e^{\frac{2\pi ^2}{3\lambda (M_\mathrm{H})}}=M_\mathrm{H}e^{\frac{16\pi ^2M_\mathrm{W}^2}{3g^2M_\mathrm{H}^2}}$$ (105) the condition $`\mu _c>2M_\mathrm{H}`$ leads to the upper limit $$M_H<700\mathrm{G}\mathrm{e}\mathrm{V}.$$ (106) If the Landau pole is pushed up to the Planck scale we get the more stringent limit of $`M_\mathrm{H}<170\mathrm{GeV}`$. This is a tentative estimate since perturbation theory is used beyond its range of validity. Note that non-perturbative lattice studies give very similar value $`M_\mathrm{H}<650\mathrm{GeV}`$. One can redo the analysis at two loop order when the two loop beta function has a metastable fixed point. With the assumption that the theory is meaningful up to a scale where the coupling constant is half of the value of the coupling at the metastable fixed point one, gets again similar upper limit . #### 3.1.2 Lower limits on $`M_\mathrm{H}`$ The requirement of stability of the Higgs potential $`V(\varphi )>0`$ leads to a lower limit on the Higgs mass. The $`\beta `$-function of the scalar self-interaction is $`\beta _\lambda `$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}[12\lambda ^23g_t^4+6\lambda g_t^2{\displaystyle \frac{\lambda }{2}}(9g^{}_{}{}^{}2+3g^2)`$ (107) $`+{\displaystyle \frac{3}{16}}(3g^{}_{}{}^{}4+2g^{}_{}{}^{}2g^2+g^4)]`$ where $`g_t=gm_t/(2M_\mathrm{W})`$ denotes the top quark Yukawa coupling. At large top mass and small $`\lambda `$ the second term will dominate, as we get negative $`\beta `$ function and $`\lambda `$ will decrease with increasing top mass and $`V(\mathrm{\Phi })`$ can become negative. A coupled channel analysis at two loop order leads to the approximate relation (assuming that the theory is meaningful up to the Planck scale) $`M_\mathrm{H}>1.95m_t190\mathrm{GeV}`$ and for $`m_t=175\mathrm{GeV}`$ one gets the lower limit $`M_\mathrm{H}>150\mathrm{GeV}`$. So the Higgs boson should not be found at LEP2 in this case. In Fig. 8 theoretical upper and lower bounds on $`M_H`$ are shown as a function of the scale charaterizing the range of validity (cut-off scale) of the Standard Model. We can see that if $`M_\mathrm{H}165195\mathrm{GeV}`$ the cut-off scale can extend up to the Planck scale. I should recall, however, that this is rather unlikely since the mass term of the scalar boson is a relevant operator, therefore, it is linearly sensitive of the scale of new physics. The question “why is the Higgs mass so light with respect to the Planck scale” is referred to in literature as the gauge hierarchy problem. ### 3.2 Search at LHC One of the most important physics goal at the LHC is to obtain decisive experimental test on the Higgs sector of the Standard Model . The experimental prospects are summarized with the so called “no loose scenario * either the Higgs boson will be discovered at LEP2 or LEP2 will establish a lower limit of $`m_H>107\mathrm{GeV}`$; * assuming the validity of the minimal Higgs sector the precision data obtained at LEP1 and SLC constrain the value of $`M_\mathrm{H}`$ to be less than $`295\mathrm{GeV}`$. * LHC will be able to discover the Standard Model Higgs boson in the interval $`800\mathrm{GeV}>m_H>107\mathrm{GeV}`$ or will find clear evidence for deviations from Standard Model predictions. To be able to make maximal use of the results of the LHC experiments one should calculate the Standard Model predictions for LHC processes as precisely as possible. It is unsatisfactory that the present experimental simulation studies are based on leading order cross sections. Nevertheless at the parton level most of the NLO corrections are available and are public in form of program packages . HDECAY generates all branching fractions of the Standard Model Higgs boson and the Higgs bosons of the Minimal Supersymmetric Standard Model (MSSM), while HGLUE provides the production cross sections of the SM and MSSM Higgs bosons via gluon fusion including the NLO QCD corrections. #### 3.2.1 Higgs branching ratios The branching ratios of the Higgs boson have been studied in many papers. A useful compilation of the early works on this subject can be found in Ref. , where the most relevant formulae for on-shell decays are summarized. Higher-order corrections to most of the decay processes have been computed (for up-to-date reviews see Refs. and references therein). The bulk of the QCD corrections to $`Hq\overline{q}`$ can be absorbed into a ‘running’ quark mass $`m_q(\mu )`$, evaluated at the energy scale $`\mu =M_H`$. The importance of this effect for the case $`q=b`$, with respect to intermediate-mass Higgs searches at the LHC, has been discussed already in Ref. . For sake of illustration, results on the Higgs branching ratios are summarized in Fig. 9. Branching ratios are depicted as function of $`M_\mathrm{H}(M_\mathrm{H}200\mathrm{GeV})`$ for channels: (a) $`b\overline{b}`$, $`c\overline{c}`$, $`\tau ^+\tau ^{}`$, $`\mu ^+\mu ^{}`$ and $`gg`$; and for channels (b) $`WW`$, $`ZZ`$, $`\gamma \gamma `$ and $`Z\gamma `$. The patterns of the various curves are not significantly different from those presented in Ref. . The inclusion of the QCD corrections in the quark-loop induced decays give a change of at most a few per cent for the decays $`H\gamma \gamma `$ and $`HZ\gamma `$, while for $`Hgg`$ differences are of order 50–60%. However, this latter result has little phenomenological relevance, since this decay width makes a negligible contribution to the total width, and in practice it is an unobservable channel. #### 3.2.2 Higgs production cross sections and event rates There are only four Higgs production mechanisms which lead to detectable cross sections at the LHC: a) gluon-gluon fusion , b) $`WW`$, $`ZZ`$ fusion , c) associated production with $`W`$, $`Z`$ bosons , d) associated production with $`t\overline{t}`$ pairs . Each mechanism involves heavy particles. Representative Feynman diagrams are shown in Fig. 10. Again for illustrative purpose in Fig. 11 total cross-section values are depicted for LHC energies$`\sqrt{s}=14\mathrm{TeV}`$ . There are various uncertainties in the rates of the above processes, although none is particularly large. They are given by the lack of precise knowledge of the gluon distribution at small $`x`$ and by effects of unknown higher-order perturbative QCD corrections . The next-to-leading order QCD corrections are known for processes (a), (b) and (c) and are included. By far the most important of these are the corrections to the gluon fusion process calculated in Ref. . Within the limit where the Higgs mass is far below the $`2m_t`$ threshold, these corrections are calculable analytically . In fact, it turns out that the analytic result is a good approximation up to the threshold $`M_\mathrm{H}<2m_t`$ . In Ref. the impact of the next-to-leading order QCD corrections for the gluon fusion process on LHC cross sections has been investigated, both for the SM and for the MSSM. Overall, the next-to-leading order correction increases the leading order result by a factor of about 2, when the normalization and factorization scales are set equal to $`\mu =M_H`$. This ‘$`K`$-factor’ can be traced to a large constant piece in the next-to-leading correction $$K1+\frac{\alpha _s(M_H)}{\pi }\left[\pi ^2+\frac{11}{2}\right].$$ (108) Such a large $`K`$-factor usually implies a non-negligible scale dependence of the theoretical cross section. To judge the quality of the various signals of Higgs production, we must know event rates both for the signals and the backgrounds. Considering all the possible combinations of production mechanisms and decay channels , the best chance of discovering a Higgs at the LHC are given by the signatures: (i) $`ggH\gamma \gamma `$, (ii) $`q\overline{q}^{}WH\mathrm{}\nu _{\mathrm{}}\gamma \gamma `$ and (iii) $`ggHZ^{()}Z^{()}\mathrm{}^+\mathrm{}^{}\mathrm{}^+\mathrm{}^{}`$, where $`\mathrm{},\mathrm{}^{}=e`$ or $`\mu `$. By exploiting techniques of flavor identification of $`b`$-jets, thereby reducing the huge QCD background from light-quark and gluon jets, the modes (iv) $`q\overline{q}^{}WH\mathrm{}\nu _{\mathrm{}}b\overline{b}`$ and (v) $`gg,q\overline{q}t\overline{t}Hb\overline{b}b\overline{b}WWb\overline{b}b\overline{b}\mathrm{}\nu _{\mathrm{}}X`$, can also be used to search for the Higgs . Another potentially important channel, particularly for the mass range $`2M_W\stackrel{<}{}M_H\stackrel{<}{}2M_Z`$, is (vi) $`HW^{()}W^{()}\mathrm{}^+\nu _{\mathrm{}}\mathrm{}^{}\overline{\nu }_{\mathrm{}^{}}`$ . Here the lack of a measurable narrow resonant peak is compensated by a relatively large branching ratio, since for this mass range $`HWW`$ is the dominant decay mode. Again for sake of illustration we show in Fig. 12 Higgs production times branching ratios for various decay modes at two different energies. The potential of the ATLAS experiment for the discovery of the Standard Model Higgs boson in the mass range $`80\mathrm{GeV}<M_\mathrm{H}<1\mathrm{TeV}`$ is summarized in Fig. 13. The significance of the signal depends on the signal (S) and background events and is given by $`S/\sqrt{B}`$. The results shown in the figure are obtained from calculating the event rates both for the background and the signal in the Born approximation. It is argued that since the QCD corrections are not known for all signal and background processes, it is more consistent to neglect them everywhere. Hopefully this shortcomings will be eliminated soon. One can consider these result conservative since the QCD corrections are large for the signal and there is no reason to assume that they are even larger for the background.
warning/0004/hep-lat0004018.html	ar5iv	text	# A Lattice Study of the Glueball Spectrum† ## 1 Introduction It is believed that QCD is the theory which describes strong interactions among quarks and gluons. A direct consequence of this is the existence of excitations of gluonic degrees of freedom, i.e. glueballs. However, due to their non-perturbative nature, the spectrum of glueballs can only be investigated reliably with non-purterbative methods like lattice QCD . Recently, it has become clear that such a calculation can be performed on a relatively coarse lattice using an improved gluonic action on asymmetric lattices . In this paper, we present our results for a glueball spectrum calculation. The spatial lattice spacing in our simulations ranges from $`0.08fm`$ to $`0.25fm`$ which enables us to extrapolate more reliably to the continuum limit. The improved gluonic action we used is the tadpole improved gluonic action on asymmetric lattices as described in . It is given by: $`S=`$ $``$ $`\beta {\displaystyle \underset{i>j}{}}\left[{\displaystyle \frac{5}{9}}{\displaystyle \frac{TrP_{ij}}{\xi u_s^4}}{\displaystyle \frac{1}{36}}{\displaystyle \frac{TrR_{ij}}{\xi u_s^6}}{\displaystyle \frac{1}{36}}{\displaystyle \frac{TrR_{ji}}{\xi u_s^6}}\right]`$ (1) $``$ $`\beta {\displaystyle \underset{i}{}}\left[{\displaystyle \frac{4}{9}}{\displaystyle \frac{\xi TrP_{0i}}{u_s^2}}{\displaystyle \frac{1}{36}}{\displaystyle \frac{\xi TrR_{i0}}{u_s^4}}\right].`$ In the above expression, $`\beta `$ is related to the bare gauge coupling, $`\xi =a_s/a_t`$ is the (bare) aspect ratio of the asymmetric lattice with $`a_s`$ and $`a_t`$ being the lattice spacing in spatial and temporal direction respectively. The parameter $`u_s`$ is the tadpole improvement parameter to be determined self-consistently from the spatial plaquettes in the simulation. $`P_{ij}`$ and $`P_{0i}`$ are the spatial and temporal plaquette variables. $`R_{ij}`$ designates the $`2\times 1`$ Wilson loop ($`2`$ in direction $`i`$ and $`1`$ in direction $`j`$). Using spatially coarse and temporally fine lattices helps to enhance signals in the glueball correlation functions. Therefore, the bare aspect ratio is taken to be some value larger than one. In our simulation, we have used $`\xi =3`$ for our glueball calculation. It turns out that, using the non-perturbatively determined tadpole improvement parameter $`u_s`$, the renormalization effects of the aspect ratio is small , typically of the order of a few percent for practical values of $`\beta `$ in the simulation. This could also be verified by measuring corresponding Wilson loops, which will directly yield the renormalized aspect ratio. For the moment, we will ignore their difference and simply use the bare value of $`\xi `$. It is also noticed that, without tadpole improvement, this renormalization effect could be as large as $`30`$ percent . This paper is organized as follows. In the next section, we will explain our calculation of the Wilson loops, glueball correlation functions which give us the mass values of the glueballs in various symmetry sectors. Finite volume errors are discussed and more importantly, finite lattice spacing errors are analyzed. Special attention is paid to the scalar sector where the extrapolation used to be troublesome at coarse lattices. In the last section, we will have some discussion of out result and conclude. ## 2 Monte Carlo Simulations Glueball spectrum calculations in pure gauge theory basically involve three stages. The first stage, gauge field configurations are generated using some algorithm. We have utilized a Hybrid Monte Carlo algorithm to update gauge field configurations. Several lattice sizes have been simulated and the detailed information can be found in Table.1. For each lattice with fixed bare parameters, order of $`10^3`$ configurations have been accumulated. Each gauge field configuration is separated from the previous one by several Hybrid Monte Carlo trajectories, typically $`510`$, to make sure that they are sufficiently decorrelated. Further binning of the data has been performed and no noticeable remaining autocorrelation has been observed. The second stage of the calculation is to perform measurements of physical observables using the accumulated gauge field configurations. In fact, two independent measurements have to be done. One is to set the scale in physical units, i.e. to determine the lattice spacing $`a_s`$ in physical units. This is necessary since there is no scale in a pure gauge theory. The second is to measure glueball mass values in lattice units. This is done by measuring various glueball correlation functions. With the scale being set in the first step, the mass values of the glueballs can then be converted into physical units. In the final stage of the calculation, glueball mass values obtained from a finite lattice have to be extrapolated to the continuum limit. This means that finite volume effects and finite lattice spacing errors have to be eliminated. Typically, finite volume errors are found to be small in these calculations . It is the finite lattice spacing errors that are more difficult to handle, especially for the scalar glueball mass. It was observed that, the scalar glueball sector exhibits a dip in the extrapolation, making the extrapolation less dependable compared with other channels . To remedy this situation, a simulation at a smaller values of $`a_s`$ has been performed. In our calculation, we also performed a simulation at a smaller lattice spacing where $`a_s0.08fm`$. These two simulations now makes the extrapolation in the scalar sector more dependable and less sensitive to the form of the functions used in the extrapolation. ### 2.1 Setting the scale In our simulation, the scale is set by measuring Wilson loops from which the static quark anti-quark potential $`V(R)`$ is obtained. Using the static potential between quarks, we are able to determine the lattice spacing in physical units by measuring $`r_0`$, a pure gluonic scale determined from the static potential . The definition of the scale $`r_0`$ is given by: $`R^2dV(R)/dR]_{R=r_0}=1.65`$. In physical units, $`r_0`$ is roughly $`0.5fm`$ which is determined by comparison with potential models. For a recent determination of $`r_0`$, please consult Ref. . In order to measure the Wilson loops accurately, it is the standard procedure to perform single link smearing on the spatial gauge links of the configurations. In this procedure, each spatial gauge link is replaced by a linear combination of the original link and its spatial staples. Each spatial gauge staple is weighted by a parameter $`\lambda _W`$ relative to the original gauge link. The final result is then projected back into $`SU(3)`$. This smearing scheme can be performed iteratively on the spatial links of gauge fields for as many as $`n_W`$ times. The effect of smearing is to projects out higher excitation contaminations from the Wilson loop measurements. Then, Wilson loops are constructed using these smeared links. For a Wilson loop of size $`R\times T`$, it is fitted against: $$W(R,T)\stackrel{T\mathrm{}}{}Z(R)e^{V(R)T}.$$ (2) Therefore, by extracting the effective mass plateau at large temporal separation, the static quark potential $`V(R)`$ is obtained. We have measured the Wilson loops along different lattice axis. It is seen that the measured data points for the quark anti-quark potential along different lattice axis lie on a universal line which is an indication that the improved action restores the rotational symmetry quite well. The smearing parameter $`\lambda _W`$ used in this process are also listed in Table. 1. Using smearing, we have been able to obtain descent plateau in the effective mass and the potential $`a_tV(R)`$ is thus determined for various values of $`R/a_s`$. The static quark potential is fitted according to a Coulomb term plus a linear confining potential which is known to work well at these lattice spacings . The potential is parameterized as: $$V(R)=V_0+e_c/R+\sigma R.$$ (3) From this and the definition of $`r_0`$, it follows that: $$r_0/a_s=\sqrt{\frac{1.65+e_c}{\sigma a_s^2}}.$$ (4) To convert the measured result to $`r_0/a_s`$, we have also used the value of $`\xi `$ taken as the bare value. As explained earlier, the renormalization effects for this parameter is small. The result of the spatial lattice spacing in physical units is also included in Table. 1. The errors for the ratio $`r_0/a_s`$ are obtained by blocking the whole data set into smaller blocks and obtain the error from different blocks. ### 2.2 Glueball mass measurement To obtain glueball mass values, it is necessary to constructed glueball operators in various symmetry sectors of interests. On a lattice, the full rotational symmetry is broken down to only cubic symmetry, which is a finite point group with $`5`$ irreducible representations <sup>1</sup><sup>1</sup>1Since we only consider the glueball states with zero momentum, it suffices to study the cubic group. For glueball states with non-vanishing momenta, the whole space group has to be considered.. They are labeled as : $`A_1`$, $`A_2`$, $`E`$, $`T_1`$ and $`T_2`$. The first two irreducible representations are one-dimensional. The third is two-dimensional while the last two are three-dimensional. In practice, one is interested in the scalar, tensor and vector glueballs in the continuum limit. The correspondence is the following: scalar glueball is in the $`A_1`$ representation of the cubic group; tensor glueball is in representation $`E+T_2`$, which forms a $`5`$-dimensional representation; vector glueballs are in the representation $`T_1`$. Of course, this correspondence is not one to one but infinite to one. Therefore, what we can measure is the lowest excitation in the corresponding representation of the cubic group . Glueball correlation functions are notoriously difficult to measure. They die out so quickly and it is very difficult to get a clear signal. In order to enhance the signal of glueball correlation functions, smearing and fuzzying have to be performed on spatial links of the gauge fields . These techniques greatly enhance the overlap of the glueball states and thus provide possibility of measuring the mass values. In our simulations, after performing single link smearing and double link fuzzying on spatial links, we first construct raw operators at a given time slice $`t`$: $`\{_t^{(i)}=_𝐱𝒲_{𝐱,t}^{(i)}\}`$, where $`𝒲_{t,𝐱}^{(i)}`$ are closed Wilson loops originates from a given lattice point $`x=(t,𝐱)`$. The loop-shapes studied in this calculation are shown in Fig.1. Then, elements from the cubic group is applied to these raw operators and the resulting set of loops now forms a basis for a representation of the cubic group. Suitable linear combinations of these operators are constructed to form a basis for a particular irreducible representation of interest . We denote these operators as $`\{𝒪_\alpha ^{(R)}(t)\}`$ where $`R`$ labels a specific irreducible representation and $`\alpha `$ labels different operators at a given time slice $`t`$. In order to maximize the overlap with one glueball state, we construct a glueball operator $`𝒢^{(R)}(t)=_\alpha v_\alpha ^{(R)}\overline{𝒪}_\alpha ^{(R)}(t)`$, where $`\overline{𝒪}_\alpha ^{(R)}(t)=𝒪_\alpha ^{(R)}(t)0𝒪_\alpha ^{(R)}(t)0`$. The coefficients $`v_\alpha ^{(R)}`$ are to be determined from a variational calculation. To do this, we construct the correlation matrix: $$𝐂_{\alpha \beta }(t)=\underset{\tau }{}0\overline{𝒪}_\alpha ^{(R)}(t+\tau )\overline{𝒪}_\alpha ^{(R)}(\tau )0.$$ (5) The coefficients $`v_\alpha ^{(R)}`$ are chosen such that they minimize the effective mass $$m_{eff}(t_C)=\frac{1}{t_C}\mathrm{log}\left[\frac{\underset{\alpha \beta }{}v_\alpha ^{(R)}v_\beta ^{(R)}𝐂_{\alpha \beta }(t_C)}{_{\alpha \beta }v_\alpha ^{(R)}v_\beta ^{(R)}𝐂_{\alpha \beta }(0)}\right],$$ (6) where $`t_C`$ is time separation for the optimization. In our simulation $`t_C=1`$ is taken. If we denote the optimal values of $`v_\alpha ^{(R)}`$ by a column vector $`𝐯^{(R)}`$, this minimization is equivalent to the following eigenvalue problem: $$𝐂(t_C)𝐯^{(R)}=e^{t_Cm_{eff}(t_C)}𝐂(0)𝐯^{(R)}.$$ (7) The eigenvector $`𝐯_0^{(R)}`$ with the lowest effective mass then yields the coefficients $`v_{0\alpha }^{(R)}`$ for the operator $`𝒢_0^{(R)}(t)`$ which best overlaps the lowest lying glueball in the channel with symmetry $`R`$. Higher-mass eigenvectors of this equation will then overlap predominantly with excited glueball states of a given symmetry channel. With these techniques, the glueball mass values are obtained in lattice units and the final results are listed in Table.2. The errors are obtained by binning the total data sets into several blocks and doing jackknife on the blocks. ### 2.3 Extrapolation to the continuum limit As has been mentioned, finite volume errors are eliminated by performing simulations at the same lattice spacing but different physical volumes. This also helps to purge away the possible toleron states whose energy are sensitive to the size of the volume. A simulation at a larger volume is done for the smallest lattice spacing in our calculation. We found that the mass of the scalar glueball remains unchanged when the size of the volume is increased. The mass of the tensor glueball seems to be affected, which is consistent with the known result that tensor glueballs have a rather large size and therefore feel the finiteness of the volume more heavily. The infinite volume is obtained by extrapolating the finite volume results using the relation : $$a_tM^{(R)}(L_s)=a_tM^{(R)}(\mathrm{})\left(1\lambda ^{(R)}\mathrm{exp}(\sqrt{3}z/2)/z\right),$$ (8) where $`z=M^{(A_1^{++})}L_s`$. Using the results for the mass of the $`E^{++}`$ and $`T_2`$ glueballs on $`8^324`$ and $`10^330`$ lattices for the same value of $`\beta `$, the final result for the mass of these glueball states are obtained. Glueball mass values for other symmetry sectors are not so sensitive to the finite volume effects. Therefore, in Table.2, only the extrapolated values for the smallest physical volume are tabulated. Other entries are obtained from $`8^324`$ lattice results. As for the finite lattice spacing errors, special attention is paid to the scalar glueball sector where the continuum limit extrapolation was known to have problems. Due to the simulation points at small lattice spacings, around $`0.1fm`$ and below, the ambiguity in this extrapolation is greatly reduced. We have tried to extrapolate the result using different formula suggested in Ref. , the extrapolated results are all consistent within statistical errors. For definiteness, we take the simple form: $$r_0M_G(a_s)=r_0M_G(0)+c_1(a_s/r_0)^2+c_2(a_s/r_0)^4,$$ (9) and the result is illustrated in Fig.2. The final extrapolated results for the glueball mass values are also listed in Table.2. The data points from our simulation results are shown with solid symbols and the corresponding extrapolations are plotted as solid lines. It is also noticed that the extrapolated mass values for $`E^{++}`$ and $`T_2^{++}`$ channels coincide within statistical errors, indicating that in the continuum limit, they form the tensor representation of the rotational group. For comparison, simulation results from are also shown with open symbols and the corresponding extrapolation are represented by the dashed lines. We also tried to extrapolate linearly in $`(a_s/r_0)^2`$ using three data points with smallest lattice spacing. The results are statistically consistent with the results using the extrapolation (9) within errors. It is seen that, due to data points at lattice spacings around $`0.1fm`$ and below, the uncertainties in the extrapolation for the glueball mass values are reduced. Finally, to convert our simulation results on glueball masses into physical units, we use the result $`r_0^1=410MeV`$. The errors for the hadronic scale $`r_0`$ is neglected. For the scalar glueball we obtain $`M_G(0^{++})=1730(90)MeV`$. For the tensor glueball mass in the continuum, we combine the results for the $`T_2^{++}`$ and $`E^{++}`$ channels and obtain $`M_G(2^{++})=2400(95)MeV`$ for the tensor glueball mass. ## 3 Discussions and Conclusions We have studied the glueball spectrum at zero momentum in the pure $`SU(3)`$ gauge theory using Monte Carlo simulations on asymmetric lattices with the lattice spacing in the spatial directions ranging from $`0.08fm`$ to $`0.25fm`$. This helps to make extrapolations to the continuum limit with more confidence for the scalar and tensor glueball states. The mass values of the glueballs are converted to physical units in terms of the hadronic scale $`r_0`$. We obtain the mass for the scalar glueball and tensor glueball to be: $`m_G(0^{++})=1730(90)MeV`$ and $`m_G(2^{++})=2400(95)MeV`$. It is interesting to note that, around these two mass values, experimental glueball candidates exist. Of course, in order to compare with the experiments other issues like the mixing effects and the the effects of quenching have to be studied.
warning/0004/cond-mat0004177.html	ar5iv	text	# Pinning phenomena in the Ginzburg-Landau Model of Superconductivity ## I Introduction Superconducting materials have the property of expelling an applied magnetic field. In fact, the behaviour of a superconducting sample varies according to the value of the applied field and the value of the Ginzburg-Landau parameter $`\kappa `$ which is characteristic of the material. When $`\kappa `$ is large, the superconductors are known as type-II and display vortex patterns for intermediate fields: for high magnetic fields, the material is normal and the magnetic field penetrates into the sample, for low fields, the material is superconducting, that is the magnetic field is expelled from the sample and for intermediate fields, there are vortices. The vortex state is a state where the superconducting and the normal phases coexist: at the center of the vortex, the material is normal and the vortex is circled by a superconducting current carrying a quantized amount of magnetic flux. The motion of vortices generates an electric field hence energy-dissipation. In order to have the desired property of dissipation-free current flow, the vortices have to be held fixed or pinned. In practice, attempts are made to pin vortices either by varying the thickness of the material or by introducing impurities or normal inclusions. Sufficiently strong pinning is necessary for functional superconductors capable of sustaining strong currents and high magnetic fields. The new high-temperature (high $`T_c`$) superconductors are strongly type-II superconductors, that is their phenomenology is dominated by the presence and properties of vortices when an exterior magnetic field is applied. The pinning problem is particularly intricate in high-$`T_c`$ superconductors where it depends on specific structures such as layering and structural defects. In this paper, we will be concerned with the case where the vortices are pinned by impurities in the framework of the Ginzburg-Landau model. We will study the behaviour of global minimizers of the Ginzburg-Landau energy when a term modelling the pinning of vortices by impurities is added, in the limit of a large Ginzburg-Landau parameter $`\kappa `$, which describes extreme type-II materials. ### I.1 The Ginzburg-Landau model with a pinning term Recall that in the framework of the Ginzburg-Landau theory (see \[T\] for more details), the state of the material is completely described by a vector potential $`A`$ and a complex-valued function $`u`$, which can be thought of as a wave-function of the superconducting electrons, and is nondimensionalized such that $`\|u\|1`$. The type of material is characterized by the Ginzburg-Landau parameter $`\kappa `$ and in the case of type II, $`\kappa `$ is large so that we define $`\epsilon =1/\kappa `$, which will be small. The energy is the following: (I.1) $$J_\epsilon (u,A)=\frac{1}{2}_\mathrm{\Omega }\|(iA)u\|^2+\frac{1}{2\epsilon ^2}\left(a_\epsilon (x)\|u\|^2\right)^2+\|hh_{\mathrm{ex}}\|^2.$$ Here, $`\mathrm{\Omega }`$ is the domain occupied by the superconductor, $`h=\mathrm{curl}A`$ is the magnetic field and $`h_{\mathrm{ex}}`$ is the exterior magnetic field which is constant in our problem. A common simplification is to restrict to a two-dimensional problem corresponding to an infinite cylindrical domain of section $`\mathrm{\Omega }^2`$ (smooth and simply connected), for an applied field parallel to the axis of the cylinder. Then $`A:\mathrm{\Omega }^2`$, $`h`$ is real-valued and all the quantities are translation-invariant. The energy $`J_\epsilon `$ that we are going to study here is slightly different from the classical Ginzburg-Landau energy in the sense that there is a term penalizing the variations of the order parameter $`u`$. We denote this function by $`a_\epsilon (x)`$. In the case originally studied by Ginzburg and Landau, $`a_\epsilon 1`$. In this paper, a typical example for $`a_\epsilon `$ would be to oscillate between $`1/2`$ and 1 in the domain, with a typical scale $`\eta `$ which may tend to 0 with $`\epsilon `$. The minima of $`a_\epsilon `$ correspond to the impurities in the material. Hence it is expected that these minima will be the pinning sites for the vortices. The modified Ginzburg-Landau functional (I.1) was first written down by Likharev \[L\]. Then, this model has been used and developed in \[CR\] and \[CDG\]. Review articles on the topic include \[BFGLV\], \[C1\], \[C2\] and \[P\]. Computational evidence that the vortices are attracted by the impurities, that is the points of minimum of $`a_\epsilon (x)`$ can be found in \[CDG\] or \[DGP\]. In this paper, we want to address the question of how the term $`a_\epsilon `$ will modify the properties of the superconductor in the presence of an exterior magnetic field. The method and techniques that we are going to use are inspired from those of \[SS3\] (in which the case $`a_\epsilon 1`$ was treated) and based on energy estimates, convergence of measures and construction of approximate solutions. Because of the term $`a_\epsilon (x)`$ in the equations, which can be a rapidly oscillating function, we will also need homogenization theory (\[CD\], \[JKO\], \[MuT\]) to describe the fact that the impurities, hence the vortices, form a homogenized medium in the material. ### I.2 The equation for the magnetic field The Ginzburg-Landau equations associated to the functional (I.1) when minimizing for $`\{(u,A)H^1(\mathrm{\Omega },)\times H^1(\mathrm{\Omega },^2)\}`$ are $`(\mathrm{G}.\mathrm{L}.)\{`$ $$\begin{array}{c}(iA)^2u=\frac{1}{\epsilon ^2}u(a_\epsilon (x)\|u\|^2)\hfill \\ ^{}h=<iu,(iA)u>,\hfill \end{array}$$ with the boundary conditions $$\{\begin{array}{cc}h=h_{\mathrm{ex}}\hfill & \text{on}\mathrm{\Omega }\hfill \\ (uiAu)n=0\hfill & \text{on}\mathrm{\Omega }.\hfill \end{array}$$ Here $`^{}`$ denotes $`(_{x_2},_{x_1})`$, and $`<z,w>=Re(z\overline{w})`$ for $`z,w`$ in $``$. Recall that the problem is invariant under the gauge transformations $$\{\begin{array}{c}uue^{i\mathrm{\Phi }}\hfill \\ AA+\mathrm{\Phi },\hfill \end{array}$$ where $`\mathrm{\Phi }H^2(\mathrm{\Omega },)`$. Physically meaningful quantities are gauge invariant. These include the energy $`J_\epsilon `$, the magnetic field $`h`$ and the superconducting current $`j=<iu,(iA)u>`$. Let us describe the properties of a superconductor. These phenomena are described for instance in \[T\]. The state of the material depends on the applied field $`h_{\mathrm{ex}}`$. In the absence of pinning, that is when $`a_\epsilon 1`$, there are two critical fields $`H_{c_1}`$ and $`H_{c_2}`$ for which a phase transition occurs. Above $`H_{c_2}=O(\frac{1}{\epsilon ^2})`$, superconductivity is destroyed and the material is in the normal phase $`(u0,hh_{\mathrm{ex}})`$. Below $`H_{c_1}=O(\|\mathrm{log}\epsilon \|)`$, the material is superconducting everywhere, that is $`\|u\|1`$. This is the Meissner phase characterized by complete expulsion of the magnetic field : in the limit when $`\epsilon `$ goes to zero, the magnetic field satisfies the London equation (I.2) $$\{\begin{array}{cc}\mathrm{\Delta }h+h=0\hfill & \text{in}\mathrm{\Omega }\hfill \\ h=h_{\mathrm{ex}}\hfill & \mathrm{on}\mathrm{\Omega }.\hfill \end{array}$$ Between $`H_{c_1}`$ and $`H_{c_2}`$, the material is in the mixed phase defined by the coexistence of the normal and superconducting phases in the form of vortex filaments: the magnetic field penetrates into the material in the form of flux lines at the center of which $`u`$ vanishes. The induced magnetic field approximately satisfies (I.3) $$\{\begin{array}{cc}\mathrm{\Delta }h+h=2\pi _id_i\delta _{p_i}\hfill & \text{in}\mathrm{\Omega }\hfill \\ h=h_{\mathrm{ex}}\hfill & \mathrm{on}\mathrm{\Omega },\hfill \end{array}$$ where the $`p_i`$’s are the centers of the vortices, and the $`d_i`$’s their degrees, that is the topological degree of the map $`u/\|u\|`$. These filaments are of characteristic size $`\epsilon `$. They are surrounded by a superconducting region in which $`\|u\|1`$. In order to minimize their repulsion, the flux lines form a triangular lattice, called the “Abrikosov lattice”. With increasing fields, the density of flux lines increase until the vortices overlap and $`H_{c_2}`$ is reached. The generation of vortices by the external field has been mathematically studied very recently in \[S1, S2, S3, SS1, SS2, SS3\]. In \[SS3\], it is proved among other things that, in the limit when $`\epsilon `$ tends to 0, equation (I.3) is replaced by (I.4) $$\mathrm{\Delta }h_{}+h_{}=\mu _{}$$ where $`\mu _{}`$ is the density of vortices in units of $`h_{\mathrm{ex}}`$ and $`h_{}=h/h_{\mathrm{ex}}`$. The measure $`\mu _{}`$ is supported in an inner region $`\omega `$ depending on the value of $`h_{\mathrm{ex}}`$ and is of uniform density in $`\omega `$. Our aim is to give a rigorous proof that in the small $`\epsilon `$ limit, stable configurations should correspond to vortices pinned at the minimum of $`a_\epsilon `$ and to derive the limiting homogenized free-boundary problem which arises for the magnetic field in replacement of the London equation (I.4). Using the second equation in (G.L.), we notice that the energy can be rewritten (I.5) $$J_\epsilon (u,A)=\frac{1}{2}_\mathrm{\Omega }\frac{1}{\|u\|^2}\|h\|^2+\|hh_{\mathrm{ex}}\|^2+\frac{1}{2}_\mathrm{\Omega }\|\|u\|\|^2+\frac{1}{2\epsilon ^2}(a_\epsilon (x)\|u\|^2)^2.$$ We will show that for a sequence of minimizers $`(u_\epsilon ,A_\epsilon )`$, the second integral in (I.5) is negligible. Then, when $`\epsilon `$ tends to 0, $`\|u\|^2a_\epsilon (x)`$ outside the vortices, and our main result will state that $`h_\epsilon =\mathrm{curl}A_\epsilon `$ satisfies roughly the following equivalent of (I.3) in the case of pinning: (I.6) $$\mathrm{div}\left(\frac{1}{a_\epsilon }h_\epsilon \right)+h_\epsilon =2\pi \underset{i}{}d_i\delta _{p_i}.$$ The existence of pinning will modify the locations $`p_i`$ of the vortices and the value of $`H_{c_1}`$. Since $`a_\epsilon `$ is a rapidly oscillating function describing impurities, the framework for passing to the limit when $`\epsilon `$ is small is that of homogenization theory. When passing to the limit in (I.6), we obtain a different limiting operator from (I.4), that is (I.7) $$\mathrm{div}\left(𝓐_\mathrm{𝟎}h_{}\right)+h_{}=\mu _{}$$ where $`\mu _{}`$ is a positive measure which is supported in an inner domain $`\omega _\mathrm{\Lambda }`$ and $`𝓐_\mathrm{𝟎}`$ is the homogenized limit of the matrix $`𝓐_𝜺={\displaystyle \frac{1}{a_\epsilon }}𝓘`$ in the sense of $`H`$-convergence, see definition below. ###### Definition 1 We say that the family of $`2\times 2`$ matrices $`𝓐_𝛆`$ $`H`$-converges to $`𝓐_\mathrm{𝟎}`$ when $`\epsilon `$ tends to 0, if and only if, for any $`f`$ in $`H^1(\mathrm{\Omega })`$, the solution $`v_\epsilon `$ in $`H_0^1(\mathrm{\Omega })`$ of $$\mathrm{div}(𝓐_𝜺v_\epsilon )+v_\epsilon =f$$ satisfies $$\begin{array}{cc}v_\epsilon v_0\text{weakly in}H_0^1(\mathrm{\Omega }),\hfill & \\ 𝓐_𝜺v_\epsilon 𝓐_\mathrm{𝟎}v_0\text{weakly in}\left(L^2(\mathrm{\Omega })\right)^2,\hfill & \end{array}$$ where $`v_0`$ is the $`H_0^1(\mathrm{\Omega })`$ solution of $$\mathrm{div}(𝓐_\mathrm{𝟎}v_0)+v_0=f.$$ We refer to the work of Murat and Tartar \[MuT\] for more details on the notion of $`H`$-convergence; one can also see \[CD, JKO\]. In the following, we will always let $`𝓐_𝜺={\displaystyle \frac{1}{a_\epsilon }}𝓘`$. Then $`𝓐_\mathrm{𝟎}`$ is also a diagonal matrix. In the general case, the computation of $`𝓐_\mathrm{𝟎}`$ is hard and not always known, see \[JKO\] for examples. But in some simple cases, this definition allows to compute $`𝓐_\mathrm{𝟎}`$. For instance, if $`a_\epsilon (x)=a(x/\epsilon )`$, and $`a(x)=a_1(x_1)a_2(x_2)`$ where $`a_1`$ and $`a_2`$ are periodic, then $$𝓐_\mathrm{𝟎}=diag(\frac{1}{a_1^0},\frac{1}{a_2^0}),\text{with}a_i^0=\overline{a_i}\overline{\left(\frac{1}{a_j}\right)}$$ where $`\overline{a_i}`$ denotes the mean of $`a_i`$ over a period (see \[JKO\]). Note that even though the sequence $`a_\epsilon `$ has no pointwise limit, the limiting problem and $`𝓐_\mathrm{𝟎}`$ are well defined. An important property of $`H`$-convergence (see \[MuT\]) is that if the sequence $`a_\epsilon `$ is bounded from below and above by positive constants independent of $`\epsilon `$, then there exists a subsequence $`𝓐_𝜺^{\mathbf{}}`$ and a matrix $`𝓐_\mathrm{𝟎}`$ for which $`𝓐_𝜺^{\mathbf{}}`$ $`H`$-converges to $`𝓐_\mathrm{𝟎}`$. For us, it will imply in the following that up to the extraction of a subsequence, the family $`𝓐_𝜺`$ $`H`$-converges to some limit $`𝓐_\mathrm{𝟎}`$, thus leading to the limiting problem (I.7). ### I.3 Main results Let us now state our hypotheses and results. We assume that $`h_{\mathrm{ex}}`$ is a function of $`\epsilon `$ and that the following limit exists and is finite: (I.8) $$\mathrm{\Lambda }=\underset{\epsilon 0}{lim}\frac{\|\mathrm{log}\epsilon \|}{h_{\mathrm{ex}}(\epsilon )}.$$ Moreover, we make the following hypotheses on the function $`a_\epsilon (x)`$: * There exists a constant $`b_0>0`$ such that $`b_0a_\epsilon (x)1`$. * There exist a constant $`C`$ and a sequence $`\eta (\epsilon )`$ (which may tend to 0 with $`\epsilon `$) such that $`1/\eta (\epsilon )h_{\mathrm{ex}}`$ and $`\|a_\epsilon \|{\displaystyle \frac{C}{\eta (\epsilon )}}`$. * There exist a continuous function $`b(x)`$ and a nonnegative functions $`\beta _\epsilon (x)`$ such that $`a_\epsilon (x)=b(x)+\beta _\epsilon (x)`$ and for any $`\epsilon >0`$ and any $`x\mathrm{\Omega }`$, $`\mathrm{min}_{B(x,\delta (\epsilon ))}\beta _\epsilon =0`$, where $$\delta (\epsilon )\frac{1}{(\mathrm{log}\|\mathrm{log}\epsilon \|)^{\frac{1}{2}}}.$$ * The family of matrices $`𝓐_𝜺`$ $`H`$-converges to $`𝓐_\mathrm{𝟎}`$. Note that, as we mentioned earlier, it follows from hypothesis (H1) and the compactness of the set of matrices bounded from above and below that there exists a subsequence of $`𝓐_𝜺`$ which $`H`$-converges to $`𝓐_\mathrm{𝟎}`$ \[MuT\]. Our hypothesis (H4) is there to restrict to this subsequence for ease of notation and to impose that the whole sequence converges. Moreover, (H2) means that $`a_\epsilon `$ can be a constant independent of $`\epsilon `$ but can also oscillate very quickly with $`\epsilon `$ (but not too quickly, i.e. not quicker than $`h_{\mathrm{ex}}`$). Note that in the case where $`a_\epsilon `$ does not depend on $`\epsilon `$, then $`𝓐_𝜺=𝓐_\mathrm{𝟎}`$ is constant. Let us emphasize that because $`\beta _\epsilon 0`$, $`b`$ can be thought of as the lower envelope of $`a_\epsilon `$ and the local minima of $`a_\epsilon `$ are the local minima of $`b`$. Hence $`b`$ will be related to the pinning sites of vortices and the oscillations of $`a_\epsilon `$ are those of $`\beta _\epsilon `$. Moreover, the hypotheses imply that $`bb_0`$. First, let us state the result concerning the limiting problem (I.7). We relate $`h_{}`$ and $`\mu _{}`$ to the minimum of a variational problem. Let $``$ denote the space of Radon measures in $`\mathrm{\Omega }`$. ###### Theorem 1 Let us assume that (H1) to (H4) are satisfied. Let us define for any $`\mathrm{\Lambda }0`$, (I.9) $$E(f)=\frac{\mathrm{\Lambda }}{2}_\mathrm{\Omega }b(x)\left\|\mathrm{div}(𝓐_\mathrm{𝟎}f)+f\right\|+\frac{1}{2}_\mathrm{\Omega }f𝓐_\mathrm{𝟎}f+\left\|f1\right\|^2,$$ over $$V=\{f\text{ s.t. }f1H_0^1(\mathrm{\Omega }),\text{ and }\mathrm{div}(𝓐_\mathrm{𝟎}f)+f\}.$$ The minimizer $`h_{}`$ of $`E`$ over $`V`$ exists and is unique. It satisfies $`(\mathrm{P})\{`$ $$\begin{array}{c}h_{}1H_0^1(\mathrm{\Omega })\hfill \\ \mu _{}=\mathrm{div}(𝓐_\mathrm{𝟎}h_{})+h_{}\hfill \\ h_{}1\frac{\mathrm{\Lambda }b}{2}\mathrm{in}\mathrm{\Omega }\hfill \\ \mu _{}\left(h_{}(1\frac{\mathrm{\Lambda }b}{2})\right)=0\mathrm{in}\mathrm{\Omega }.\hfill \end{array}$$ Moreover $`\mu _{}0`$ and $`\mu _{}H^1(\mathrm{\Omega })`$. Problem (P) is a free-boundary problem, called in the literature an “obstacle problem” (see \[R\]). Another way of considering problem (P) is to define the subset of $`\mathrm{\Omega }`$ (I.10) $$\omega _\mathrm{\Lambda }=\{x\mathrm{\Omega },\text{ s.t. }h_{}=1\mathrm{\Lambda }b/2\}.$$ Then $`\mu _{}=0`$ in $`\mathrm{\Omega }\overline{\omega _\mathrm{\Lambda }}`$, and $`h_{}=1\mathrm{\Lambda }b/2`$ in $`\omega _\mathrm{\Lambda }`$, $`\omega _\mathrm{\Lambda }`$ being called the “free-boundary”, because $`\omega _\mathrm{\Lambda }`$ is unknown and uniquely determined by the set of equations (P). Note that if $`𝓐_\mathrm{𝟎}`$ and $`b`$ are smooth enough then $`h_{}`$ is $`C^{1,\alpha }`$ ($`\alpha <1`$), $`\mu _{}`$ is in $`L^{\mathrm{}}`$, the free-boundary $`\omega _\mathrm{\Lambda }`$ is regular for almost every $`\mathrm{\Lambda }`$ (see \[BM\]) and then we can write $$\mu _{}=1\frac{\mathrm{\Lambda }b}{2}+\frac{\mathrm{\Lambda }}{2}\mathrm{div}(𝓐_\mathrm{𝟎}b)\text{in}\omega _\mathrm{\Lambda }.$$ Once we have proved Theorem 1 concerning the limiting problem, we can get convergence for any sequence of minimizers $`(u_\epsilon ,A_\epsilon )`$ of the energy $`J_\epsilon (u_\epsilon ,A_\epsilon )`$ to $`E(h_{})`$ in a sense similar to $`\mathrm{\Gamma }`$-convergence. ###### Theorem 2 Let us assume that (I.8) and (H1) to (H4) are satisfied. Let $`(u_\epsilon ,A_\epsilon )`$ be a family of minimizers of $`J_\epsilon `$, and $`h_\epsilon =\mathrm{curl}A_\epsilon `$ the associated magnetic field. Then, as $`\epsilon `$ tends to 0, $$\frac{h_\epsilon }{h_{\mathrm{ex}}}h_{}\text{weakly in}H^1(\mathrm{\Omega }),$$ where $`h_{}`$ is the minimizer of $`E`$. Moreover, (I.11) $`\underset{\epsilon 0}{lim}{\displaystyle \frac{J_\epsilon (u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}}=E(h_{})={\displaystyle \frac{\mathrm{\Lambda }}{2}}{\displaystyle _\mathrm{\Omega }}b\|\mu _{}\|+{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Omega }}h_{}𝓐_\mathrm{𝟎}h_{}+\|h_{}1\|^2,`$ (I.12) $`{\displaystyle \frac{\|h_\epsilon \|^2}{h_{\mathrm{ex}}^2a_\epsilon }}h_{}𝓐_\mathrm{𝟎}h_{}+\mathrm{\Lambda }b\mu _{},\text{in the sense of measures.}`$ One can easily notice that if $`\mathrm{\Lambda }=0`$ (i.e. if $`h_{\mathrm{ex}}\|\mathrm{log}\epsilon \|`$), the solution of (P) is $`h_{}=1`$, and $`E(h_{})=0`$. In this case, Theorem 2 asserts that $$\frac{h_\epsilon }{h_{\mathrm{ex}}}1\text{strongly in}H^1,\text{and}\underset{\epsilon 0}{lim}\frac{\mathrm{min}J_\epsilon }{h_{\mathrm{ex}}^2}=0.$$ The proof of Theorem 2 is the main part of the paper (see Section I.6 for a sketch). ### I.4 The case $`\mathrm{\Lambda }>0`$ Let us now present some stronger results in the case where $`\mathrm{\Lambda }`$ is positive, i.e. $`h_{\mathrm{ex}}`$ is of the order of $`\|\mathrm{log}\epsilon \|`$. The first issue is to determine mathematically the location of vortices. From the physics, we know that vortices are the zeroes of $`u_\epsilon `$ with non-zero winding number. Instead of defining vortices, we isolate them in disjoint vortex balls covering the set where $`\|u_\epsilon \|`$ is small. The centers of these balls can be thought of as being the centers of the vortices. ###### Proposition I.1 Let us assume that $`\mathrm{\Lambda }>0`$ and that (H1) to (H4) are satisfied, then there exists $`\epsilon _0`$ such that if $`\epsilon <\epsilon _0`$ and $`(u_\epsilon ,A_\epsilon )`$ is a minimizer of $`J_\epsilon `$, there exists a family of balls of disjoint closures (depending on $`\epsilon `$) $`(B_i)_{iI_\epsilon }=(B(p_i,r_i))_{iI_\epsilon }`$ satisfying : (I.13) $`\left\{x\mathrm{\Omega },\|\sqrt{a_\epsilon (x)}\|u_\epsilon (x)\|\|{\displaystyle \frac{1}{\|\mathrm{log}\epsilon \|}}\right\}{\displaystyle \underset{iI_\epsilon }{}}B(p_i,r_i).`$ (I.14) $`{\displaystyle \underset{iI_\epsilon }{}}r_i{\displaystyle \frac{1}{e^{\sqrt{\|\mathrm{log}\epsilon \|}}}}`$ (I.15) $`{\displaystyle \frac{1}{2}}{\displaystyle _{B_i}}{\displaystyle \frac{\|h_\epsilon \|^2}{\|u\|^2}}\pi b(p_i)\|d_i\|\|\mathrm{log}\epsilon \|(1o(1)),`$ where $`h_\epsilon =\mathrm{curl}A_\epsilon `$, and $`d_i=deg(\frac{u_\epsilon }{\|u_\epsilon \|},B_i)`$ if $`\overline{B_i}\mathrm{\Omega }`$, and 0 otherwise. This proposition will be proved at the beginning of Section II. Here is the meaning of the different inequalities: (I.13) locates the set where $`\|u_\epsilon \|`$ differs from $`a_\epsilon `$, which is contained in a union of disjoint balls; these balls represent the vortices or clusters of vortices. (I.14) gives a control on the size of the balls and (I.15) gives a lower bound on the energy, which is the contribution of vortices according to their degree $`d_i`$ and their location $`p_i`$, appearing through the value $`b(p_i)`$. As opposed to the case of $`a_\epsilon 1`$ (see \[SS3\]), the least energy is attained for $`p_i`$ at the minimum of $`b`$. Using this proposition, Theorem 1 can be made more precise: ###### Theorem 3 Let us assume that $`\mathrm{\Lambda }>0`$ and that (H1) to (H4) are satisfied. For any balls $`B(p_i,r_i)`$ and integers $`d_i`$ which satisfy (I.13)-(I.14)-(I.15), then (I.16) $`\underset{\epsilon 0}{lim}{\displaystyle \frac{2\pi }{h_{\mathrm{ex}}}}{\displaystyle \underset{iI_\epsilon }{}}d_ia_\epsilon (p_i)`$ $`={\displaystyle _\mathrm{\Omega }}b\|\mu _{}\|,`$ (I.17) $`{\displaystyle \frac{2\pi }{h_{\mathrm{ex}}}}{\displaystyle \underset{iI_\epsilon }{}}d_i\delta _{p_i}`$ $`\underset{\epsilon 0}{}\mu _{},`$ (I.18) $`{\displaystyle \frac{2\pi }{h_{\mathrm{ex}}}}{\displaystyle \underset{iI_\epsilon }{}}\|d_i\|\delta _{p_i}`$ $`\underset{\epsilon 0}{}\mu _{},`$ in the sense of measures, where $$\mu _{}=\mathrm{div}(𝓐_\mathrm{𝟎}h_{})+h_{}.$$ ### I.5 Physical interpretations and consequences Our results show that $`h_{}h_{\mathrm{ex}}`$ is a good approximation of $`h_\epsilon `$ and that, in the limit $`\epsilon 0`$, the vortices are scattered in an inner region $`\omega _\mathrm{\Lambda }`$ with density $`\mu _{}`$, where $`h_{}=1\mathrm{\Lambda }b(x)/2`$. In the outer region $`\mathrm{\Omega }\overline{\omega _\mathrm{\Lambda }}`$, there are no vortices and $`h_{}`$ satisfies $`\mathrm{div}(𝓐_\mathrm{𝟎}h_{})+h_{}=0`$. Unlike the case $`a_\epsilon 1`$, the vortex-density in $`\overline{\omega _\mathrm{\Lambda }}`$ is non-uniform in general. Moreover, as $`\mathrm{\Lambda }`$ decreases, the vortex-region first appears at the minimum of $`\psi `$ as defined by problem (I.19) below: as in \[SS3\], we can derive a necessary and sufficient condition for $`\omega _\mathrm{\Lambda }`$ to be nonempty. ###### Proposition I.2 Let $`\psi `$ be the solution of (I.19) $$\{\begin{array}{cccc}\hfill \mathrm{div}(𝓐_\mathrm{𝟎}\psi )+\psi & =1\hfill & & \mathrm{in}\mathrm{\Omega }\hfill \\ \hfill \psi & =0\hfill & & \mathrm{on}\mathrm{\Omega },\hfill \end{array}$$ then $$\omega _\mathrm{\Lambda }\mathrm{}\underset{\epsilon 0}{lim}\frac{h_{\mathrm{ex}}}{\|\mathrm{log}\epsilon \|}\frac{1}{2\mathrm{max}\|\psi \|}.$$ If we define $`H_{c_1}`$ as the field such that for $`h_{\mathrm{ex}}H_{c_1}`$, the minimizer of the energy has no vortex (i.e. $`\|u\|b_0/2`$) and for $`h_{\mathrm{ex}}H_{c_1}`$, there exists a minimizer with vortices; then Proposition I.2 gives a hint that $$H_{c_1}\frac{\|\mathrm{log}\epsilon \|}{2\mathrm{max}\|\psi \|}.$$ Thus the presence of pinning modifies the values of the first critical field (see \[S1, SS1\] for the case without pinning). In fact, we could adjust the proof of \[SS1\] to obtain: there exists $`k_\epsilon =O(\|\mathrm{log}\|\mathrm{log}\epsilon \|\|)`$ such that for $`\epsilon `$ small enough and $$h_{\mathrm{ex}}\frac{\|\mathrm{log}\epsilon \|}{2\mathrm{max}\|\psi \|}k_\epsilon $$ then any minimizer has no vortex. Furthermore, the position of the minimum of $`\psi `$ depends on the pinning potential $`a_\epsilon (x)`$. As $`\mathrm{\Lambda }`$ further decreases, corresponding to $`h_{\mathrm{ex}}`$ increasing, the vortex-region $`\omega _\mathrm{\Lambda }`$ grows, until, for $`\mathrm{\Lambda }=0`$ ($`h_{\mathrm{ex}}\|\mathrm{log}\epsilon \|`$), $`\omega _\mathrm{\Lambda }=\mathrm{\Omega }`$. At this point there are so many vortices that the macroscopic density of vortices and the induced magnetic field are no longer influenced by $`a_\epsilon `$. In other words, the strength of flux pinning is 0 for $`h_{\mathrm{ex}}\|\mathrm{log}\epsilon \|`$. In the case where $`a_\epsilon (x)=a(x)`$ is independent of $`\epsilon `$, $`a(x)=b(x)`$ and $`𝓐_\mathrm{𝟎}=a^1𝓘`$. Hence the limiting problem is a London equation with weight. We would like to point out that it is natural to define a vortex velocity by $`v=\frac{1}{\|u\|^2}h`$ (see \[CyP\]). In particular $$v_{}=\frac{1}{a}h_{}$$ can be defined as a limiting velocity (per unit of $`h_{\mathrm{ex}}`$). Note that in $`\omega _\mathrm{\Lambda }`$, since $`h_{}=1\frac{1}{2}\mathrm{\Lambda }a`$, then $`v_{}=\frac{1}{2}\mathrm{\Lambda }\mathrm{log}a`$. It implies that when $`a`$ is constant, $`v_{}=0`$ and there is no mean current in the vortex region. But when $`a`$ varies spatially, there is a nonzero limiting mean current and a nonzero limiting velocity $`v_{}`$. Hence $`vh_{\mathrm{ex}}v_{}`$ that is $`\frac{1}{2}\mathrm{log}\kappa \mathrm{log}a`$. This is the result of Chapman-Richardson \[CR\] in the case where the three-dimensional vortex line has no curvature. They describe the phenomenon saying that the variation in $`a`$ acts as a pinning potential. When $`\mathrm{\Lambda }=0`$, the velocity $`v_{}`$ is zero as well. Decreasing $`\mathrm{\Lambda }`$ means increasing the field. So when $`a`$ varies spatially, there is a critical exterior magnetic field above which the pinning potential has no role and the current is destroyed. In the general case where $`a_\epsilon `$ depends on $`\epsilon `$, it would be interesting to prove a convergence of the mean vortex velocity $`v_\epsilon =\frac{1}{\|u_\epsilon \|^2}h_\epsilon `$. Still, one can observe two different effects coming from the presence of pinning in the term $`\|h_\epsilon \|^2/a_\epsilon `$ and resulting in the energy $`E(h_{})`$ in the homogenization process: – One effect is related to the concentration of energy in the vortices and the location of the vortices. It appears through the term $$\frac{\mathrm{\Lambda }}{2}_\mathrm{\Omega }b\|\mu _{}\|$$ in the limiting energy $`E`$. This term is smaller if $`\mu _{}`$ is non-zero at points where $`b`$ is minimal. (I.16) implies that vortices go to points where $`\beta _\epsilon =0`$. These points will be called pinning sites in the following. Because $`\delta (\epsilon )`$ tends to 0, the number of such points is big. The effect on the position of vortices is to see $`b`$ and the minima of $`b`$. Moreover, since (I.17) and (I.18) have the same limit, it means that vortices tend to have positive degrees. If $`b`$ does not depend on $`x`$ then $`h_{}`$ and $`\mu _{}`$ are constant in $`\omega _\mathrm{\Lambda }`$, and there is no change for the location of vortices from the case $`a_\epsilon 1`$. On the other hand, if $`b`$ is non-uniform, then $`h_{}`$ is non-constant in $`\omega _\mathrm{\Lambda }`$ and there is a pinning current. If for example the domain is a disc and the minima of $`b`$, that is the impurities, are located at sites different from the center of the disc, one expects that vortices, or the vortex-region $`\omega _\mathrm{\Lambda }`$ will be closer to the minima of $`b`$, but it seems difficult to give a rigorous proof of this qualitative fact. – The other effect is due to the rapid oscillations of $`a_\epsilon `$ with $`\epsilon `$ and comes from the energy outside the vortices, converging to the homogenized term $$\frac{1}{2}_\mathrm{\Omega }h_{}𝓐_\mathrm{𝟎}h_{}+\|h_{}1\|^2$$ in $`E`$. It changes the equation for the magnetic field $`h`$ from the usual London equation. If $`\beta _\epsilon 0`$, then the homogenization effect can be anisotropic. The size $`\delta (\epsilon )`$ (which can be related to $`\eta `$ if $`\beta _\epsilon `$ is not identically 0) cannot be taken bigger than in (H3), otherwise each pinning site would be too large and the vortices could push one another outside the pinning site. Let us also point out that we cannot allow stronger oscillations of $`a_\epsilon `$ than in (H2), because the second integral in (I.5) would become the dominant term. It would be interesting to investigate what happens if (H2)-(H3) are relaxed. ### I.6 Main steps of the proof Let us now state the two steps of the proof of Theorem 2. It is obtained as in \[SS3\] by getting first a lower bound on the energy, Proposition I.3, proved in Section II, and then an upper bound, Proposition I.4, proved in Section III. ###### Proposition I.3 Let us assume that $`\mathrm{\Lambda }>0`$ and that (H1) to (H4) are satisfied. Let $`(u_\epsilon ,A_\epsilon )`$ be a minimizer of $`J_\epsilon `$. Then (I.20) $$\underset{\epsilon 0}{lim\; inf}\frac{1}{h_{\mathrm{ex}}^2}J_\epsilon (u_\epsilon ,A_\epsilon )\frac{\mathrm{\Lambda }}{2}_\mathrm{\Omega }b\|\mu _{}\|+\frac{1}{2}_\mathrm{\Omega }h_{}𝓐_\mathrm{𝟎}h_{}+\|h_{}1\|^2,$$ where $`h_{}`$ is the solution of $`(P)`$. ###### Proposition I.4 Let us assume that $`\mathrm{\Lambda }>0`$ and that (H1) to (H4) are satisfied. Let $`\mu `$ be a positive Radon measure, and let $`(u_\epsilon ,A_\epsilon )`$ be a minimizer of $`J_\epsilon `$. Then (I.21) $$\underset{\epsilon 0}{lim\; sup}\frac{1}{h_{\mathrm{ex}}^2}J_\epsilon (u_\epsilon ,A_\epsilon )\frac{\mathrm{\Lambda }}{2}_\mathrm{\Omega }b𝑑\mu +\frac{1}{2}_\mathrm{\Omega }h𝓐_\mathrm{𝟎}h+\|h1\|^2,$$ where $`h`$ is the solution of (I.22) $$\{\begin{array}{c}\hfill \mathrm{div}(𝓐_\mathrm{𝟎}h)+h=\mu \text{in}\mathrm{\Omega },\\ \hfill h=1\text{on}\mathrm{\Omega }.\end{array}$$ Section II is devoted to the proof of Proposition I.3. Let $`(u_\epsilon ,A_\epsilon )`$ be a sequence of minimizers and $`h_\epsilon =\mathrm{curl}A_\epsilon `$. The energy $`J_\epsilon (u_\epsilon ,A_\epsilon )`$ gives two contributions: inside the vortex balls and outside. Thus, first we prove Proposition I.1 where the vortex balls $`B_i`$ with centers $`p_i`$ are constructed and where the vortex energy is bounded from below. We define (I.23) $$\mu _\epsilon =\frac{2\pi }{h_{\mathrm{ex}}}\underset{iI_\epsilon }{}d_i\delta _{p_i}.$$ Then, Proposition I.1 implies (I.24) $$\frac{1}{h_{\mathrm{ex}}^2}_{_{iI}B_i}\frac{1}{\|u\|^2}\|h_\epsilon \|^2\frac{\|\mathrm{log}\epsilon \|}{h_{\mathrm{ex}}}_\mathrm{\Omega }b\|\mu _\epsilon \|,$$ which gives the lower bound inside the vortex balls. The next step is to pass to the limit in the energy outside the vortex balls. Letting $`h_0`$ be the weak $`H^1`$ limit of $`h_\epsilon /h_{\mathrm{ex}}`$, we obtain the following, which is similar to a standard result in homogenization theory (I.25) $$\underset{\epsilon 0}{lim\; inf}_{\mathrm{\Omega }\backslash _iB_i}\frac{\|h\|^2}{a_\epsilon h_{\mathrm{ex}}^2}_\mathrm{\Omega }h_0𝓐_\mathrm{𝟎}h_0.$$ This requires to introduce an auxiliary problem before applying the homogenization theory result and it works because the vortex balls are small and thus can be taken out of the first integral. Finally we derive from the Ginzburg-Landau equations the crucial fact that $`h_\epsilon `$ satisfies (I.26) $$\frac{1}{h_{\mathrm{ex}}}\left(\mathrm{div}\left(\frac{h_\epsilon }{a_\epsilon }\right)+h_\epsilon \right)=\mu _\epsilon +\psi _\epsilon $$ where $`\psi _\epsilon `$ tends to 0 and $`\mu _\epsilon `$ defined in (I.23) tends to some $`\mu _0`$, both convergences being strong in $`W^{1,r}`$ for $`r<2`$. The notion of $`H`$-convergence and a priori estimates allow us to pass to the limit in (I.26) in order to get that the weak $`H^1`$ limit of $`h_\epsilon /h_{\mathrm{ex}}`$, that we call $`h_0`$, solves (I.27) $$\mathrm{div}(𝓐_\mathrm{𝟎}h_0)+h_0=\mu _0.$$ Combining the lower bounds of the energy inside and outside the vortex balls (I.24)-(I.25), we find $$\underset{\epsilon 0}{lim\; inf}\frac{1}{h_{\mathrm{ex}}^2}J_\epsilon (u_\epsilon ,A_\epsilon )E(h_0)E(h_{}).$$ The last inequality is true because (I.27) implies that $`h_0`$ is in $`V`$. Section III is devoted to the proof of Proposition I.4. The proof holds for any positive Radon measure $`\mu `$. We apply it to $`\mu _{}`$ to get that $$\underset{\epsilon 0}{lim\; sup}\frac{1}{h_{\mathrm{ex}}^2}J_\epsilon (u_\epsilon ,A_\epsilon )E(h_{}),$$ which will imply the desired results of convergence. The upper bound of Proposition I.4 is obtained by constructing test configurations as follows. First, given a positive Radon measure $`\mu `$, we construct approximate measures $`\mu _\epsilon `$ which converge weakly to $`\mu `$: $$\mu _\epsilon =\frac{1}{h_{\mathrm{ex}}}\underset{i=1}{\overset{n_\epsilon }{}}\mu _\epsilon ^i,$$ where $`\mu _\epsilon ^i`$ is the line element on the circle $`B(p_\epsilon ^i,\epsilon )`$ normalized so that $`\mu _\epsilon ^i(B(p_\epsilon ^i,\epsilon ))=2\pi `$. The measure $`\mu _\epsilon `$ describes the vortices of our test-configuration. The difficulty is to choose the points $`p_\epsilon ^i`$ satisfying a number of properties. We tile $`\mathrm{\Omega }`$ with squares $`K`$ of size $`\delta (\epsilon )`$. In each square, there is at least a point $`p_K`$ where $`\beta _\epsilon =0`$. We choose $`n_𝒦`$ points $`p_\epsilon ^i`$ regularly scattered around $`p_K`$ in a ball of radius $`1/h_{\mathrm{ex}}`$. The number $`n_K`$ is chosen depending on $`\mu (K)`$ so that $`\mu _\epsilon `$ converge to $`\mu `$. Once the vortices are constructed, the rest follows easily: the magnetic field $`h_\epsilon `$ is defined to be the solution of (I.28) $$\frac{1}{h_{\mathrm{ex}}}\left(\mathrm{div}\left(\frac{h_\epsilon }{a_\epsilon }\right)+h_\epsilon \right)=\mu _\epsilon .$$ Then, we are the able to construct a configuration $`(u_\epsilon ,A_\epsilon )`$ such that $`\mathrm{curl}A_\epsilon =h_\epsilon `$ and $`u_\epsilon `$ has vortices at the points $`p_\epsilon ^i`$. Moreover, we obtain $$J_\epsilon (u_\epsilon ,A_\epsilon )\frac{1}{2}_\mathrm{\Omega }\frac{1}{a_\epsilon }\|h_\epsilon \|^2+\|h_\epsilon 1\|^2.$$ Finally we are able to show that $$\underset{\epsilon 0}{lim\; sup}\frac{1}{2h_{\mathrm{ex}}^2}_\mathrm{\Omega }\frac{1}{a_\epsilon }\|h_\epsilon \|^2+\|h_\epsilon 1\|^2\frac{\mathrm{\Lambda }}{2}_\mathrm{\Omega }b𝑑\mu +\frac{1}{2}_\mathrm{\Omega }h𝓐_\mathrm{𝟎}h+\|h1\|^2,$$ where $`h`$ solves $`\mathrm{div}(𝓐_\mathrm{𝟎}h)+h=\mu `$ and $`h=1`$ on $`\mathrm{\Omega }`$. ## II Lower bound In the following, we will denote $`_Au=uiAu`$. We will often drop the subscripts $`\epsilon `$. We consider $`(u_\epsilon ,A_\epsilon )`$ a family of minimizers of $`J_\epsilon `$, thus a family of solutions of (G.L.). We can state a few a priori bounds. Firstly, by the maximum principle, $`\|u_\epsilon \|\mathrm{max}a_\epsilon 1`$. Secondly, by minimality, comparing with $`(a_\epsilon ,0)`$, we get $$J_\epsilon (u_\epsilon ,A_\epsilon )J_\epsilon (a_\epsilon ,0).$$ But, by hypothesis (H2) on $`a_\epsilon `$, $$J_\epsilon (a_\epsilon ,0)=\frac{1}{2}_\mathrm{\Omega }\|a_\epsilon \|^2+O(h_{\mathrm{ex}}^2)\frac{C}{\eta ^2}+O(h_{\mathrm{ex}}^2)Ch_{\mathrm{ex}}^2.$$ Hence, we have the a-priori estimate (II.1) $$J_\epsilon (u_\epsilon ,A_\epsilon )Ch_{\mathrm{ex}}^2.$$ In addition, by applying a gauge-transformation to $`(u_\epsilon ,A_\epsilon )`$, we can choose the Coulomb-gauge $`\mathrm{div}A_\epsilon =0`$ in $`\mathrm{\Omega }`$, with $`A_\epsilon .n=0`$ on $`\mathrm{\Omega }`$. With this choice of gauge, we are easily lead (see \[S1, SS1\]) to the a priori bounds (II.2) $`A_\epsilon _{L^{\mathrm{}}(\mathrm{\Omega })}`$ $`Ch_{\mathrm{ex}}`$ (II.3) $`u_\epsilon _{L^2(\mathrm{\Omega })}`$ $`Ch_{\mathrm{ex}}.`$ We begin with the proof of Proposition I.1. ### II.1 Proof of Proposition I.1 \- Step 1 : Let $`(u,A)`$ be an energy-minimizer. Denoting $`\|u\|`$ by $`\rho `$, since $`_\mathrm{\Omega }\|u\|^2_\mathrm{\Omega }\|\rho \|^2`$, we deduce from (II.1) : (II.4) $$_\mathrm{\Omega }\|\rho \|^2+\frac{1}{2\epsilon ^2}(\rho ^2a_\epsilon )^2Ch_{\mathrm{ex}}^2.$$ But, $`{\displaystyle _\mathrm{\Omega }}\|\rho \|^2`$ $`=`$ $`{\displaystyle _\mathrm{\Omega }}\|(\rho \sqrt{a_\epsilon })\|^2+\|\sqrt{a_\epsilon }\|^22(\rho \sqrt{a_\epsilon })\sqrt{a_\epsilon }`$ $``$ $`{\displaystyle _\mathrm{\Omega }}\|(\rho \sqrt{a_\epsilon })\|^22\|(\rho \sqrt{a_\epsilon })\|\|\sqrt{a_\epsilon }\|.`$ Hence, in view of (II.4), $`{\displaystyle _\mathrm{\Omega }}\|(\rho \sqrt{a_\epsilon })\|^2`$ $``$ $`Ch_{\mathrm{ex}}^2+(\rho \sqrt{a_\epsilon })_{L^2}\sqrt{a_\epsilon }_{L^2}`$ $``$ $`Ch_{\mathrm{ex}}^2+{\displaystyle \frac{C}{\eta (\epsilon )}}(\rho \sqrt{a_\epsilon })_{L^2},`$ and, since $`\frac{1}{\eta (\epsilon )}h_{\mathrm{ex}}`$, $$_\mathrm{\Omega }\|(\rho \sqrt{a_\epsilon })\|^2\mathrm{max}(Ch_{\mathrm{ex}}^2,\frac{C}{\eta ^2})Ch_{\mathrm{ex}}^2.$$ In view of (II.4), we thus have (II.5) $$\frac{1}{2}_\mathrm{\Omega }\|(\rho \sqrt{a_\epsilon })\|^2+\frac{1}{2\epsilon ^2}(a_\epsilon \rho ^2)^2Ch_{\mathrm{ex}}^2C\|\mathrm{log}\epsilon \|^2.$$ \- Step 2 : For any $`t`$, let $`\mathrm{\Omega }_t=\{x\mathrm{\Omega }/\|\rho \sqrt{a_\epsilon }\|(x)>t\}`$ and $`\gamma _t=\mathrm{\Omega }_t.`$ Applying the coarea formula and arguing as in Lemma IV.2 of \[SS2\], $`C\|\mathrm{log}\epsilon \|^2{\displaystyle _\mathrm{\Omega }}\|(\rho \sqrt{a_\epsilon })\|^2+{\displaystyle \frac{1}{2\epsilon ^2}}(a_\epsilon \rho ^2)^2`$ $``$ $`{\displaystyle \frac{C}{\epsilon }}{\displaystyle _\mathrm{\Omega }}\|(\rho \sqrt{a_\epsilon })\|\|a_\epsilon \rho ^2\|`$ $``$ $`{\displaystyle \frac{C}{\epsilon }}{\displaystyle _0^+\mathrm{}}r(\gamma _t)t𝑑t.`$ Here, as in \[SS2\], $`r(\gamma _t)`$ is defined as the infimum over all finite coverings of $`\gamma _t`$ by balls $`B_1,\mathrm{},B_k`$ of the sum $`r_1+\mathrm{}+r_k`$ where $`r_i`$ is the radius of $`B_i`$. Combining the previous inequality with the mean-value theorem, we find that there exists a $`t[0,\frac{1}{\|\mathrm{log}\epsilon \|}]`$ such that $`r(\gamma _t)<C\epsilon \|\mathrm{log}\epsilon \|^3.`$ \- Step 3 : The next step is to construct the vortex-balls : starting from the chosen $`\gamma _t`$, covered by balls $`B_1,\mathrm{},B_k`$ (whose sum of the radii is controlled by $`C\epsilon \|\mathrm{log}\epsilon \|^3`$), we use the method of growing and merging of balls used in \[Sa, SS2\] : one needs to grow these balls $`B_i`$, keeping a suitable lower bound on the energy they contain, until the desired size is reached, with the desired lower bound. When some balls happen to intersect during the growth process, they are merged into a larger one. We refer the reader to \[SS2\], and here we only need to apply the result of Proposition IV.1 of \[SS2\] to $`A_\epsilon `$ and $`v=\frac{u}{\|u\|}=e^{i\phi }`$ in $`\mathrm{\Omega }\backslash \mathrm{\Omega }_t`$, $`\sigma =e^{\sqrt{\|\mathrm{log}\epsilon \|}}`$. We then obtain the existence of balls $`B_i=B(p_i,r_i)`$ such that (I.13) and (I.14) hold, and (II.6) $$\frac{1}{2}_{B_i\backslash \mathrm{\Omega }_t}\|\phi A\|^2+\frac{1}{2}_{B_i}\|hh_{\mathrm{ex}}\|^2\pi \|d_i\|\|\mathrm{log}\epsilon \|(1o(1)),$$ with $`d_i=\text{deg}(u,B_i)`$ if $`\overline{B_i}\mathrm{\Omega }`$, and $`0`$ otherwise. But we also have, from the Ginzburg-Landau equation $`^{}h=\rho ^2(\phi A)`$, and from $`\rho 1`$, $$_\mathrm{\Omega }\|h\|^2=_\mathrm{\Omega }\rho ^4\|\phi A\|^2_\mathrm{\Omega }\|_Au\|^2Ch_{\mathrm{ex}}^2,$$ hence $`{\displaystyle _{B_i}}\|hh_{\mathrm{ex}}\|^2`$ $``$ $`Cr_ihh_{\mathrm{ex}}_{L^4(\mathrm{\Omega })}^2Cr_ihh_{\mathrm{ex}}_{H^1(\mathrm{\Omega })}^2`$ $``$ $`Ch_{\mathrm{ex}}^2e^{\sqrt{\|\mathrm{log}\epsilon \|}}=o(1).`$ Thus, (II.6) becomes (II.7) $$\frac{1}{2}_{B_i\backslash \mathrm{\Omega }_t}\|\phi A\|^2\pi \|d_i\|\|\mathrm{log}\epsilon \|(1o(1)).$$ Now, $`{\displaystyle \frac{1}{2}}{\displaystyle _{B_i\backslash \mathrm{\Omega }_t}}\|_Au\|^2`$ $``$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{B_i\backslash \mathrm{\Omega }_t}}\rho ^2\|\phi A\|^2`$ $``$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{B_i\backslash \mathrm{\Omega }_t}}a_\epsilon \|\phi A\|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{B_i\backslash \mathrm{\Omega }_t}}(\rho ^2a_\epsilon )\|\phi A\|^2`$ $``$ $`{\displaystyle \frac{1}{2}}\left(\underset{B_i}{\mathrm{min}}a_\epsilon \right){\displaystyle _{B_i\backslash \mathrm{\Omega }_t}}\|\phi A\|^2{\displaystyle \frac{C}{\|\mathrm{log}\epsilon \|}}{\displaystyle _{B_i\backslash \mathrm{\Omega }_t}}\|\phi A\|^2,`$ where we have used (I.13). In view of (II.7), $$\frac{1}{2}_{B_i\backslash \mathrm{\Omega }_t}\|_Au\|^2\pi \left(\underset{B_i}{\mathrm{min}}a_\epsilon \right)\|d_i\|\|\mathrm{log}\epsilon \|(1o(1)).$$ So, using the hypotheses (H2) and (H3) on $`a_\epsilon `$, we are led to the two following lower bounds (II.8) $`{\displaystyle \frac{1}{2}}{\displaystyle _{B_i\backslash \mathrm{\Omega }_t}}\|_Au\|^2`$ $``$ $`\pi a_\epsilon (p_i)\|d_i\|\|\mathrm{log}\epsilon \|(1o(1))`$ (II.9) $`{\displaystyle \frac{1}{2}}{\displaystyle _{B_i\backslash \mathrm{\Omega }_t}}\|_Au\|^2`$ $``$ $`\pi b(p_i)\|d_i\|\|\mathrm{log}\epsilon \|(1o(1)).`$ This proves (I.15). $``$ ### II.2 Deriving the limiting equation For any $`(p_i,d_i)`$ satisfying (I.13)—(I.15), we can define (II.10) $$\mu _\epsilon =\frac{2\pi }{h_{\mathrm{ex}}}\underset{iI_\epsilon }{}d_i\delta _{p_i},$$ a measure of vorticity per unit of applied field. We will see that it remains a bounded family of measures. ###### Lemma II.1 If $`\mathrm{\Lambda }>0`$, and $`(u_\epsilon ,A_\epsilon )`$ is a family of minimizers of $`J_\epsilon `$ with $`h_\epsilon =\mathrm{curl}A_\epsilon `$, we can extract a sequence $`\epsilon _n0`$ such that there exists $`h_01H_0^1(\mathrm{\Omega })`$, and $`\mu _0`$ with $$\frac{h_{\epsilon _n}}{h_{\mathrm{ex}}}1h_01\text{ in }H_0^1(\mathrm{\Omega }),$$ $$\mu _{\epsilon _n}\mu _0\text{in the sense of measures}.$$ Proof : As seen in the previous proof, since $`(u_\epsilon ,A_\epsilon )`$ is a solution of the second Ginzburg-Landau equation $$_\mathrm{\Omega }\|h_\epsilon \|^2_\mathrm{\Omega }\|_{A_\epsilon }u_\epsilon \|^2Ch_{\mathrm{ex}}^2$$ and $$_\mathrm{\Omega }\|h_\epsilon h_{\mathrm{ex}}\|^2Ch_{\mathrm{ex}}^2.$$ Hence, $`\frac{h_\epsilon }{h_{\mathrm{ex}}}1`$ is bounded in $`H_0^1(\mathrm{\Omega })`$, and we can find a sequence $`\epsilon _n0`$ such that $`\frac{h_{\epsilon _n}}{h_{\mathrm{ex}}}`$ converges weakly in $`H_0^1`$ to some $`h_01`$. On the other hand, from Proposition I.1, $`Ch_{\mathrm{ex}}{\displaystyle \frac{\|\mathrm{log}\epsilon \|}{\mathrm{\Lambda }}}J_\epsilon (u_\epsilon ,A_\epsilon )`$ $``$ $`{\displaystyle \underset{iI_\epsilon }{}}\pi \|d_i\|b(p_i)\|\mathrm{log}\epsilon \|(1o(1))`$ $``$ $`b_0{\displaystyle \underset{i}{}}\pi \|d_i\|\|\mathrm{log}\epsilon \|(1o(1)),`$ where $`b_0`$ is given by hypothesis (H1) on $`a_\epsilon `$. Hence, $$\frac{1}{2}_\mathrm{\Omega }\|\mu _{\epsilon _n}\|=\frac{\pi _i\|d_i\|}{h_{\mathrm{ex}}}C,$$ thus $`(\mu _{\epsilon _n})`$ is a bounded sequence of measures, and extracting again if necessary, we can assume that $`\mu _{\epsilon _n}`$ converges to some $`\mu _0`$ in the sense of measures. $``$ ###### Proposition II.1 Let $`\mu _0`$ and $`h_0`$ be the measures and fields defined in Lemma II.1. Then there exists $`r_0<2`$ such that $`\mu _0W^{1,r}(\mathrm{\Omega })`$ $`r(r_0,2)`$, and $`h_0`$ is the unique solution in $`W^{1,r}`$ of (II.11) $$\{\begin{array}{cc}\mathrm{div}\left(𝓐_\mathrm{𝟎}h_0\right)+h_0=\mu _0\hfill & \text{ in }\mathrm{\Omega }\hfill \\ h_0=1\hfill & \text{ on }\mathrm{\Omega }.\hfill \end{array}$$ The proof of this proposition requires the following lemma, a slight refinement of the result stated in \[SS1\], Lemma II.3. ###### Lemma II.2 Under the hypotheses of Lemma II.1, for any $`q>2`$, $$\frac{1}{h_{\mathrm{ex}}}\mathrm{curl}\frac{(iu_\epsilon ,u_\epsilon )}{a_\epsilon }\mu _\epsilon \underset{\epsilon 0}{}0\text{ strongly in }(W_0^{1,q}(\mathrm{\Omega }))^{}.$$ Proof : Denote $`\stackrel{~}{\mathrm{\Omega }}=\mathrm{\Omega }\backslash _iB_i`$. On $`\stackrel{~}{\mathrm{\Omega }}`$, $`\|u_\epsilon \|b_0>0`$ and $`v_\epsilon =\frac{u_\epsilon }{\|u_\epsilon \|}`$ is well-defined. Let $`q>2`$, and $`\xi W_0^{1,q}`$. We need to show that $$\left\|\frac{1}{h_{\mathrm{ex}}}_\mathrm{\Omega }\xi \mathrm{curl}\frac{(iu_\epsilon ,u_\epsilon )}{a_\epsilon }\frac{2\pi }{h_{\mathrm{ex}}}\underset{i}{}d_i\xi (p_i)\right\|o(1)\xi _{W_0^{1,q}(\mathrm{\Omega })}.$$ Dropping again some of the subscripts, we have (II.12) $$\frac{1}{h_{\mathrm{ex}}}_\mathrm{\Omega }\xi \mathrm{curl}\frac{(iu,u)}{a_\epsilon }=\frac{1}{h_{\mathrm{ex}}}_\mathrm{\Omega }^{}\xi \frac{(iu,u)}{a_\epsilon }.$$ Then, the method consists in splitting this integral into the integral over the vortex-balls (which is going to be negligible because the balls are small enough) and the integral over $`\stackrel{~}{\mathrm{\Omega }}`$, the complement of the balls. \- Step 1 : We prove that (II.13) $$\left\|_{_iB_i}\frac{1}{h_{\mathrm{ex}}}^{}\xi \frac{(iu,u)}{a_\epsilon }\right\|=o(1)\xi _{L^q(\mathrm{\Omega })}.$$ Indeed, since $`a_\epsilon b_0>0`$, $$\left\|_{_iB_i}\frac{1}{h_{\mathrm{ex}}}^{}\xi \frac{(iu,u)}{a_\epsilon }\right\|\frac{1}{b_0}\frac{u_{L^2(\mathrm{\Omega })}}{h_{\mathrm{ex}}}\xi _{L^q}(\text{vol}(_iB_i))^{\frac{1}{p}},$$ where $`\frac{1}{p}+\frac{1}{q}=\frac{1}{2}`$ and we have used Hölder’s inequality twice. Using (II.3), $$\left\|_{_iB_i}\frac{1}{h_{\mathrm{ex}}}^{}\xi \frac{(iu,u)}{a_\epsilon }\right\|C(\underset{i}{}r_i^2)^{\frac{1}{p}}\xi _{L^q(\mathrm{\Omega })}.$$ In addition, $`(_ir_i^2)^{\frac{1}{p}}(_ir_i)^{\frac{2}{p}}=o(1)`$ since we know that $`_ir_i0`$. Therefore, (II.13) is proved. \- Step 2 : We observe that (II.14) $`{\displaystyle \frac{1}{h_{\mathrm{ex}}}}{\displaystyle _{\stackrel{~}{\mathrm{\Omega }}}}^{}\xi {\displaystyle \frac{(iu,u)}{a_\epsilon }}`$ $`=`$ $`{\displaystyle \frac{1}{h_{\mathrm{ex}}}}{\displaystyle _{\stackrel{~}{\mathrm{\Omega }}}}{\displaystyle \frac{\|u\|^2}{a_\epsilon }}(iv,v)^{}\xi `$ $`=`$ $`{\displaystyle \frac{1}{h_{\mathrm{ex}}}}{\displaystyle _{\stackrel{~}{\mathrm{\Omega }}}}(iv,v)^{}\xi +{\displaystyle \frac{1}{h_{\mathrm{ex}}}}{\displaystyle _{\stackrel{~}{\mathrm{\Omega }}}}\left({\displaystyle \frac{\|u\|^2}{a_\epsilon }}1\right)(iv,v)^{}\xi .`$ We claim that (II.15) $$\frac{1}{h_{\mathrm{ex}}}\left\|_{\stackrel{~}{\mathrm{\Omega }}}\left(\frac{\|u\|^2}{a_\epsilon }1\right)(iv,v)^{}\xi \right\|o(1)\xi _{L^q}.$$ Indeed, $`{\displaystyle \frac{1}{h_{\mathrm{ex}}}}\left\|{\displaystyle _{\stackrel{~}{\mathrm{\Omega }}}}\left({\displaystyle \frac{\|u\|^2}{a_\epsilon }}1\right)(iv,v)^{}\xi \right\|`$ $``$ $`{\displaystyle \frac{1}{b_0h_{\mathrm{ex}}}}\left\|{\displaystyle _{\stackrel{~}{\mathrm{\Omega }}}}(\|u\|^2a_\epsilon )\|v\|\|\xi \|\right\|`$ $``$ $`C{\displaystyle \frac{v_{L^2(\stackrel{~}{\mathrm{\Omega }})}}{h_{\mathrm{ex}}}}\xi _{L^q(\mathrm{\Omega })}\|u\|^2a_\epsilon _{L^p(\mathrm{\Omega })},`$ with $`\frac{1}{p}+\frac{1}{q}=\frac{1}{2}`$. From the a priori estimate (II.1), $$_\mathrm{\Omega }(\|u\|^2a_\epsilon )^pC_\mathrm{\Omega }(\|u\|^2a_\epsilon )^2C\epsilon ^2h_{\mathrm{ex}}^2=o(1),$$ hence, using $`v_{L^2(\stackrel{~}{\mathrm{\Omega }})}Cu_{L^2(\mathrm{\Omega })}Ch_{\mathrm{ex}}`$, we obtain (II.15). Combining (II.12)—(II.15), we have (II.16) $$\frac{1}{h_{\mathrm{ex}}}_\mathrm{\Omega }\mathrm{curl}\frac{(iu,u)}{a_\epsilon }\xi =\frac{1}{h_{\mathrm{ex}}}_{\stackrel{~}{\mathrm{\Omega }}}(iv,v)^{}\xi +o(1)\xi _{W_0^{1,q}}.$$ \- Step 3 : We evaluate $`{\displaystyle _{\stackrel{~}{\mathrm{\Omega }}}}(iv,v)^{}\xi .`$ Noticing that $`\mathrm{curl}(iv,v)0`$ on $`\stackrel{~}{\mathrm{\Omega }}`$, we have $$_{\stackrel{~}{\mathrm{\Omega }}}(iv,v)^{}\xi =_{\stackrel{~}{\mathrm{\Omega }}}\xi (iv,\frac{v}{\tau })=\underset{i}{}_{B_i\mathrm{\Omega }}\xi (iv,\frac{v}{\tau }).$$ There remains to prove that (II.17) $$\underset{i}{}_{B_i\mathrm{\Omega }}\xi (iv,\frac{v}{\tau })=2\pi \underset{i}{}d_i\xi (a_i)+o(h_{\mathrm{ex}})\xi _{W_0^{1,q}(\mathrm{\Omega })}.$$ Let $`f`$ be a $`C^1`$ function defined on $`_+`$ such that (II.18) $$\{\begin{array}{cc}f(x)=x\hfill & \text{ for }x\frac{b_0}{2}\hfill \\ f(x)=1\hfill & \text{ for }xb_0\hfill \\ \|f^{}(x)\|C\hfill & \text{ for any }x0.\hfill \end{array}$$ We can define the complex-valued function (II.19) $$w=f(\|u\|)v.$$ It has a meaning everywhere by setting $`w=u`$ where $`\|u\|\frac{b_0}{2}`$. Then, it is easy to check that (II.20) $$\|w\|C\|u\|\text{ in }\mathrm{\Omega },$$ and (II.21) $$\underset{i}{}_{B_i\mathrm{\Omega }}\xi (iv,\frac{v}{\tau })=\underset{i}{}_{B_i\mathrm{\Omega }}\xi (iw,\frac{w}{\tau }).$$ Using Stokes theorem, we have (II.22) $$\left\|\underset{i}{}_{B_i}(\xi \xi (p_i))(iw,\frac{w}{\tau })\right\|=\left\|\underset{i}{}_{B_i}^{}\xi (iw,w)+(\xi \xi (p_i))\mathrm{curl}(iw,w)\right\|.$$ But, on the one hand, (II.23) $`{\displaystyle \frac{1}{h_{\mathrm{ex}}}}\left\|{\displaystyle \underset{i}{}}{\displaystyle _{b_i}}^{}\xi (iw,w)\right\|`$ $``$ $`C{\displaystyle \frac{w_{L^2}}{h_{\mathrm{ex}}}}\xi _{L^q}\left({\displaystyle \underset{i}{}}\text{vol}(B_i)\right)^{\frac{1}{p}}`$ $``$ $`C{\displaystyle \frac{u_{L^2}}{h_{\mathrm{ex}}}}\xi _{L^q}\left({\displaystyle \underset{i}{}}r_i^2\right)^{\frac{1}{p}}`$ $``$ $`o(1)\xi _{L^q}`$ as in the proof of (II.13). On the other hand, using the fact that, since $`q>2`$, $`W_0^{1,q}`$ embeds in $`C^{0,\beta }`$ for some $`\beta <1`$, and $`\|\mathrm{curl}(iw,w)\|C\|w\|^2C\|u\|^2,`$ we have (II.24) $`\left\|{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{h_{\mathrm{ex}}}}{\displaystyle _{B_i}}(\xi \xi (p_i))\mathrm{curl}(iw,w)\right\|`$ $``$ $`(\underset{i}{\mathrm{max}}r_i)^\beta \xi _{C^{0,\beta }(\mathrm{\Omega })}{\displaystyle \underset{i}{}}{\displaystyle _{U_i}}{\displaystyle \frac{\|u\|^2}{h_{\mathrm{ex}}}}`$ $``$ $`e^{\beta \sqrt{\|\mathrm{log}\epsilon \|}}{\displaystyle \frac{u_{L^2(\mathrm{\Omega })}^2}{h_{\mathrm{ex}}}}\xi _{W_0^{1,q}}`$ $``$ $`h_{\mathrm{ex}}e^{\beta \sqrt{\|\mathrm{log}\epsilon \|}}\xi _{W_0^{1,q}}=o(1)\xi _{W_0^{1,q}},`$ since $`h_{\mathrm{ex}}C\|\mathrm{log}\epsilon \|`$. As in \[SS1\], the proof remains valid even if $`B_i`$ intersects $`\mathrm{\Omega }`$. Combining (II.23), (II.24), (II.21), and (II.22), (II.17) is proved. Consequently, in view of (II.16), we can conclude that $$\left\|\frac{1}{h_{\mathrm{ex}}}_\mathrm{\Omega }\xi \mathrm{curl}\frac{(iu,u)}{a_\epsilon }\frac{2\pi }{h_{\mathrm{ex}}}\underset{i}{}d_i\xi (p_i)\right\|o(1)\xi _{W_0^{1,q}},$$ hence that $`\frac{1}{h_{\mathrm{ex}}}\mathrm{curl}\frac{(iu,u)}{a_\epsilon }\mu _\epsilon 0`$ strongly in $`(W_0^{1,q})^{}`$ as stated. $``$ Proof of Proposition II.1 : For the sake of simplicity, we write $`\epsilon `$ instead of $`\epsilon _n`$. \- Step 1 : We prove that $`h_\epsilon `$ satisfies (II.25) $$\frac{1}{h_{\mathrm{ex}}}\left(\mathrm{div}\left(\frac{h_\epsilon }{a_\epsilon }\right)+h_\epsilon \right)=f_\epsilon ,$$ with $`f_\epsilon =\mu _\epsilon +\psi _\epsilon `$, where $`\psi _\epsilon 0`$ strongly in $`(W_0^{1,q})^{}`$ for $`q>2`$. Indeed, we start from the second Ginzburg-Landau equation : $$^{}h_\epsilon =(iu_\epsilon ,_{A_\epsilon }u_\epsilon ),$$ divide it by $`a_\epsilon `$ and take the curl : $$\mathrm{div}\left(\frac{h_\epsilon }{a_\epsilon }\right)=\mathrm{curl}\left(\frac{(iu_\epsilon ,u_\epsilon )}{a_\epsilon }A_\epsilon \frac{\|u_\epsilon \|^2}{a_\epsilon }\right),$$ hence (II.26) $$\mathrm{div}\left(\frac{h_\epsilon }{a_\epsilon }\right)+h_\epsilon =\mathrm{curl}\frac{(iu_\epsilon ,u_\epsilon )}{a_\epsilon }+\mathrm{curl}\left(A_\epsilon \left(1\frac{\|u_\epsilon \|^2}{a_\epsilon }\right)\right).$$ Now consider a test-function $`\xi W_0^{1,q}(\mathrm{\Omega })`$, $`q>2`$, $`\left\|{\displaystyle _\mathrm{\Omega }}\xi \mathrm{curl}\left(A_\epsilon \left(1{\displaystyle \frac{\|u\|^2}{a_\epsilon }}\right)\right)\right\|`$ $`=`$ $`\left\|{\displaystyle _\mathrm{\Omega }}^{}\xi A_\epsilon \left(1{\displaystyle \frac{\|u\|^2}{a_\epsilon }}\right)\right\|`$ $``$ $`CA_\epsilon _{L^{\mathrm{}}(\mathrm{\Omega })}\xi _{L^2(\mathrm{\Omega })}a_\epsilon \|u\|^2_{L^2(\mathrm{\Omega })}.`$ The a-priori bound (II.2), $`A_\epsilon _{L^{\mathrm{}}(\mathrm{\Omega })}O(h_{\mathrm{ex}})`$ and the energy bound, $`a_\epsilon \|u\|^2_{L^2}C\epsilon h_{\mathrm{ex}},`$ yield $$\left\|_\mathrm{\Omega }\xi \mathrm{curl}\left(A_\epsilon \left(1\frac{\|u\|^2}{a_\epsilon }\right)\right)\right\|o(1)\xi _{L^2}.$$ Consequently, $`\mathrm{curl}\left(A_\epsilon \left(1\frac{\|u\|^2}{a_\epsilon }\right)\right)0`$ strongly in $`(W_0^{1,q})^{}`$ for $`q>2`$. Combining this with (II.26) and Lemma II.2, we get the desired result. \- Step 2 : We prove that $`f_\epsilon `$ converges to $`\mu _0`$, the weak limit of $`\mu _\epsilon `$, in $`W^{1,r}(\mathrm{\Omega })`$ for any $`r<2`$. Indeed, from the upper bound on the energy, we know that $`\frac{1}{a_\epsilon h_{\mathrm{ex}}}h_\epsilon `$ is bounded in $`L^2(\mathrm{\Omega })`$, hence, in view of (II.25), $`f_\epsilon `$ is bounded in $`H^1`$, hence in $`W^{1,p}`$ for $`p<2`$. But, on the other hand, $`f_\epsilon =\mu _\epsilon +\psi _\epsilon `$, with $`\psi _\epsilon `$ bounded in $`W^{1,p}`$ for $`p<2`$, hence $`\mu _\epsilon `$ remains bounded in $`W^{1,p}`$ for $`p<2`$. Furthermore, $`\mu _\epsilon `$ is also bounded in the sense of measures, therefore we can apply a theorem of Murat (see \[Mu1\]) which asserts that such a $`\mu _\epsilon `$, bounded in the sense of measures and in $`W^{1,p}`$ for $`p<2`$, is necessarily compact in $`W^{1,r}`$ for $`r<p`$. Since this is also the case for $`\psi _\epsilon `$, which converges to zero, this implies that $`f_\epsilon `$ is compact in $`W^{1,r}`$ for $`r<2`$. In addition, its limit in the sense of distributions is $`\mu _0`$, hence it must converge to $`\mu _0`$ in $`W^{1,r}`$. \- Step 3 : We wish to pass to the limit in (II.25), but it is not possible directly because the $`H`$-convergence requires a right-hand side in $`H^1`$. So we are going to pass to the limit in the duality sense for a fixed right-hand side. Let $`gW^{1,q}`$ for $`q>2`$. Using the hypothesis (H1) on $`a_\epsilon `$, (which implies in particular the uniform ellipticity of $`\frac{1}{a_\epsilon }𝓘`$), we can apply a theorem of Meyers \[Me\] : there exists a $`q_0>2`$, such that if $`g`$ is in $`W^{1,q}`$ with $`2<qq_0`$, then equation (II.27) $$\{\begin{array}{cc}\mathrm{div}\left(\frac{v_\epsilon }{a_\epsilon }\right)+v_\epsilon =g\hfill & \text{ in }\mathrm{\Omega }\hfill \\ v_\epsilon =0\hfill & \text{ on }\mathrm{\Omega },\hfill \end{array}$$ has a unique solution $`v_\epsilon `$ in $`W_0^{1,q}`$. Thus, we have (II.28) $$_{W_0^{1,q^{}}}<\frac{h_\epsilon }{h_{\mathrm{ex}}}1,g>_{W^{1,q}}=_{W^{1,q^{}}}<f_\epsilon 1,v_\epsilon >_{W_0^{1,q}},$$ where $`\frac{1}{q^{}}+\frac{1}{q}=1`$, and we want to pass to the limit. More precisely, Meyers’ theorem yields that the operator $`R_\epsilon `$ which maps $`g`$ to $`v_\epsilon `$, is a bounded linear operator from $`W^{1,q}`$ to $`W_0^{1,q}`$ (for $`2<qq_0`$), hence up to extraction of a subsequence, $`v_\epsilon `$ has a weak limit $`v_0`$ in $`W_0^{1,q}`$. We assumed in hypothesis (H4) that $`\frac{1}{a_\epsilon }𝓘`$ $`H`$-converges to $`𝓐_\mathrm{𝟎}`$. By the definition of $`H`$-convergence (see \[MuT\]), and since $`W_0^{1,q}H_0^1`$, this implies that $`v_0`$ is the solution of (II.29) $$\{\begin{array}{cc}\mathrm{div}\left(𝓐_\mathrm{𝟎}v_0\right)+v_0=g\hfill & \text{ in }\mathrm{\Omega }\hfill \\ v_0=0\hfill & \text{ on }\mathrm{\Omega },\hfill \end{array}$$ Since this possible weak limit $`v_0`$ is unique, the whole sequence $`v_\epsilon `$ converges to $`v_0`$ weakly in $`W_0^{1,q}`$. In addition, $`f_\epsilon `$ converges strongly to $`\mu _0`$ in $`W^{1,q^{}}`$, thus we have $${}_{W^{1,q^{}}}{}^{}<f_\epsilon 1,v_\epsilon >_{W_0^{1,q}}<\mu _01,v_0>.$$ On the other hand, $`\frac{h_\epsilon }{h_{\mathrm{ex}}}1`$ converges weakly to $`h_01`$ in $`H_0^1`$. Thus, $${}_{W_0^{1,q^{}}}{}^{}<\frac{h_\epsilon }{h_{\mathrm{ex}}}1,g>_{W^{1,q}}<h_01,g>.$$ Therefore, we can pass to the limit in (II.28), and we are led to (II.30) $$_{W_0^{1,q^{}}}<h_01,g>_{W^{1,q}}=_{W^{1,q^{}}}<\mu _01,v_0>_{W_0^{1,q}}.$$ Meyers’ aforementioned theorem, also yields that for $`q_0^{}q^{}<2`$, (II.11) has a unique solution in $`W^{1,q^{}}`$. Since (II.30) holds for any $`g`$ in $`W^{1,q}`$, it implies that $`h_0`$ is this solution. $``$ ### II.3 Deriving a lower bound outside the vortex balls Next, we would like to deduce from (II.11) a lower bound like $$\underset{\epsilon 0}{lim\; inf}_{\mathrm{\Omega }\backslash _iB_i}\frac{\|h\|^2}{a_\epsilon h_{\mathrm{ex}}^2}_\mathrm{\Omega }h_0𝓐_\mathrm{𝟎}h_0.$$ But this is impossible to derive straightforwardly because the domain of integration in the left-hand side integral is not $`\mathrm{\Omega }`$. To remedy this, we replace $`h_\epsilon `$ by an auxiliary field $`\overline{h_\epsilon }`$, a sort of truncated of $`h_\epsilon `$ in the balls. This is a trick that was already used in \[SS2\] Proposition IV.1, Step 1. ###### Lemma II.3 There exists $`\overline{h_\epsilon }`$ such that $`\overline{h_\epsilon }1H_0^1(\mathrm{\Omega })`$ and 1) $`\frac{\overline{h_\epsilon }}{h_{\mathrm{ex}}}1h_01`$ in $`H_0^1(\mathrm{\Omega })`$, 2) $$_{\mathrm{\Omega }\backslash _iB_i}\frac{\|h\|^2}{a_\epsilon }+_\mathrm{\Omega }\|h_\epsilon h_{\mathrm{ex}}\|^2_\mathrm{\Omega }\frac{\|\overline{h_\epsilon }\|^2}{a_\epsilon }+\|\overline{h_\epsilon }h_{\mathrm{ex}}\|^2o(1),$$ 3) $$\underset{\epsilon 0}{lim\; inf}_\mathrm{\Omega }\frac{\|\overline{h_\epsilon }\|^2}{a_\epsilon }_\mathrm{\Omega }h_0𝓐_\mathrm{𝟎}h_0.$$ Proof : We consider $`\overline{A_\epsilon }`$ a solution of the following minimization problem : (II.31) $$\underset{AH^1(\mathrm{\Omega },^2),\text{div }A=0}{min}_{\mathrm{\Omega }\backslash _iB_i}a_\epsilon \|\phi A\|^2+_\mathrm{\Omega }\|\mathrm{curl}Ah_{\mathrm{ex}}\|^2,$$ where $`\phi `$ denotes the gradient of the phase of $`u_\epsilon `$ which is well-defined in $`\mathrm{\Omega }\backslash _iB_i`$. If we write $`\overline{h_\epsilon }=\mathrm{curl}\overline{A_\epsilon }`$, and we test (II.31) with $`h_\epsilon `$, we have (II.32) $$_{\mathrm{\Omega }\backslash _iB_i}a_\epsilon \|\phi \overline{A_\epsilon }\|^2+_\mathrm{\Omega }\|\overline{h_\epsilon }h_{\mathrm{ex}}\|^2_{\mathrm{\Omega }\backslash _iB_i}a_\epsilon \|\phi A_\epsilon \|^2+_\mathrm{\Omega }\|h_\epsilon h_{\mathrm{ex}}\|^2Ch_{\mathrm{ex}}^2.$$ In addition, $`\overline{h_\epsilon }`$ and $`\overline{A_\epsilon }`$ satisfy the following equations : (II.33) $$\{\begin{array}{cc}^{}\overline{h_\epsilon }=a_\epsilon (\phi \overline{A_\epsilon })\hfill & \text{ in }\mathrm{\Omega }\backslash _iB_i\hfill \\ \overline{h_\epsilon }=cst=c_i\hfill & \text{ on }B_i,i\hfill \\ \overline{h_\epsilon }=h_{\mathrm{ex}}\hfill & \text{ on }\mathrm{\Omega }.\hfill \end{array}$$ Thus, it satisfies (II.34) $$\mathrm{div}\left(\frac{\overline{h_\epsilon }}{a_\epsilon h_{\mathrm{ex}}}\right)+\frac{\overline{h_\epsilon }}{h_{\mathrm{ex}}}=\nu _\epsilon ,$$ where $`\nu _\epsilon `$ is the measure defined by (II.35) $$\xi W_0^{1,q}(\mathrm{\Omega }),(q>2),_\mathrm{\Omega }\nu _\epsilon \xi =\underset{i}{}\frac{1}{h_{\mathrm{ex}}}_{B_i}\xi \frac{\phi }{\tau }+\underset{i}{}\frac{1}{h_{\mathrm{ex}}}_{B_i}c_i\xi .$$ On the other hand, using Cauchy-Schwartz inequality, $$\left\|\frac{1}{h_{\mathrm{ex}}}\underset{i}{}_{B_i}c_i\xi \right\|=\left\|\frac{1}{h_{\mathrm{ex}}}_{_iB_i}\overline{h_\epsilon }\xi \right\|\xi _L^{\mathrm{}}\frac{\overline{h_\epsilon }}{h_{\mathrm{ex}}}_{L^2}\left(\underset{i}{}r_i\right)^{\frac{1}{2}}.$$ In view of (II.32), $`\frac{\overline{h_\epsilon }}{h_{\mathrm{ex}}}_{L^2}`$ is bounded, and $`(_ir_i)^{\frac{1}{2}}_ir_i0`$ from Proposition I.1. Hence, $$\left\|\frac{1}{h_{\mathrm{ex}}}\underset{i}{}_{B_i}c_i\xi \right\|=o(1)\xi _L^{\mathrm{}}.$$ On the other hand, the same proof as for Lemma II.2 shows that $$\left\|\underset{i}{}\frac{1}{h_{\mathrm{ex}}}_{B_i}\frac{\phi }{\tau }\xi _\mathrm{\Omega }\xi 𝑑\mu _\epsilon \right\|=o(1)\xi _{W_0^{1,q}}.$$ Hence, in view of (II.35), $`\nu _\epsilon \mu _\epsilon `$ converges strongly to $`0`$ in $`(W_0^{1,q})^{}`$. The same argument as in Proposition II.1 allows to conclude from (II.34) that $$\frac{\overline{h_\epsilon }}{h_{\mathrm{ex}}}1h_01\text{ in }H_0^1(\mathrm{\Omega }),$$ using the uniqueness of the solution of (II.11). Using (II.32) and (II.33), we get $`{\displaystyle _\mathrm{\Omega }}{\displaystyle \frac{\|\overline{h_\epsilon }\|^2}{a_\epsilon }}+\|\overline{h_\epsilon }h_{\mathrm{ex}}\|^2`$ $`=`$ $`{\displaystyle _{\mathrm{\Omega }\backslash _iB_i}}a_\epsilon \|\phi \overline{A_\epsilon }\|^2+{\displaystyle _\mathrm{\Omega }}\|\overline{h_\epsilon }h_{\mathrm{ex}}\|^2`$ $``$ $`{\displaystyle _{\mathrm{\Omega }\backslash _iB_i}}a_\epsilon \|\phi A_\epsilon \|^2+{\displaystyle _\mathrm{\Omega }}\|h_\epsilon h_{\mathrm{ex}}\|^2.`$ As in the proof of Proposition I.1, we have $$_{\mathrm{\Omega }\backslash _iB_i}a_\epsilon \|\phi A_\epsilon \|^2_{\mathrm{\Omega }\backslash _iB_i}\frac{\|h_\epsilon \|^2}{a_\epsilon }+o(1).$$ Thus, assertion 2) is proved. In addition, $`\frac{\overline{h_\epsilon }}{h_{\mathrm{ex}}}1`$ is bounded in $`H_0^1(\mathrm{\Omega })`$ and the convergence to $`h_01`$ is weak in $`H_0^1`$. There remains to prove the third assertion. But it is a classical result in homogenization theory (see \[JKO\]) that, since $`\frac{\overline{h_\epsilon }}{h_{\mathrm{ex}}}1h_01`$ in $`H_0^1(\mathrm{\Omega })`$ and $`\frac{1}{a_\epsilon }𝓘`$ $`H`$-converges to $`𝓐_\mathrm{𝟎}`$, $$\underset{\epsilon 0}{lim\; inf}_\mathrm{\Omega }\frac{1}{a_\epsilon }\left\|\left(\frac{\overline{h_\epsilon }}{h_{\mathrm{ex}}}\right)\right\|^2_\mathrm{\Omega }h_0𝓐_\mathrm{𝟎}h_0.$$ This completes the proof of the lemma. $``$ We recall that we defined $`E`$ in (I.9). ###### Lemma II.4 With the same notations, $$\underset{\epsilon 0}{lim\; inf}\frac{J_\epsilon (u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}\frac{\mathrm{\Lambda }}{2}_\mathrm{\Omega }b\|\mu _0\|+\frac{1}{2}_\mathrm{\Omega }h_0𝓐_\mathrm{𝟎}h_0+\|h_01\|^2=E(h_0).$$ Proof : The energy can easily be bounded from below as follows, splitting between the contribution inside the vortex-balls and the contribution outside : $`J_\epsilon (u_\epsilon ,A_\epsilon )`$ $``$ $`{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Omega }}\|_Au\|^2+\|hh_{\mathrm{ex}}\|^2`$ $``$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{_{iI}B_i}}\|_Au\|^2+{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Omega }\backslash _iB_i}}\rho ^2\|\phi A\|^2+{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Omega }}\|hh_{\mathrm{ex}}\|^2.`$ As previously, since for the energy-minimizers $`^{}h=(iu,_Au)`$, and $`\|\rho ^2a_\epsilon \|\frac{C}{\|\mathrm{log}\epsilon \|}`$ in $`\mathrm{\Omega }\backslash _iB_i`$, we have $$_{\mathrm{\Omega }\backslash _iB_i}\rho ^2\|\phi A\|^2=_{\mathrm{\Omega }\backslash _iB_i}\frac{\|h\|^2}{a_\epsilon }(1o(1)).$$ Therefore, in view of Proposition I.1, $$J_\epsilon (u_\epsilon ,A_\epsilon )\pi \underset{i}{}\|d_i\|b(p_i)\|\mathrm{log}\epsilon \|(1o(1))+_{\mathrm{\Omega }\backslash _iB_i}\frac{\|h\|^2}{a_\epsilon }(1o(1))+_\mathrm{\Omega }\|hh_{\mathrm{ex}}\|^2,$$ and with assertion 2) of Lemma II.3, $$\frac{J_\epsilon (u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}\frac{1}{2}\frac{\|\mathrm{log}\epsilon \|}{h_{\mathrm{ex}}}_\mathrm{\Omega }b\|\mu _\epsilon \|+\frac{1}{h_{\mathrm{ex}}^2}_\mathrm{\Omega }\frac{\|\overline{h_\epsilon }\|^2}{a_\epsilon }+_\mathrm{\Omega }\left\|\frac{\overline{h_\epsilon }}{h_{\mathrm{ex}}}1\right\|^2o(1).$$ We thus obtain, using assertion 3) of Lemma II.3 that (II.36) $$lim\; inf\frac{J_\epsilon (u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}lim\; inf\frac{1}{2}\left(\frac{\|\mathrm{log}\epsilon \|}{h_{\mathrm{ex}}}_\mathrm{\Omega }b\|\mu _\epsilon \|\right)+_\mathrm{\Omega }h_0𝓐_\mathrm{𝟎}h_0+\|h_01\|^2.$$ Similarly, using (II.8), we obtain (II.37) $$lim\; inf\frac{J_\epsilon (u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}lim\; inf\frac{1}{2}\left(\frac{\|\mathrm{log}\epsilon \|}{h_{\mathrm{ex}}}_\mathrm{\Omega }a_\epsilon \|\mu _\epsilon \|\right)+_\mathrm{\Omega }h_0𝓐_\mathrm{𝟎}h_0+\|h_01\|^2.$$ Then, using the weak convergence of $`\mu _\epsilon `$ to $`\mu _0`$ in $``$, and the weak lower semi-continuity of $`\mu _\mathrm{\Omega }b\|\mu \|`$, we conclude from (II.36) that $$lim\; inf\frac{J_\epsilon (u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}\frac{\mathrm{\Lambda }}{2}_\mathrm{\Omega }b\|\mu _0\|+_\mathrm{\Omega }h_0𝓐_\mathrm{𝟎}h_0+\|h_01\|^2=E(h_0).$$ $``$ The final convergence result will then follow from the combination of this result with the upper bound of Section III, leading to the fact that necessarily $`h_0`$ has to be $`h_{}`$, the minimizer of $`E`$, and $`\mu _0=\mu _{}`$. ## III Upper Bound In this section we prove Proposition I.4. First we remark that if $`h`$ is the solution of $`\mathrm{div}(𝒜h)+h=\mu `$ with boundary value $`1`$, then $$h(x)1=G(x,y)d(\mu 1)(y),$$ where $`G(.,y)`$ is the solution of $`\mathrm{div}(𝒜h)+h=\delta _y`$ vanishing on $`\mathrm{\Omega }`$ and $`\mu 1`$ denotes the difference between the measure $`\mu `$ and the Lebesgue measure in $`\mathrm{\Omega }`$. From this it follows easily that (III.1) $$_\mathrm{\Omega }h𝒜h+\|h1\|^2=G(x,y)d(\mu 1)(x)d(\mu 1)(y).$$ This last expression will be the one we use. To prove Proposition I.4 we will then need some properties of the Green functions $`G_\epsilon `$, $`G_0`$ associated to the operators $`\mathrm{div}(𝓐_𝜺u)+u`$ and $`\mathrm{div}(𝓐_\mathrm{𝟎}u)+u`$ respectively. These properties will be proved at the end of this section. ###### Lemma III.1 Let $`a_\epsilon =b+\beta _\epsilon `$ be a sequence of functions satisfying (H1) to (H4), and $`𝓐_\mathrm{𝟎}`$ be the homogenized limit of the matrices $`𝓐_𝛆=a_{\epsilon }^{}{}_{}{}^{1}𝓘`$ as $`\epsilon `$ goes to zero. For any $`y\mathrm{\Omega }`$, let $`G_\epsilon (.,y)`$ (resp. $`G_0(.,y)`$) be the solution of $`\mathrm{div}(𝓐_𝛆G_\epsilon )+G_\epsilon =\delta _y`$ (resp. $`\mathrm{div}(𝓐_\mathrm{𝟎}G_0)+G_0=\delta _y`$) that vanishes on $`\mathrm{\Omega }`$. The following properties hold: 1) $`G_\epsilon (x,y)`$, $`G_0(x,y)`$ are positive functions, and symmetric in $`x`$ and $`y`$. 2) $`\mathrm{\Delta }`$ denoting the diagonal in $`^2`$, there exists $`C>0`$ such that $`G_\epsilon (x,y)`$, $`G_0(x,y)`$ are bounded by $$C\left(\left\|\mathrm{log}\|xy\|\right\|+1\right)$$ for all $`x,y\overline{\mathrm{\Omega }}\times \overline{\mathrm{\Omega }}\mathrm{\Delta }`$. 3) For any compact $`K\mathrm{\Omega }`$, there exists $`C>0`$ such that for any $`x,y\mathrm{\Omega }`$ $$G_\epsilon (x,y)+\frac{a_\epsilon (x)}{2\pi }\mathrm{log}\|xy\|\frac{C}{\eta (\epsilon )},$$ where $`\eta (\epsilon )`$ is defined in (H3). 4) $`G_\epsilon `$ converges to $`G_0`$ locally uniformly in $`\overline{\mathrm{\Omega }}\times \overline{\mathrm{\Omega }}\mathrm{\Delta }`$. Then we have the following easy Lemma: ###### Lemma III.2 The function (III.2) $$I(\mu )=\frac{\mathrm{\Lambda }}{2}b𝑑\mu +\frac{1}{2}G_0(x,y)d(\mu 1)(x)d(\mu 1)(y)$$ is sequentially lower semicontinuous over the set of positive Radon measures supported in $`\overline{\mathrm{\Omega }}`$, with respect to weak-\* convergence. The proof of this can be found in \[W\] for instance. Note that $`I(.)`$ is well defined over the set of positive Radon measures if we admit the value $`+\mathrm{}`$. Note also that if we restrict to measures in $`H^1(\mathrm{\Omega })`$ then (III.1) shows that $`I(\mu )`$ is a lower semicontinous functional of $`h=L^1\mu `$ where $`L`$ is the operator $`u\mathrm{div}(𝒜u)+u`$ defined on $`H_0^1(\mathrm{\Omega })`$. It follows that $`I`$ is a lower semicontinuous function of $`\mu `$ with respect to $`H^1`$ convergence. Now the proof of Proposition I.4 splits into two propositions. ###### Proposition III.1 Assume that $`\mathrm{\Lambda }>0`$ and that (H1) to (H4) are satisfied. Let $`\mu `$ be a positive Radon measure with support in $`\overline{\mathrm{\Omega }}`$ and $`(p_\epsilon ^i)_{1in_\epsilon }`$ be families of points in $`\mathrm{\Omega }`$ such that $`ij`$ (III.3) $$\|p_\epsilon ^ip_\epsilon ^j\|>4\epsilon ,d(p_\epsilon ^i,\mathrm{\Omega })>\alpha _0>0,$$ where $`\alpha _0`$ is independent of $`\epsilon `$, (III.4) $$\frac{2\pi }{h_{\mathrm{ex}}}\underset{i=1}{\overset{n_\epsilon }{}}\delta _{p_\epsilon ^i}\mu ,\text{ in the sense of measures}$$ and (III.5) $$\underset{\epsilon 0}{lim}\left(\underset{\begin{array}{c}ij\\ \|p_\epsilon ^ip_\epsilon ^j\|<\alpha \end{array}}{}\frac{\left\|\mathrm{log}\|p_\epsilon ^ip_\epsilon ^j\|\right\|}{h_{\mathrm{ex}}^2}\right)\underset{\alpha 0}{\overset{}{}}0.$$ Then there exist configurations $`(v_\epsilon ,B_\epsilon )_{\epsilon >0}`$ such that (III.6) $$\underset{\epsilon 0}{lim\; sup}\frac{J_\epsilon (v_\epsilon ,B_\epsilon )}{h_{\mathrm{ex}}^2}\frac{\mathrm{\Lambda }}{2}\underset{\epsilon 0}{lim\; sup}\frac{2\pi {\displaystyle \underset{i=1}{\overset{n_\epsilon }{}}}a_\epsilon (p_\epsilon ^i)}{h_{\mathrm{ex}}}+\frac{1}{2}G_0d(\mu 1)d(\mu 1),$$ where $`G_0`$ is defined in Lemma III.1. This proposition states that under reasonable hypotheses on points $`p_\epsilon ^i`$, one can construct a good test configuration with prescribed vortices at $`p_\epsilon ^i`$. Moreover, (III.4) implies that $`n_\epsilon /h_{\mathrm{ex}}`$ is bounded. The following Proposition asserts that the construction of points $`p_\epsilon ^i`$ is possible. ###### Proposition III.2 Assume that $`\mathrm{\Lambda }>0`$ and that (H1) to (H4) are satisfied. Then given any positive Radon measure $`\mu `$ of the form $`\sigma (x)dx`$ where $`\sigma `$ is a positive continuous function compactly supported in $`\mathrm{\Omega }`$, there exist families of points $`(p_\epsilon ^i)_{1in_\epsilon }`$ satisfying (III.3), (III.4), (III.5) and such that (III.7) $$\underset{\epsilon 0}{lim\; sup}\frac{2\pi {\displaystyle \underset{i=1}{\overset{n_\epsilon }{}}}a_\epsilon (p_\epsilon ^i)}{h_{\mathrm{ex}}}_\mathrm{\Omega }b(x)𝑑\mu (x).$$ The proof of Proposition I.4 follows easily from these two Propositions. First, taking any positive Radon measure $`\mu `$ supported in $`\overline{\mathrm{\Omega }}`$, we may approach it in the weak-\* topology by measures $`\mu _n=\sigma _n(x)dx`$ where $`\sigma _nC_c(\mathrm{\Omega })`$ is a positive function. Applying Propositions III.1 and III.2, we may construct test-configurations $`(v_\epsilon ^n,B_\epsilon ^n)_{\epsilon >0}`$ such that $$\underset{\epsilon 0}{lim\; sup}\frac{J_\epsilon (v_\epsilon ^n,B_\epsilon ^n)}{h_{\mathrm{ex}}^2}\frac{\mathrm{\Lambda }}{2}b(x)𝑑\mu _n(x)+\frac{1}{2}G_0d(\mu _n1)d(\mu _n1).$$ Therefore the same inequality is satisfied if we replace $`(v_\epsilon ^n,B_\epsilon ^n)`$ by the minimizing configuration $`(u_\epsilon ,A_\epsilon )`$. This proves that for each $`n`$, $$\underset{\epsilon 0}{lim\; sup}\frac{J_\epsilon (u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}I(\mu _n),$$ and then, using Lemma III.2, (III.8) $$\underset{\epsilon 0}{lim\; sup}\frac{J_\epsilon (u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}\frac{\mathrm{\Lambda }}{2}_\mathrm{\Omega }b𝑑\mu +\frac{1}{2}G_0(x,y)d(\mu 1)(x)d(\mu 1)(y).$$ Using (III.1) we get the conclusion of Proposition I.4. ### III.1 Proof of Proposition III.1 The method for constructing a test configuration $`(v_\epsilon ,B_\epsilon )`$ with prescribed vortices $`(p_\epsilon ^i)_{1in_\epsilon }`$ follows closely that of \[SS3\]. First we define $`h_\epsilon `$ to be the solution of (III.9) $$\{\begin{array}{cccc}\hfill \mathrm{div}(𝓐_𝜺h_\epsilon )+h_\epsilon & =\underset{i=1}{\overset{n_\epsilon }{}}\mu _\epsilon ^i\hfill & & \mathrm{in}\mathrm{\Omega }\hfill \\ \hfill h_\epsilon & =h_{\mathrm{ex}}\hfill & & \mathrm{on}\mathrm{\Omega },\hfill \end{array}$$ where $`\mu _\epsilon ^i`$ is the line element on the circle $`B(p_\epsilon ^i,\epsilon )`$ normalized so that $`\mu _\epsilon ^i(B(p_\epsilon ^i,\epsilon ))=2\pi `$. Then we let $`B_\epsilon `$ be any vector field such that $`\mathrm{curl}B_\epsilon =h_\epsilon `$. Finally, we define $`v_\epsilon =\rho _\epsilon e^{i\phi _\epsilon }`$ as follows: first we let (III.10) $$\rho _\epsilon (x)=\{\begin{array}{cc}\hfill 0& \text{if }\|xp_i^\epsilon \|\epsilon \text{ for some }i\text{,}\hfill \\ \hfill \sqrt{a_\epsilon (x)}\frac{\|xp_i^\epsilon \|\epsilon }{\epsilon }& \text{if }\epsilon <\|xa_i^\epsilon \|<2\epsilon \text{ for some }i\text{,}\hfill \\ \hfill \sqrt{a_\epsilon (x)}& \text{otherwise,}\hfill \end{array}$$ and for any $`x\mathrm{\Omega }_\epsilon =\mathrm{\Omega }_iB(p_i^\epsilon ,\epsilon )`$, (III.11) $$\phi _\epsilon (x)=_{(x_0,x)}(B_\epsilon 𝓐_𝜺^{}h_\epsilon ).\tau d\mathrm{},$$ where $`x_0`$ is a base point in $`\mathrm{\Omega }_\epsilon `$, $`(x_0,x)`$ is any curve joining $`x_0`$ to $`x`$ in $`\mathrm{\Omega }_\epsilon `$ and $`\tau `$ is the tangent vector to the curve. From (III.9), we see that this definition of $`\phi _\epsilon (x)`$ does not depend modulo $`2\pi `$ on the particular curve $`(x_0,x)`$ chosen. The fact that $`\phi _\epsilon `$ is not defined on $`_iB(p_i^\epsilon ,\epsilon )`$ is not important since $`\rho _\epsilon `$ is zero there. Thus, $`\phi _\epsilon `$ satisfies (III.12) $$𝓐_𝜺^{}h_\epsilon =\phi _\epsilon B_\epsilon $$ in $`\mathrm{\Omega }_\epsilon `$. Having defined $`v_\epsilon =\rho _\epsilon e^{i\phi _\epsilon }`$, we estimate $`J_\epsilon (v_\epsilon ,B_\epsilon )`$. Recall that (III.13) $$J_\epsilon (v_\epsilon ,B_\epsilon )=\frac{1}{2}_\mathrm{\Omega }\|\rho _\epsilon \|^2+\rho _{\epsilon }^{}{}_{}{}^{2}\|\phi _\epsilon B_\epsilon \|^2+\|h_\epsilon h_{\mathrm{ex}}\|^2+\frac{1}{2\epsilon ^2}\left(a_\epsilon \phi _{\epsilon }^{}{}_{}{}^{2}\right)^2.$$ Using the fact that $`\|a_\epsilon \|h_{\mathrm{ex}}`$ (hypothesis (H2)) and that the number of points $`p_\epsilon ^i`$ is less than $`Ch_{\mathrm{ex}}`$ — which follows from (III.4) — it is not difficult to check that (III.14) $$\frac{1}{2}_\mathrm{\Omega }\|\rho _\epsilon \|^2+\frac{1}{2\epsilon ^2}\left(a_\epsilon \rho _{\epsilon }^{}{}_{}{}^{2}\right)^2h_{\mathrm{ex}}^2.$$ Also, from (III.10), (III.12), $$\rho _{\epsilon }^{}{}_{}{}^{2}\|\phi _\epsilon B_\epsilon \|^2a_\epsilon \|\phi _\epsilon B_\epsilon \|^2=h_\epsilon 𝓐_𝜺h_\epsilon $$ in $`\mathrm{\Omega }_\epsilon `$. Therefore, replacing in (III.13) and in view of (III.14) (III.15) $$\underset{\epsilon 0}{lim\; sup}\frac{J_\epsilon (v_\epsilon ,B_\epsilon )}{h_{\mathrm{ex}}^2}\underset{\epsilon 0}{lim\; sup}\frac{1}{2h_{\mathrm{ex}}^2}_\mathrm{\Omega }h_\epsilon 𝓐_𝜺h_\epsilon +\|h_\epsilon h_{\mathrm{ex}}\|^2.$$ Because $`h_\epsilon `$ is the solution of (III.9), we may rewrite the right-hand side of this inequality as $$\underset{\epsilon 0}{lim\; sup}\frac{1}{2}G_\epsilon (x,y)d(\mu _\epsilon 1)(x)d(\mu _\epsilon 1)(y),$$ where (III.16) $$\mu _\epsilon =\frac{1}{h_{\mathrm{ex}}}\underset{i=1}{\overset{n_\epsilon }{}}\mu _\epsilon ^i,$$ and $`\mu _\epsilon ^i`$ is defined in (III.9). It follows from (III.4), (III.9) and (III.16) that $`\mu _\epsilon \mu `$ as $`\epsilon 0`$. Thus, to finish the proof of the proposition, it remains to show that (III.17) $$\begin{array}{c}\underset{\epsilon 0}{lim\; sup}\frac{1}{2}G_\epsilon d(\mu _\epsilon 1)d(\mu _\epsilon 1)\frac{\mathrm{\Lambda }}{2}\underset{\epsilon 0}{lim\; sup}\frac{2\pi {\displaystyle \underset{i=1}{\overset{n_\epsilon }{}}}a_\epsilon (p_\epsilon ^i)}{h_{\mathrm{ex}}}\hfill \\ \hfill +\frac{1}{2}G_0d(\mu 1)d(\mu 1)\end{array}$$ #### Proof of (III.17) Let $`\alpha >0`$ and let $`\mathrm{\Delta }_\alpha =\{(x,y)\|xy\|<\alpha \}`$. Recall that $`\mu _\epsilon \mu `$. Hence, it follows that $`(\mu _\epsilon 1)(\mu _\epsilon 1)(\mu 1)(\mu 1)`$ as $`\epsilon 0`$. But from Lemma II.1, $`G_\epsilon `$ tends to $`G_0`$ uniformly in $`\overline{\mathrm{\Omega }}\times \overline{\mathrm{\Omega }}\mathrm{\Delta }_\alpha `$, therefore (III.18) $$\underset{\epsilon 0}{lim}\frac{1}{2}_{\overline{\mathrm{\Omega }}\times \overline{\mathrm{\Omega }}\mathrm{\Delta }_\alpha }G_\epsilon d(\mu _\epsilon 1)d(\mu _\epsilon 1)=\frac{1}{2}_{\overline{\mathrm{\Omega }}\times \overline{\mathrm{\Omega }}\mathrm{\Delta }_\alpha }G_0d(\mu 1)d(\mu 1).$$ Now we treat the integral on $`\mathrm{\Delta }_\alpha `$. More precisely we prove that (III.19) $$\underset{\epsilon 0}{lim\; sup}_{\mathrm{\Delta }_\alpha }G_\epsilon d(\mu _\epsilon 1)d(\mu _\epsilon 1)\frac{\mathrm{\Lambda }}{2}\underset{\epsilon 0}{lim\; sup}\frac{2\pi {\displaystyle \underset{i=1}{\overset{n_\epsilon }{}}}a_\epsilon (p_\epsilon ^i)}{h_{\mathrm{ex}}}+o_\alpha (1),$$ where $`lim_{\alpha 0}o_\alpha (1)=0`$. Adding (III.18), (III.19) and letting $`\alpha 0`$ yields (III.17). We are left with proving (III.19). First we use the bound $`\|G_\epsilon (x,y)\|<C\left\|\mathrm{log}\|xy\|\right\|`$ from which one easily gets $$_{\mathrm{\Delta }_\alpha }G_\epsilon d(\mu _\epsilon 1)d(\mu _\epsilon 1)_{\mathrm{\Delta }_\alpha }G_\epsilon 𝑑\mu _\epsilon 𝑑\mu _\epsilon +C\alpha ^2\|\mathrm{log}\alpha \|.$$ Therefore (III.19) will follow if we prove (III.20) $$\underset{\epsilon 0}{lim\; sup}_{\mathrm{\Delta }_\alpha }G_\epsilon 𝑑\mu _\epsilon 𝑑\mu _\epsilon \frac{\mathrm{\Lambda }}{2}\underset{\epsilon 0}{lim\; sup}\frac{2\pi {\displaystyle \underset{i=1}{\overset{n_\epsilon }{}}}a_\epsilon (p_\epsilon ^i)}{h_{\mathrm{ex}}}+o_\alpha (1).$$ To prove this, we come back to the definition of $`\mu _\epsilon `$. From this definition, we have (III.21) $$_{\mathrm{\Delta }_\alpha }G_\epsilon 𝑑\mu _\epsilon 𝑑\mu _\epsilon \frac{1}{h_{\mathrm{ex}}^2}\left(\underset{\begin{array}{c}1ijn_\epsilon \\ \|p_\epsilon ^ip_\epsilon ^j\|<2\alpha \end{array}}{}G_\epsilon 𝑑\mu _\epsilon ^i𝑑\mu _\epsilon ^j+\underset{i=1}{\overset{n_\epsilon }{}}G_\epsilon 𝑑\mu _\epsilon ^i𝑑\mu _\epsilon ^i\right).$$ Let us first estimate the first sum on the right-hand side. If $`x\text{Supp}\mu _\epsilon ^i=B(p_\epsilon ^i,\epsilon )`$, $`y\text{Supp}\mu _\epsilon ^j`$ and $`ij`$, since $`\|p_\epsilon ^ip_\epsilon ^j\|>4\epsilon `$, then $`\|xy\|>\frac{1}{2}\|p_\epsilon ^ip_\epsilon ^j\|`$. Using the bound $`\|G_\epsilon (x,y)\|<C\left\|\mathrm{log}\|xy\|\right\|`$ together with the fact that $`\|p_\epsilon ^ip_\epsilon ^j\|<2\alpha `$ and $`\alpha `$ is small enough, we get $$G_\epsilon 𝑑\mu _\epsilon ^i𝑑\mu _\epsilon ^j<C\left\|\mathrm{log}\|p_\epsilon ^ip_\epsilon ^j\|\right\|.$$ Then, by hypothesis (III.5), (III.22) $$\underset{\epsilon 0}{lim\; sup}\frac{1}{h_{\mathrm{ex}}^2}\underset{\begin{array}{c}1ijn_\epsilon \\ \|p_\epsilon ^ip_\epsilon ^j\|<2\alpha \end{array}}{}G_\epsilon 𝑑\mu _\epsilon ^i𝑑\mu _\epsilon ^jo_\alpha (1).$$ As for the second sum in the right-hand side of (III.21), we use property 3) in Lemma III.1 to get that for any $`1in_\epsilon `$, and any $`x,y\text{Supp}\mu _\epsilon ^i`$, (III.23) $$G_\epsilon (x_,y)+\frac{a_\epsilon (x)}{2\pi }\mathrm{log}\|xy\|<\frac{C}{\eta (\epsilon )}\|\mathrm{log}\epsilon \|.$$ But $`x\text{Supp}\mu _\epsilon ^i`$ is equivalent to $`\|xp_\epsilon ^i\|=\epsilon `$. Then property (H2) of $`a_\epsilon `$ implies that $`a_\epsilon (x)a_\epsilon (p_\epsilon ^i)`$ as $`\epsilon 0`$. Replacing in (III.23) and integrating w.r.t. $`\mu _\epsilon ^i\mu _\epsilon ^i`$ yields $$G_\epsilon 𝑑\mu _\epsilon ^i𝑑\mu _\epsilon ^i2\pi a_\epsilon (p_\epsilon ^i)\|\mathrm{log}\epsilon \|\left(1+o_\epsilon (1)\right)$$ and then, summing over $`1in_\epsilon `$ and dividing by $`h_{\mathrm{ex}}`$, (III.24) $$\underset{\epsilon 0}{lim\; sup}\frac{1}{h_{\mathrm{ex}}^2}\underset{i=1}{\overset{n_\epsilon }{}}G_\epsilon 𝑑\mu _\epsilon ^i𝑑\mu _\epsilon ^i\frac{\mathrm{\Lambda }}{2}\underset{\epsilon 0}{lim\; sup}\frac{2\pi {\displaystyle \underset{i=1}{\overset{n_\epsilon }{}}}a_\epsilon (p_\epsilon ^i)}{h_{\mathrm{ex}}}.$$ Here we have used the fact that $`\|\mathrm{log}\epsilon \|\mathrm{\Lambda }h_{\mathrm{ex}}`$. Thus (III.20) is proved and the Proposition follows. $``$ ### III.2 Proof of Proposition III.2 Let $`\mu =\sigma (x)dx`$, $`C=u_{\mathrm{}}`$ and $`\alpha _0=dist(\text{supp}\mu ,\mathrm{\Omega })`$. Also, let (III.25) $$\stackrel{~}{\mathrm{\Omega }}=\{x\mathrm{\Omega }d(x,\mathrm{\Omega })>\alpha _0/2\}.$$ Recall that from hypothesis (H3) on $`a_\epsilon `$ there exists a positive function $`\delta (\epsilon )`$ such that (III.26) $$\delta (\epsilon )\frac{1}{(\mathrm{log}\|\mathrm{log}\epsilon \|)^{\frac{1}{2}}},\text{and for any }x\mathrm{\Omega }\text{}\mathrm{min}_{B(x,\delta (\epsilon ))}\beta _\epsilon =0.$$ For any $`\epsilon >0`$, we tile $`^2`$ with open squares of sidelength $`2\delta (\epsilon )`$ and let $`𝒦(\epsilon )`$ be the family of those squares that are entirely inside $`\stackrel{~}{\mathrm{\Omega }}`$. We denote by $`c_K`$ the center of a square $`K`$. Since $`\mu `$ is absolutely continuous with respect to the Lebesgue measure, we have $`\mu (K)C\delta ^2`$. Now the family of points $`(p_\epsilon ^i)_{1in_\epsilon }`$ is defined as follows: for any $`K𝒦(\epsilon )`$, we let (III.27) $$n(K,\epsilon )=\left[\frac{h_{\mathrm{ex}}(\epsilon )\mu (K)}{2\pi }\right],$$ where $`[x]`$ is the biggest integer no greater than $`x`$. Using (III.26) there is a point $`p_KB(c_K,\delta )`$ such that $`\beta _\epsilon (p_K)=0`$ ($`p_K`$ is a pinning site). We now pick $`n(K,\epsilon )`$ points evenly scattered in the ball $`B(p_K,1/h_{\mathrm{ex}})`$, and we call $`𝒫(K,\epsilon )`$ their union. By evenly scattered we mean that for any $`p,q𝒫(K,\epsilon )`$, (III.28) $$\|pq\|\frac{C}{h_{\mathrm{ex}}\sqrt{n(K,\epsilon )}}.$$ We let (III.29) $$n_\epsilon =\underset{K𝒦(\epsilon )}{}n(K,\epsilon ),\text{and }𝒫(\epsilon )=_{K𝒦(\epsilon )}𝒫(K,\epsilon )=(p_\epsilon ^i)_{1in_\epsilon }$$ be our family of points. We now check that this family satisfies (III.3), (III.4), (III.5) and (III.7). (III.3) is clear from (III.28) if $`p_\epsilon ^i,p_\epsilon ^j`$ belong to the same pinning site. It is even more true if $`p_\epsilon ^i,p_\epsilon ^j`$ do not belong to the same site since in this case their mutual distance is at least $`2\delta (\epsilon )\epsilon `$. Moreover from (III.25) we have $`d(p_\epsilon ^i,\mathrm{\Omega })>\alpha _0/2`$. For (III.4), let (III.30) $$\mu _\epsilon =\frac{2\pi }{h_{\mathrm{ex}}}\underset{i=1}{\overset{n_\epsilon }{}}\delta _{p_\epsilon ^i}$$ and $`f`$ be a continuous function in $`\overline{\mathrm{\Omega }}`$. We let $`\gamma _\epsilon =sup_{K𝒦(\epsilon )}sup_{x,yK}\|f(x)f(y)\|`$. Then since the size of the squares in $`𝒦(\epsilon )`$ tends to zero with $`\epsilon `$, so does $`\gamma _\epsilon `$. Let $`K_\epsilon `$ be the union of the squares in $`𝒦(\epsilon )`$, then for $`\epsilon `$ small enough $`\text{supp}\mu K_\epsilon `$ and $$\left\|f𝑑\mu f𝑑\mu _\epsilon \right\|f_{\mathrm{}}\underset{K𝒦(\epsilon )}{}\|\mu (K)\mu _\epsilon (K)\|+\gamma _\epsilon (\mu _\epsilon +\mu )(K_\epsilon ).$$ It is clear that the second term on the right-hand side goes to zero with $`\epsilon `$. For the first term we note that from (III.27), (III.30), we have $`\|\mu (K)\mu _\epsilon (K)\|2\pi /h_{\mathrm{ex}}`$ while the number of squares in $`𝒦(\epsilon )`$ is of the order of $`1/\delta ^2`$. From (III.26) it then follows that $`_{K𝒦(\epsilon )}\|\mu (K)\mu _\epsilon (K)\|`$ tends to zero with $`\epsilon `$. We thus have $`lim_{\epsilon 0}f𝑑\mu _\epsilon =f𝑑\mu `$ and (III.4) follows. We easily deduce (III.7) from (III.4). Indeed from (H2) and the fact that each point is at a distance at most $`1/h_{\mathrm{ex}}`$ from a pinning site, we get that $`a_\epsilon (p)b(p)`$ as $`\epsilon 0`$, uniformly in $`p𝒫(\epsilon )`$. Moreover, since $`n_\epsilon /h_{\mathrm{ex}}`$ is bounded, $$\underset{\epsilon 0}{lim}\frac{2\pi {\displaystyle \underset{i=1}{\overset{n_\epsilon }{}}}a_\epsilon (p_\epsilon ^i)}{h_{\mathrm{ex}}}=\underset{\epsilon 0}{lim}\frac{2\pi {\displaystyle \underset{i=1}{\overset{n_\epsilon }{}}}b(p_\epsilon ^i)}{h_{\mathrm{ex}}}=b(x)𝑑\mu (x),$$ by the convergence of $`\mu _\epsilon `$ to $`\mu `$. It remains to prove (III.5). We split the sum in (III.5) as follows: let $`(\epsilon )`$ be the set of pairs of indices $`(i,j)`$ such that $`1ijn_\epsilon `$ and $`p_\epsilon ^i,p_\epsilon ^j`$ belong to the same square of the subdivision $`𝒦(\epsilon )`$. Let $`𝒥(\epsilon )`$ be pairs $`(i,j)`$ such that $`p_\epsilon ^i,p_\epsilon ^j`$ belong to different squares. Then (III.31) $$\underset{\begin{array}{c}ij\\ \|p_\epsilon ^ip_\epsilon ^j\|<\alpha \end{array}}{}\left\|\mathrm{log}\|p_\epsilon ^ip_\epsilon ^j\|\right\|=\underset{\begin{array}{c}(i,j)(\epsilon )\\ \|p_\epsilon ^ip_\epsilon ^j\|<\alpha \end{array}}{}\left\|\mathrm{log}\|p_\epsilon ^ip_\epsilon ^j\|\left\|+\underset{\begin{array}{c}(i,j)𝒥(\epsilon )\\ \|p_\epsilon ^ip_\epsilon ^j\|<\alpha \end{array}}{}\right\|\mathrm{log}\|p_\epsilon ^ip_\epsilon ^j\|\right\|$$ The first sum in (III.31) is estimated as follows. For every $`K𝒦(\epsilon )`$, $`\mu (K)<C\delta ^2`$ thus the number of points of $`𝒫(\epsilon )`$ in $`K`$ is less than $`C\delta ^2h_{\mathrm{ex}}`$. The number of squares being of the order of $`\delta ^2`$, the cardinal of $`(\epsilon )`$ is less than $`C\delta ^2h_{\mathrm{ex}}^2`$. Using (III.26), (III.27) and (III.28), we find (III.32) $$\underset{\begin{array}{c}(i,j)(\epsilon )\\ \|p_\epsilon ^ip_\epsilon ^j\|<\alpha \end{array}}{}\left\|\mathrm{log}\|p_\epsilon ^ip_\epsilon ^j\|\right\|Ch_{\mathrm{ex}}^2\delta ^2\mathrm{log}\|\mathrm{log}\epsilon \|h_{\mathrm{ex}}^2.$$ To treat the second sum in (III.31), we note that if $`K`$ and $`K^{}`$ are distinct squares in $`𝒦(\epsilon )`$ and $`pK`$, $`qK^{}`$ then $$xK,yK^{},\|xy\|4\|pq\|.$$ Thus we may write, using the fact that $`\mu (K)<C\delta ^2`$, $$\underset{\begin{array}{c}ij\\ p_\epsilon ^iK,p_\epsilon ^jK^{}\end{array}}{}\left\|\mathrm{log}\|p_\epsilon ^ip_\epsilon ^j\|\right\|Ch_{\mathrm{ex}}^2_{K\times K^{}}(\left\|\mathrm{log}\|xy\|\right\|+1)𝑑x𝑑y.$$ Summing over pairs of squares $`K,K^{}𝒦(\epsilon )`$ such that $`K\times K^{}`$ intersects $`\{(x,y)\|xy\|<\alpha \}`$ we get for $`\epsilon `$ small enough (III.33) $$\underset{\begin{array}{c}(i,j)𝒥(\epsilon )\\ \|p_\epsilon ^ip_\epsilon ^j\|<\alpha \end{array}}{}\left\|\mathrm{log}\|p_\epsilon ^ip_\epsilon ^j\|\right\|Ch_{\mathrm{ex}}^2_{\|xy\|<2\alpha }(\left\|\mathrm{log}\|xy\|\right\|+1)𝑑x𝑑y.$$ Summing (III.32), (III.33), dividing by $`h_{\mathrm{ex}}^2`$ and letting $`\epsilon `$ and then $`\alpha `$ tend to zero yields (III.5). Proposition III.2 is proved. $``$ ### III.3 Proof of Lemma III.1 The fact that $`G_\epsilon `$ and $`G_0`$ are positive is a simple consequence of the maximum principle, that they are symmetric is standard and follows from Green’s identity. The inequality $$G_\epsilon (x,y),G_0(x,y)<C\mathrm{log}\|xy\|+C$$ is a well known property of Green functions for elliptic operators in divergence form, a proof can be found in \[St\]. To prove property 3), we let $$v_\epsilon (x,y)=G_\epsilon (x,y)+\frac{a_\epsilon (y)}{2\pi }\mathrm{log}\|xy\|$$ and $`L_\epsilon `$ be the operator $`u\mathrm{div}(𝓐_𝜺u)+u`$. Then letting $`f_\epsilon =L_\epsilon v_\epsilon (.,y)`$, we have (III.34) $$f_\epsilon (x,y)=\frac{a_\epsilon (y)}{2\pi }\frac{1}{a_\epsilon (x)}._x\mathrm{log}\|xy\|\frac{a_\epsilon (y)}{2\pi }\mathrm{log}\|xy\|.$$ Thus for any $`1q<2`$, there is a $`C`$ independent of $`y`$ and $`\epsilon `$, such that $`f_\epsilon (.,y)_{L^q}C/\eta (\epsilon )`$. On the other hand, $`v_\epsilon (.,y)`$ is bounded in $`W^{1,q}(\mathrm{\Omega })`$ independently of $`\epsilon `$ and $`y`$ (see \[St\]). Now, Theorem 2 of \[Me\] implies that there exist $`p>2`$ and $`p^{}<2`$ such that if $`u`$ satisfies $`L_\epsilon u=f`$, then for any compact $`K\mathrm{\Omega }`$, $$u_{L^p(K)}C(K)\left(u_{L^p^{}(\mathrm{\Omega })}+f_{W^{1,p}(\mathrm{\Omega })}\right).$$ We may choose $`q<2`$ such that $`W^{1,p}L^q`$ and $`p^{}<q`$. Thus, we find that $`v_\epsilon (.,y)`$ is bounded in $`W^{1,p}(K)`$ by $`C/\eta (\epsilon )`$. Since $`p>2`$, this yields the uniform bound $`xK,y\mathrm{\Omega }`$, $$\|v_\epsilon (x,y)\|\frac{C(K)}{\eta (\epsilon )}$$ i.e. property 3). To prove property 4), we note that for any $`\alpha >0`$, $`L_\epsilon G_\epsilon (.,y)=0`$ in $`\mathrm{\Omega }B(y,\alpha )`$ while $`G_\epsilon (.,y)`$ is bounded in $`W^{1,q}(\mathrm{\Omega })`$ independently of $`\epsilon `$ and $`y`$ (see \[St\]). Using the aforementioned result of \[Me\], we find that $`G_\epsilon (.,y)`$ is bounded in $`W_{\text{loc}}^{1,p}(\mathrm{\Omega }B(y,\alpha ))`$, for some $`p>2`$, independently of $`y`$ and $`\epsilon `$, thus $`G_\epsilon `$ converges locally uniformly in $`\mathrm{\Omega }\times \mathrm{\Omega }\mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ is the diagonal. The limit is necessarily $`G_0`$, since $`G_0(.,y)`$ satisfies $`L_0G_0(.,y)=\mathrm{div}𝓐_\mathrm{𝟎}_xG_0+G_0=\delta _y`$ and $`L_\epsilon `$ $`H`$-converges to $`L_0`$. Lemma II.1 is proved. $``$ ## IV Convergence results We can then proceed as in the rest of Section III in \[SS3\]. ###### Proposition IV.1 The minimum of $`E`$ is uniquely achieved by $`h_{}C^{1,\gamma }(\mathrm{\Omega })(\gamma <1)`$ satisfying (IV.1) $$\{\begin{array}{cc}h_{}1\frac{\mathrm{\Lambda }b}{2}\hfill & \mathrm{in}\mathrm{\Omega }\hfill \\ h_{}=1\hfill & \mathrm{on}\mathrm{\Omega }\hfill \\ \mu _{}:=\mathrm{div}(𝓐_\mathrm{𝟎}h_{})+h_{}0\hfill & \\ \left(h_{}\left(1\frac{\mathrm{\Lambda }b}{2}\right)\right)\mu _{}=0\hfill & \end{array}$$ As in \[SS3\], we divide the proof of this proposition into several lemmas. ###### Lemma IV.1 Let $`\mu _{}^+`$ and $`\mu _{}^{}`$ be the positive and negative parts of the measure $`\mu _{}`$. Then $`h_{}=1{\displaystyle \frac{\mathrm{\Lambda }b}{2}}`$ $`\mu _{}^+\text{ a.e.}`$ $`h_{}=1+{\displaystyle \frac{\mathrm{\Lambda }b}{2}}`$ $`\mu _{}^{}\text{ a.e.}`$ $$1\frac{\mathrm{\Lambda }b}{2}h_{}1+\frac{\mathrm{\Lambda }b}{2}.$$ Proof : As in \[SS3\], the minimum of $`E`$ is achieved by some $`h_{}`$, by lower semi-continuity. Performing variations $`(1+tf)\mu _{}`$ where $`fC^0(\mathrm{\Omega })`$, and looking at the first order in $`t0`$, we find similarly as in \[SS3\] that $$\frac{\mathrm{\Lambda }b}{2}\|\mu _{}\|+(h_{}1)\mu _{}=0.$$ Hence, $`h_{}=1{\displaystyle \frac{\mathrm{\Lambda }b}{2}}`$ $`\mu _{}^+\text{ a.e.}`$ $`h_{}=1+{\displaystyle \frac{\mathrm{\Lambda }b}{2}}`$ $`\mu _{}^{}\text{ a.e.}.`$ As in \[SS3\], considering variations $`\mu _{}+\nu `$, where $`\nu H^1`$ and $`\nu `$ and $`\mu _{}`$ are mutually singular, we are led to $`1\frac{\mathrm{\Lambda }b}{2}h_{}1+\frac{\mathrm{\Lambda }b}{2}.`$ $``$ ###### Lemma IV.2 $`\mu _{}`$ is a positive measure. Proof : $$_\mathrm{\Omega }\mu _{}(h_{}1)_+=_\mathrm{\Omega }\mu _{}^+(h_{}1)_+_\mathrm{\Omega }\mu _{}^{}(h_{}1)_+.$$ Since $`(h_{}1)_+=0`$ $`\mu _{}^+`$-a.e., we have $`{\displaystyle _\mathrm{\Omega }}\mu _{}(h_{}1)_+`$ $`=`$ $`{\displaystyle _\mathrm{\Omega }}\mu _{}^{}(h_{}1)_+`$ $`=`$ $`{\displaystyle _\mathrm{\Omega }}(\mathrm{div}(𝓐_\mathrm{𝟎}h_{})+h_{})(h_{}1)_+`$ $`=`$ $`{\displaystyle _{h_{}>1}}h_{}(𝓐_\mathrm{𝟎}h_{})+h_{}(h_{}1)0,`$ because $`𝓐_\mathrm{𝟎}`$ is a symmetric positive matrix (this follows from the compactness of the set of matrices bounded from above and below). We deduce that $$_\mathrm{\Omega }\mu _{}^{}(h_{}1)_+=0,$$ but since $`h_{}1=\frac{\mathrm{\Lambda }b}{2}`$, $`\mu _{}^{}`$ a.e., we have $$_\mathrm{\Omega }\frac{\mathrm{\Lambda }b}{2}\mu _{}^{}=0,$$ hence $`\mu _{}^{}=0`$, and $`\mu _{}0`$. $``$ Thus, $`h_{}`$ satisfies all the properties listed in (IV.1). We can now complete the convergence results. From the upper bound of Proposition I.4 and Lemma II.4, we deduce that for our family of minimizers $`(u_\epsilon ,A_\epsilon )`$, $$\underset{V}{\mathrm{min}}E=E(h_{})\underset{\epsilon 0}{lim\; inf}\frac{J_\epsilon (u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}E(h_0)E(h_{}).$$ $`h_{}`$ being the unique minimizer of $`E`$, we conclude that $`h_0=h_{}`$ and thus $`\mu _0=\mu _{}`$. We also obtain (IV.2) $$\underset{\epsilon 0}{lim}\frac{J_\epsilon (u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}=E(h_{}).$$ Since the possible limits are unique, the whole family $`\frac{h_\epsilon }{h_{\mathrm{ex}}}`$ converges to $`h_{}`$, and the same for $`\mu _\epsilon `$. In view of (II.37), we have $`\underset{\epsilon 0}{lim\; inf}{\displaystyle \frac{J(u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}}`$ $``$ $`\underset{\epsilon 0}{lim\; inf}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\|\mathrm{log}\epsilon \|}{h_{\mathrm{ex}}}}{\displaystyle _\mathrm{\Omega }}a_\epsilon \|\mu _\epsilon \|\right)+{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Omega }}h_{}𝓐_\mathrm{𝟎}h_{}+\|h_{}1\|^2.`$ $``$ $`{\displaystyle \frac{\mathrm{\Lambda }}{2}}{\displaystyle _\mathrm{\Omega }}b\|\mu _{}\|+{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Omega }}h_{}𝓐_\mathrm{𝟎}h_{}+\|h_{}1\|^2,`$ while $$\underset{\epsilon 0}{lim\; sup}\frac{J(u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}\frac{\mathrm{\Lambda }}{2}_\mathrm{\Omega }b\|\mu _{}\|+\frac{1}{2}_\mathrm{\Omega }h_{}𝓐_\mathrm{𝟎}h_{}+\|h_{}1\|^2.$$ Thus, we deduce that $$\underset{\epsilon 0}{lim}_\mathrm{\Omega }a_\epsilon \|\mu _\epsilon \|=_\mathrm{\Omega }b\mu _{}.$$ On the other hand, $$\underset{\epsilon 0}{lim\; inf}_\mathrm{\Omega }a_\epsilon \|\mu _\epsilon \|\underset{\epsilon 0}{lim\; inf}_\mathrm{\Omega }b\|\mu _\epsilon \|_\mathrm{\Omega }b\|\mu _{}\|,$$ hence $`_\mathrm{\Omega }b\|\mu _\epsilon \|_\mathrm{\Omega }b\mu _{}`$, while $`_\mathrm{\Omega }b\mu _\epsilon _\mathrm{\Omega }b\mu _{}`$. We conclude that $`_\mathrm{\Omega }b(\|\mu _\epsilon \|\mu _\epsilon )0`$ and thus $`\|\mu _\epsilon \|`$ and $`\mu _\epsilon `$ have the same limiting measure $`\mu _{}`$. This proves (I.16), (I.17), and (I.18). Following \[SS3\], Section IV, we can also prove easily the following : ###### Proposition IV.2 If $`\mathrm{\Lambda }=0`$, then $`h_{}=1`$ and $`\frac{h_\epsilon }{h_{\mathrm{ex}}}10`$ strongly in $`H_0^1(\mathrm{\Omega })`$. If $`\mathrm{\Lambda }>0`$, then $`\frac{h_\epsilon }{h_{\mathrm{ex}}}1h_{}1`$ in $`H_0^1(\mathrm{\Omega })`$, the convergence is not strong and $$\frac{\|h_\epsilon \|^2}{h_{\mathrm{ex}}^2a_\epsilon }h_{}𝓐_\mathrm{𝟎}h_{}+\mathrm{\Lambda }b\mu _{}\text{in }.$$ Proof : First, it is easy to get, as seen in Lemma II.4 for example, that $$_\mathrm{\Omega }\|_{A_\epsilon }u_\epsilon \|^2_\mathrm{\Omega }\frac{\|h_\epsilon \|^2}{a_\epsilon }(1o(1)),$$ thus, we have (IV.3) $`\underset{\epsilon 0}{lim\; inf}{\displaystyle \frac{J(u_\epsilon ,A_\epsilon )}{h_{\mathrm{ex}}^2}}`$ $``$ $`\underset{\epsilon 0}{lim\; inf}{\displaystyle \frac{1}{h_{\mathrm{ex}}^2}}\left({\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Omega }}{\displaystyle \frac{\|h_\epsilon \|^2}{a_\epsilon }}+\|h_\epsilon h_{\mathrm{ex}}\|^2\right)`$ (IV.4) $``$ $`{\displaystyle \frac{\mathrm{\Lambda }}{2}}{\displaystyle _\mathrm{\Omega }}b\mu _{}+{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Omega }}h_{}𝓐_\mathrm{𝟎}h_{}+\|h_{}1\|^2.`$ The case $`\mathrm{\Lambda }=0`$ follows easily from the upper bound $`\mathrm{min}J_\epsilon (u_\epsilon ,A_\epsilon )o(h_{\mathrm{ex}}^2)`$ of Section II combined with (IV.4). The convergence of $`\frac{h_\epsilon }{h_{\mathrm{ex}}}`$ to $`h_{}`$ is weak in $`H^1`$, in general, thus strong in $`L^2(\mathrm{\Omega })`$, and $$\underset{\epsilon 0}{lim}_\mathrm{\Omega }\left\|\frac{h_\epsilon }{h_{\mathrm{ex}}}1\right\|^2=_\mathrm{\Omega }\|h_{}1\|^2.$$ Combining this to the convergence result (IV.2), we have (IV.5) $$\underset{\epsilon 0}{lim}\frac{1}{2}_\mathrm{\Omega }\frac{\|h_\epsilon \|^2}{h_{\mathrm{ex}}^2a_\epsilon }=\frac{\mathrm{\Lambda }}{2}_\mathrm{\Omega }b\mu _{}+\frac{1}{2}_\mathrm{\Omega }h_{}𝓐_\mathrm{𝟎}h_{}.$$ Then, we argue as in \[SS3\], Proposition IV.1. Roughly speaking, one considers any open set $`U\mathrm{\Omega }`$, and gets a lower bound $`\underset{\epsilon 0}{lim\; inf}{\displaystyle _U}{\displaystyle \frac{\|h_\epsilon \|^2}{h_{\mathrm{ex}}^2a_\epsilon }}`$ $`=`$ $`\underset{\epsilon 0}{lim\; inf}{\displaystyle _{U(_iB_i)}}{\displaystyle \frac{\|h_\epsilon \|^2}{h_{\mathrm{ex}}^2a_\epsilon }}+{\displaystyle _{U\backslash _iB_i}}{\displaystyle \frac{\|h_\epsilon \|^2}{h_{\mathrm{ex}}^2a_\epsilon }}`$ $``$ $`\mathrm{\Lambda }{\displaystyle _U}b\|\mu _\epsilon \|+{\displaystyle _U}h_{}𝓐_\mathrm{𝟎}h_{}`$ $``$ $`\mathrm{\Lambda }{\displaystyle _U}b\mu _{}+{\displaystyle _U}h_{}𝓐_\mathrm{𝟎}h_{}.`$ Since this is true for any $`U\mathrm{\Omega }`$, comparing this to (IV.4) and (IV.5), we obtain as in \[SS3\], $$\frac{\|h_\epsilon \|^2}{h_{\mathrm{ex}}^2a_\epsilon }h_{}𝓐_\mathrm{𝟎}h_{}+\mathrm{\Lambda }b\mu _{}\text{in }.$$ $``$ This completes the proof of Theorems 1, 2 and 3. Acknowledgments : The authors are very grateful to François Murat for taking time explaining the basis of homogenization and pointing out the good references. They would also like to thank very much Jon Chapman for fruitful discussions on pinning models and Alano Ancona for pointing out references on Green functions.
warning/0004/math0004066.html	ar5iv	text	# Generalized Toric Varieties for Simple Non-Rational Convex Polytopes ## Introduction Consider a vector space d of dimension $`n`$. To each simple convex polytope $`\mathrm{\Delta }\text{d}^{}`$ that is rational with respect to a lattice in d there corresponds a toric variety with at worst quotient singularities. What happens in the case that the simple convex polytope is no longer rational? To answer this question we consider a special class of spaces, called quasifolds, which were first introduced by one of the authors in \[P\]. A quasifold is not necessarily a Hausdorff space: it is locally modeled by orbit spaces of the action of discrete, possibly infinite, groups on open subsets of $`\text{}^k`$. A quasitorus, on the other hand, is the natural replacement of a torus in this geometry. In this article we define the notions of complex quasifold and complex quasitorus and we associate to each simple convex polytope $`\mathrm{\Delta }\text{d}^{}`$ a family of compact complex quasifolds of dimension $`n`$, each endowed with the holomorphic action of a complex quasitorus $`D_{\text{}}`$ having a dense open orbit. Our construction is explicit: each space $`M`$ is the topological quotient of a suitable open subset of $`\text{}^d`$ by the action of a suitable subgroup $`N_{\text{}}T_{\text{}}^d=\text{}^d/\text{}^d`$, and $`D_{\text{}}`$ is isomorphic to $`T_{\text{}}^d/N_{\text{}}`$, $`d`$ being the number of facets of the polytope. We show that $`M`$ is a complex quasifold by covering it with mutually compatible local models of the type $`\text{}^n`$ modulo the action of a discrete group. If the polytope is rational our procedure matches the standard one for constructing toric varieties as quotients that is described in Appendix 1 of Guillemin’s book \[G\]. It is shown in \[P\] that to the same simple convex polytope $`\mathrm{\Delta }`$ one can also associate a family of symplectic quasifolds of dimension $`2n`$, each endowed with the effective Hamiltonian action of a quasitorus $`D`$, having the property that the image of the corresponding moment mapping is exactly $`\mathrm{\Delta }`$. The construction extends the one given by Delzant \[D\] in the smooth case: it is explicit and uses the symplectic quotient operation. In the last section of this paper we show that each complex quotient $`M`$ is equivariantly diffeomorphic to one of these symplectic quotients, that the induced symplectic structure is compatible with the complex one, and thus defines on $`M`$ the structure of a Kähler quasifold; of course $`D_{\text{}}`$ here is the complexification of the corresponding quasitorus $`D`$. For these reasons these spaces may well be thought of as a natural generalization of toric varieties for simple convex polytopes that are not rational. ## 1 Complex quasifolds and complex quasitori This section is devoted to defining complex quasifolds and their geometry. We will not repeat the remarks and results that are in common with the real case, for which we refer the reader to the article \[P\]. The local model for complex quasifolds is a complex manifold acted on holomorphically by a discrete group. ###### Definition 1.1 (Complex model) Let $`\stackrel{~}{V}`$ be a connected, simply connected complex manifold of dimension $`k`$ and let $`\mathrm{\Gamma }`$ be a discrete group acting on $`\stackrel{~}{V}`$ holomorphically so that the set of points, $`\stackrel{~}{V}_0`$, where the action is free, is connected and dense. Consider the space of orbits, $`\stackrel{~}{V}/\mathrm{\Gamma }`$, of the action of the group $`\mathrm{\Gamma }`$ on the manifold $`\stackrel{~}{V}`$, endowed with the quotient topology, and the canonical projection $`p:\stackrel{~}{V}\stackrel{~}{V}/\mathrm{\Gamma }`$. A complex model of dimension $`k`$ is the triple $`(\stackrel{~}{V}/\mathrm{\Gamma },p,\stackrel{~}{V})`$, shortly $`\stackrel{~}{V}/\mathrm{\Gamma }`$. ###### Remark 1.2 We remark that the assumption in Definition 1.1 that the manifold $`\stackrel{~}{V}`$ be simply connected is not as strong as one may think. Consider the triple $`(\stackrel{~}{V}/\mathrm{\Gamma },p,\stackrel{~}{V})`$ as defined above but assume that the manifold $`\stackrel{~}{V}`$ is not simply connected; consider its universal cover, $`\pi :V^\mathrm{\#}\stackrel{~}{V}`$, and its fundamental group, $`\mathrm{\Pi }`$. The manifold $`V^\mathrm{\#}`$ is connected and simply connected, the mapping $`\pi `$ is holomorphic, the discrete group $`\mathrm{\Pi }`$ acts holomorphically, freely and properly on the manifold $`V^\mathrm{\#}`$ and $`\stackrel{~}{V}=V^\mathrm{\#}/\mathrm{\Pi }`$. Consider the extension of the group $`\mathrm{\Gamma }`$ by the group $`\mathrm{\Pi }`$, $`1\mathrm{\Pi }\mathrm{\Lambda }\mathrm{\Gamma }1`$, defined as follows $$\mathrm{\Lambda }=\left\{\lambda \text{Diff}(V^\mathrm{\#})\right\|\gamma \mathrm{\Gamma }\text{s. t.}\pi (\lambda (u^\mathrm{\#}))=\gamma \pi (u^\mathrm{\#})u^\mathrm{\#}V^\mathrm{\#}\}.$$ It is easy to verify that $`\mathrm{\Lambda }`$ is a discrete group, that it acts on the manifold $`V^\mathrm{\#}`$ according to the assumptions of Definition 1.1 and that $`\stackrel{~}{V}/\mathrm{\Gamma }=V^\mathrm{\#}/\mathrm{\Lambda }`$. ###### Definition 1.3 (Submodel) Consider a model $`(\stackrel{~}{V}/\mathrm{\Gamma },p,\stackrel{~}{V})`$ and let $`W`$ be an open subset of $`\stackrel{~}{V}/\mathrm{\Gamma }`$. We will say that $`W`$ is a submodel of $`(\stackrel{~}{V}/\mathrm{\Gamma },p,\stackrel{~}{V})`$, if $`(p^1(W),p,W)`$ defines a model according to Remark 1.2. ###### Definition 1.4 (Holomorphic mapping, biholomorphism of models) Given two models $`(\stackrel{~}{V}/\mathrm{\Gamma },p,\stackrel{~}{V})`$ and $`(\stackrel{~}{W}/\mathrm{\Delta },q,\stackrel{~}{W})`$, a mapping $`f:\stackrel{~}{V}/\mathrm{\Gamma }\stackrel{~}{W}/\mathrm{\Delta }`$ is said to be holomorphic if there exists a holomorphic mapping $`\stackrel{~}{f}:\stackrel{~}{V}\stackrel{~}{W}`$ such that $`q\stackrel{~}{f}=fp`$; we will say that $`\stackrel{~}{f}`$ is a lift of $`f`$. We will say that the holomorphic mapping $`f`$ is a biholomorphism of models if it is bijective and if the lift $`\stackrel{~}{f}`$ is a biholomorphism. If the mapping $`\stackrel{~}{f}`$ is a lift of a holomorphic mapping of models $`f:\stackrel{~}{U}/\mathrm{\Gamma }\stackrel{~}{V}/\mathrm{\Delta }`$ so are the mappings $`\stackrel{~}{f}^\gamma ()=\stackrel{~}{f}(\gamma )`$, for all elements $`\gamma `$ in $`\mathrm{\Gamma }`$ and $`{}_{}{}^{\delta }\stackrel{~}{f}()=\delta \stackrel{~}{f}()`$, for all elements $`\delta `$ in $`\mathrm{\Delta }`$. If the mapping $`f`$ is a biholomorphism, then these are the only other possible lifts and the groups $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ are isomorphic; the proof goes exactly as in the real case, see \[P, orange and green lemmas\]. Geometric objects on a model $`\stackrel{~}{V}/\mathrm{\Gamma }`$ are defined by $`\mathrm{\Gamma }`$-invariant geometric objects on the manifold $`\stackrel{~}{V}`$. For example: ###### Definition 1.5 (Differential form on a model) A differential form of degree $`h`$, $`\omega `$, on a model $`\stackrel{~}{V}/\mathrm{\Gamma }`$ is the assignment of a $`\mathrm{\Gamma }`$-invariant differential form of degree $`h`$, $`\stackrel{~}{\omega }`$, on the complex manifold $`\stackrel{~}{V}`$. ###### Definition 1.6 (Kähler form on a model) A Kähler form on a model $`\stackrel{~}{V}/\mathrm{\Gamma }`$ is a differential form, $`\omega `$, such that $`\stackrel{~}{\omega }`$ (see Definition 1.5) is Kähler. Complex quasifolds are obtained by gluing together the models in the appropriate way: ###### Definition 1.7 (Complex quasifold) A dimension $`k`$ complex quasifold structure on a topological space $`M`$ is the assignment of an atlas, or collection of charts, $`𝒜=\{(V_\alpha ,\varphi _\alpha ,\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha )\|\alpha A\}`$ having the following properties: 1. The collection $`\{V_\alpha \|\alpha A\}`$ is a cover of $`M`$. 2. For each index $`\alpha `$ in $`A`$, the set $`V_\alpha `$ is open, the space $`\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha `$ defines a model, where the set $`\stackrel{~}{V}_\alpha `$ is an open, connected, and simply connected subset of the space $`\text{}^k`$, and the mapping $`\varphi _\alpha `$ is a homeomorphism of the space $`\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha `$ onto the set $`V_\alpha `$. 3. For all indices $`\alpha ,\beta `$ in $`A`$ such that $`V_\alpha V_\beta \mathrm{}`$, the sets $`\varphi _\alpha ^1(V_\alpha V_\beta )`$ and $`\varphi _\beta ^1(V_\alpha V_\beta )`$ are submodels of $`\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha `$ and $`\stackrel{~}{V}_\beta /\mathrm{\Gamma }_\beta `$ respectively and the mapping $$g_{\alpha \beta }=\varphi _\beta ^1\varphi _\alpha :\varphi _\alpha ^1(V_\alpha V_\beta )\varphi _\beta ^1(V_\alpha V_\beta )$$ is a biholomorphism of models. We will then say that the mapping $`g_{\alpha \beta }`$ is a change of charts and that the corresponding charts are compatible. 4. The atlas $`𝒜`$ is maximal, that is: if the triple $`(V,\varphi ,\stackrel{~}{V}/\mathrm{\Gamma })`$ satisfies property 2. and is compatible with all the charts in $`𝒜`$, then $`(V,\varphi ,\stackrel{~}{V}/\mathrm{\Gamma })`$ belongs to $`𝒜`$. We will say that a space $`M`$ with a complex quasifold structure is a complex quasifold. A complex quasifold of dimension $`k`$ has an underlying structure of real quasifold of dimension $`2k`$. ###### Remark 1.8 To each point $`mM`$ there corresponds a discrete group $`\mathrm{\Gamma }_m`$ defined as follows: take a chart $`(V_\alpha ,\varphi _\alpha ,\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha )`$ around $`m`$, then $`\mathrm{\Gamma }_m`$ is the isotropy group of $`\mathrm{\Gamma }_\alpha `$ at any point $`\stackrel{~}{v}\stackrel{~}{V}`$ which projects down to $`m`$. One can check that this definition does not depend on the choice of the chart. ###### Definition 1.9 (Holomorphic mapping, biholomorphism) Let $`M`$ and $`N`$ be two complex quasifolds. A continuous mapping $`f:MN`$ is said to be a holomorphic mapping if there exists a chart $`(V_\alpha ,\varphi _\alpha ,\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha )`$ around each point $`m`$ in the space $`M`$, a chart $`(W_\alpha ,\psi _\alpha ,\stackrel{~}{W}_\alpha /\mathrm{\Delta }_\alpha )`$ around the point $`f(m)`$, and a holomorphic mapping of models $`f_\alpha :\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha \stackrel{~}{W}_\alpha /\mathrm{\Delta }_\alpha `$ such that $`\psi _\alpha f_\alpha =f\varphi _\alpha `$. If $`f`$ is bijective, and if each $`f_\alpha `$ is a biholomorphism of models, we will say that $`f`$ is a biholomorphism. We remark that two biholomorphic quasifolds are also diffeomorphic with respect to their underlying real quasifold structure. Geometric objects on quasifolds are defined as geometric objects on the charts that respect the changes of charts. For example: ###### Definition 1.10 (Differential form) A differential form of degree $`h`$, $`\omega `$, on a complex quasifold $`M`$ is the assignment of a chart $`(V_\alpha ,\varphi _\alpha ,\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha )`$ around each point $`mM`$ and of a differential form of degree $`h`$, $`\omega _\alpha `$, on the model $`\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha `$. We require that whenever we have two such charts, $`(V_\alpha ,\varphi _\alpha ,\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha )`$ and $`(V_\beta ,\varphi _\beta ,\stackrel{~}{V}_\beta /\mathrm{\Gamma }_\beta )`$, with the property that $`V_\alpha V_\beta \mathrm{}`$, then $`(g_{\alpha \beta })^{}\omega _\beta =\omega _\alpha `$ for the corresponding change of charts $`g_{\alpha \beta }`$. ###### Definition 1.11 (Kähler form) A Kähler form, $`\omega `$, on a complex quasifold $`M`$ is a differential form, $`\omega `$, such that each $`\omega _\alpha `$ (see Definition 1.10) is a Kähler form on the model $`\stackrel{~}{V}_\alpha /\mathrm{\Gamma }_\alpha `$. A Kähler structure on a quasifold $`M`$ is the assignment of a Kähler form, $`\omega `$, and we will say that $`(M,\omega )`$, or simply $`M`$, is a Kähler quasifold. The analogue of a torus in the real setting is a quasitorus. Let d be a vector space of dimension $`n`$. We recall from \[P\] that a quasitorus of dimension $`n`$ and quasi-Lie algebra d is the quotient of the space d by a quasilattice $`Q`$, which in turn is the -span of a set of spanning vectors in d. Consider now the complexification of d, the complex vector space $`\text{d}_{\text{}}=\text{d}+i\text{d}`$; the quasilattice $`Q`$ acts naturally on $`\text{d}_{\text{}}`$: $$\begin{array}{ccccc}Q\hfill & \times & \text{d}_{\text{}}& & \text{d}_{\text{}}\\ (A\hfill & ,& X+iY)& & (X+A)+iY.\end{array}$$ (1) ###### Definition 1.12 (Complex quasitorus, quasi-Lie algebra, exponential) Let d be a vector space of dimension $`n`$ and let $`Q`$ be a quasilattice in d. Then we call the quotient $`D_{\text{}}=\text{d}_{\text{}}/Q`$ complex quasitorus of dimension $`n`$ and quasi-Lie algebra $`\text{d}_{\text{}}`$. We call the corresponding projection $`\text{d}_{\text{}}D_{\text{}}`$ exponential mapping and we denote it by $`\mathrm{exp}`$. If the quasilattice $`Q`$ is a lattice $`L`$ we obtain the honest complex torus $`\text{d}_{\text{}}/L`$, which is the complexification of the torus $`\text{d}/L`$. The complex quasitorus $`D_{\text{}}`$ is a quasifold of one chart and may be naturally thought of as the complexification of the quasitorus $`D=\text{d}/Q`$. The main result that we will be needing is the following: ###### Proposition 1.13 Let $`T`$ be a torus and $`N`$ a Lie subgroup. Then $`T_{\text{}}/N_{\text{}}`$ is a complex quasitorus of complex dimension $`n=dimTdimN`$. Proof. This proposition has a real analogue, the proof of which is quite similar (compare with \[P, Proposition 2.5\]). Denote by n the Lie algebra of $`N`$, and by t the Lie algebra of $`T`$. Let d be a complement of n in t, then the complex vector space $`\text{d}_{\text{}}=\text{d}+i\text{d}`$ is a complement of $`\text{n}_{\text{}}`$ in $`\text{t}_{\text{}}`$. We define a projection $`p:\text{d}_{\text{}}T_{\text{}}/N_{\text{}}`$ by the rule $`p(Z)=\mathrm{\Pi }(\mathrm{exp}Z),Z\text{d}_{\text{}}`$, where $`\mathrm{\Pi }:T_{\text{}}T_{\text{}}/N_{\text{}}`$ denotes the canonical projection. Notice that, by definition of $`\mathrm{exp}`$ the kernel of $`p`$ is a quasilattice $`Q`$ and $`p`$ induces a group isomorphism $`\text{d}_{\text{}}/QT_{\text{}}/N_{\text{}}`$. $`\mathrm{}`$ We conclude this section with the definition of holomorphic action of a complex quasitorus on a complex quasifold. ###### Definition 1.14 (Holomorphic action) A holomorphic action of a complex quasitorus $`D_{\text{}}`$ on a complex quasifold $`M`$ is a holomorphic mapping $`\tau :D_{\text{}}\times MM`$ such that $`\tau (d_1d_2,m)=\tau (d_1,\tau (d_2,m))`$ and $`\tau (1_D_{\text{}},m)=m`$ for all elements $`d_1,d_2`$ in the quasitorus $`D_{\text{}}`$ and for each point $`m`$ in the space $`M`$. ## 2 From simple polytopes to complex quasifolds Let d be a real vector space of dimension $`n`$, and let $`\mathrm{\Delta }`$ be a simple convex polytope of dimension $`n`$ in the dual space $`\text{d}^{}`$ (we recall that a polytope of dimension $`n`$ is simple if there are exactly $`n`$ edges stemming from each of its vertices). It is our intention to associate to this polytope a family of complex quasifolds, in much the same way that one associates a toric variety to a simple convex polytope that is rational. To do so we follow and extend the procedure for constructing toric varieties as global quotients that is described by Guillemin in \[G\]. Write the polytope as $$\mathrm{\Delta }=\underset{j=1}{\overset{d}{}}\{\mu \text{d}^{}\|\mu ,X_j\lambda _j\}$$ (2) for some elements $`X_1,\mathrm{},X_d`$ in the vector space d and some real numbers $`\lambda _1,\mathrm{},\lambda _d`$. Let $`Q`$ be a quasilattice in the space d containing the elements $`X_j`$ (for example the one that is generated by these elements) and let $`\{e_1,\mathrm{},e_d\}`$ denote the standard basis of $`\text{}^d`$ and $`\text{}^d`$; consider the surjective linear mapping $$\begin{array}{cccc}\pi :& \text{}^d& & \text{d}\\ & e_j& & X_j,\end{array}$$ and its complexification $$\begin{array}{cccc}\pi _{\text{}}:& \text{}^d& & \text{d}_{\text{}}\\ & e_j& & X_j.\end{array}$$ Consider the quasitorus $`\text{d}/Q`$ and its complexification $`\text{d}_{\text{}}/Q`$. The mappings $`\pi `$ and $`\pi _{\text{}}`$ each induce a group homomorphism, $$\mathrm{\Pi }:T^d=\text{}^d/\text{}^d\text{d}/Q$$ and $$\mathrm{\Pi }_{\text{}}:T_{\text{}}^d=\text{}^d/\text{}^d\text{d}_{\text{}}/Q.$$ We define $`N`$ to be the kernel of the mapping $`\mathrm{\Pi }`$ and $`N_{\text{}}`$ to be the kernel of the mapping $`\mathrm{\Pi }_{\text{}}`$. Notice that neither $`N`$ nor $`N_{\text{}}`$ is a torus unless $`Q`$ is a honest lattice. The mapping $`\mathrm{\Pi }_{\text{}}`$ defines an isomorphism $$T_{\text{}}^d/N_{\text{}}\text{d}_{\text{}}/Q$$ (3) ###### Remark 2.1 For the complexified group $`N_{\text{}}`$ the polar decomposition holds, namely $$N_{\text{}}=NA,$$ (4) where $`A=\mathrm{exp}(i\text{n})`$. In other words every element $`wN_{\text{}}`$ can be written uniquely as $`x\mathrm{exp}(iY)`$ where $`xN`$ and $`Y\text{n}`$. This follows from the definition of $`N`$ and $`N_{\text{}}`$, indeed $`N_{\text{}}=\{\mathrm{exp}(Z)\|Z\text{}^d\text{and}\pi _{\text{}}(Z)Q\}`$. Write $`Z=X+iY`$, then $`\pi _{\text{}}(Z)Q`$ if and only if $`\pi (X)Q`$ and $`\pi (Y)=0`$, which implies (4). Let $`F`$ denote a codimension-$`k`$ face of the polytope; this face is defined by a system of $`k`$ equalities: $`\mu ,X_j=\lambda _j`$, for $`jI\{1,\mathrm{},d\}`$. Then we consider the $`T_{\text{}}^d`$-orbit $`\text{}_F^d=\{(z_1,\mathrm{},z_d)\text{}^d\|z_j=0\text{iff}jI\}`$ and we take the union over all the possible faces of the polytope $$\text{}_\mathrm{\Delta }^d=\underset{F}{}\text{}_F^d.$$ The group $`N_{\text{}}`$ acts on the space $`\text{}_\mathrm{\Delta }^d`$. Consider the space of orbits $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$. We then have: ###### Theorem 2.2 Let d be a vector space of dimension $`n`$, and let $`\mathrm{\Delta }\text{d}^{}`$ be a simple convex polytope. Choose inward-pointing normals to the facets of $`\mathrm{\Delta }`$, $`X_1,\mathrm{},\text{X}_d\text{d}`$, and let $`Q`$ be a quasilattice containing these vectors. Then the corresponding quotient space $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$ is a complex quasifold of dimension $`n`$. The complex quasitorus $`\text{d}_{\text{}}/Q`$ acts on $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$, this action is holomorphic and has a dense open orbit. Before we prove this theorem we need a lemma that will be crucial not only now but also throughout the rest of the article. Consider any face $`F`$ of the polytope and denote by $`S_{\text{}}^F`$ the stabilizer of the $`T_{\text{}}^d`$-action on $`\text{}_F^d`$ and by $`\text{s}_{\text{}}^F`$ its Lie algebra. Then the fact that the polytope is simple implies that $$\text{n}_{\text{}}\text{s}_{\text{}}^F=\{0\},$$ (5) and this in turn implies that $$AS_{\text{}}^F=\{\text{Id}\},$$ (6) and that the following group is discrete $$N_{\text{}}S_{\text{}}^F=NS_{\text{}}^F.$$ (7) Equation (5) also implies that $`\pi `$, restricted to $`\text{s}_{\text{}}^F`$, is injective for any face $`F`$. In the special case that $`F`$ is equal to a vertex, $`\mu `$, we get that $`\text{s}_{\text{}}^\mu `$ is a complement of $`\text{n}_{\text{}}`$ in $`\text{}^d`$, and that the linear mapping $`\pi _{\text{}}`$, when restricted to $`\text{s}_{\text{}}^\mu `$, defines a (very useful) isomorphism $$\pi _\mu :\text{s}_{\text{}}^\mu \text{d}_{\text{}}.$$ (8) One can then deduce: ###### Lemma 2.3 Let $`\mu `$ be a vertex of the polytope $`\mathrm{\Delta }`$, consider the stabilizer $`S_{\text{}}^\mu `$ of the orbit $`\text{}_\mu ^d`$, its Lie algebra $`\text{s}_{\text{}}^\mu `$, and the discrete group $`\mathrm{\Gamma }_\mu =S_{\text{}}^\mu N_{\text{}}`$. Then we have that (i) $`T_{\text{}}^d/S_{\text{}}^\mu N_{\text{}}/\mathrm{\Gamma }_\mu `$; (ii) $`N_{\text{}}=\mathrm{\Gamma }_\mu \mathrm{exp}(\text{n}_{\text{}})`$; (iii) given any complement $`\text{b}_{\text{}}`$ of $`\text{s}_{\text{}}^\mu `$ in $`\text{}^d`$, we have that $$\text{n}_{\text{}}=\{V\pi _\mu ^1(\pi _{\text{}}(V))\|V\text{b}_{\text{}}\}.$$ Proof. (i) Consider the group homomorphism $$\begin{array}{cccc}\lambda _\mu :& N_{\text{}}& & T_{\text{}}^d/S_{\text{}}^\mu \\ & n& & [n].\end{array}$$ Since $`\text{n}_{\text{}}`$ and $`\text{s}_{\text{}}^\mu `$ are complementary, we have that $`T_{\text{}}^d=S_{\text{}}^\mu N_{\text{}}`$, therefore $`\lambda _\mu `$ is surjective. The kernel of $`\lambda _\mu `$ is given by $`\mathrm{\Gamma }_\mu `$, so $`\lambda _\mu `$ induces an isomorphism $`T_{\text{}}^d/S_{\text{}}^\mu N_{\text{}}/\mathrm{\Gamma }_\mu `$. (ii) Every element in $`N_{\text{}}`$ can be written in the form $`\mathrm{exp}(Z)`$, where $`Z\text{}^d`$ is such that $`\pi _{\text{}}(Z)Q`$. Write now $`Z=Z\pi _\mu ^1(\pi _{\text{}}(Z))+\pi _\mu ^1(\pi _{\text{}}(Z))`$; it is easy to check that $`Z\pi _\mu ^1(\pi _{\text{}}(Z))\text{n}_{\text{}}`$, and that $`\mathrm{exp}(\pi _\mu ^1(\pi _{\text{}}(Z)))\mathrm{\Gamma }_\mu `$. The group $`\mathrm{\Gamma }_\mu \mathrm{exp}(\text{n}_{\text{}})`$ is not necessarily trivial, so the decomposition is not necessarily unique. (iii) Every element of the form $`V\pi _\mu ^1(\pi _{\text{}}(V))`$, $`V\text{b}_{\text{}}`$ clearly belongs to $`\text{n}_{\text{}}`$. Conversely, write every element $`Z\text{n}_{\text{}}`$ as $`Z=U+V`$ according to the decomposition $`\text{}^d=\text{s}_{\text{}}^\mu \text{b}_{\text{}}`$, and notice that $`\pi _{\text{}}(Z)=0`$ implies that $`U=\pi _\mu ^1(\pi _{\text{}}(V))`$. $`\mathrm{}`$ Proof of Theorem 2.2. We want to define a complex quasifold atlas on the topological space $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$. We start by considering a rather special covering of $`\text{}_\mathrm{\Delta }^d`$. To do so we restrict our attention to each vertex $`\mu `$ of the polytope (the corresponding orbit $`\text{}_\mu ^d`$ has the smallest possible dimension $`dn`$) and we take the following $`T_{\text{}}^d`$-invariant, open and connected neighborhood of the orbit $`\text{}_\mu ^d`$ in $`\text{}^d`$: $$\widehat{V}_\mu =\{(z_1,\mathrm{},z_d)\text{}^d\|z_j0\text{if}jI\},$$ where $`I`$ is the index set corresponding to the vertex $`\mu `$ according to the recipe given below Remark 2.1. Notice that we have $`\text{}_\mathrm{\Delta }^d=_\mu \widehat{V}_\mu `$, where $`\mu `$ ranges over all the vertices of the polytope $`\mathrm{\Delta }`$. Indeed, take $`(z_1,\mathrm{},z_d)\text{}_\mathrm{\Delta }^d`$; then $`(z_1,\mathrm{},z_d)\text{}_G^d`$ for some face $`G`$, and $`(z_1,\mathrm{},z_d)\widehat{V}_\mu `$ for any vertex $`\mu `$ contained in the face $`G`$. The opposite inclusion holds because the polytope is simple. The neighborhoods $`\widehat{V}_\mu `$ are rather special, they are tubular. Let us check this. Fix a point $`\underset{¯}{z}^\mu `$ in the orbit $`\text{}_\mu ^d`$, for example $`\underset{¯}{z}^\mu =(z_1^\mu ,\mathrm{},z_d^\mu )`$ with $$\{\begin{array}{ccc}z_j^\mu =0& \text{if}& jI\\ z_j^\mu =1& \text{if}& jI.\end{array}$$ (9) Then there is a $`T_{\text{}}^d`$-equivariant biholomorphism given by $$\begin{array}{ccccc}T_{\text{}}^d& \times _{S_{\text{}}^\mu }& \stackrel{~}{V}_\mu & & \widehat{V}_\mu \\ [t& :& \underset{¯}{z}]& & t(\underset{¯}{z}+\underset{¯}{z}^\mu ),\end{array}$$ where $$\stackrel{~}{V}_\mu =\underset{iI}{}\text{}e_i\text{}^n$$ is the holomorphic slice at $`\underset{¯}{z}^\mu `$ for the $`T_{\text{}}^d`$-action on $`\text{}_\mathrm{\Delta }^d`$ and $$S_{\text{}}^\mu =\{(z_1,\mathrm{},z_d)T_{\text{}}^d\|z_j=1\text{if}jI\}$$ is the stabilizer of the action at $`\underset{¯}{z}^\mu `$. Notice that $`\stackrel{~}{V}_\mu =\text{s}_{\text{}}^\mu `$. We now prove that the subsets $`V_\mu =\widehat{V}_\mu /N_{\text{}}`$ are complex charts for the quotient space $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$. Consider the continuous surjective mapping $$\begin{array}{cccc}p_\mu :& \stackrel{~}{V}_\mu & & V_\mu \\ & \underset{¯}{z}& & [\underset{¯}{z}^\mu +\underset{¯}{z}].\end{array}$$ The discrete group $`\mathrm{\Gamma }_\mu =N_{\text{}}S_{\text{}}^\mu `$ acts on the set $`\stackrel{~}{V}_\mu `$ as follows: $$\begin{array}{cccccc}& \mathrm{\Gamma }_\mu & \times & \stackrel{~}{V}_\mu & & \stackrel{~}{V}_\mu \\ & (t& ,& \underset{¯}{z})& & t\underset{¯}{z}.\end{array}$$ This action is holomorphic, free on a connected dense open set and the mapping $`p_\mu `$ induces a continuous bijection $$\varphi _\mu :\stackrel{~}{V}_\mu /\mathrm{\Gamma }_\mu V_\mu .$$ To show that $`\varphi _\mu `$ is actually a homeomorphism, we only need to show that it is an open mapping. This amounts to showing that, given a $`\mathrm{\Gamma }_\mu `$-invariant open subset $`W`$ of $`\stackrel{~}{V}_\mu `$, then $`N_{\text{}}(\underset{¯}{z}^\mu +W)`$ is an open subset of $`\text{}_\mathrm{\Delta }^d`$. Since $`W`$ is $`\mathrm{\Gamma }_\mu `$-invariant, by Lemma 2.3, (ii), we have that $`N_{\text{}}(\underset{¯}{z}^\mu +W)=\mathrm{exp}(\text{n}_{\text{}})(\underset{¯}{z}^\mu +W)`$. Applying Lemma 2.3, (iii), we construct the surjective mapping $$\begin{array}{ccccc}W\hfill & \times & \hfill \mathrm{\Pi }_{jI}\text{}e_j& & \mathrm{exp}(\text{n}_{\text{}})(\underset{¯}{z}^\mu +W)\\ (\underset{¯}{w}\hfill & ,& \hfill V)& & \mathrm{exp}(V)\underset{¯}{z}^\mu +\mathrm{exp}(\pi _\mu ^1(\pi _{\text{}}(V)))\underset{¯}{w}.\end{array}$$ The determinant of its Jacobian matrix is given by $`_{j=1}^ne^{K_j(Y)}+4\pi ^2_{hI}e^{4\pi Y_h}`$, where $`Y=\frac{i}{2}(\overline{V}V)`$ and the $`K_j`$’s are real linear functions. This implies, by the inverse function theorem, that $`\mathrm{exp}(\text{n}_{\text{}})(\underset{¯}{z}^\mu +W)`$ is an open subset of $`\text{}_\mathrm{\Delta }^d`$. Let us now show that these charts are compatible. Consider two vertices of $`\mathrm{\Delta }`$, $`\mu `$ and $`\nu `$, and let $`I`$ and $`J`$ be the corresponding subsets of $`\{1,\mathrm{},d\}`$. Assume that the corresponding charts, $`V_\mu `$ and $`V_\nu `$, have non-empty intersection, and consider the two sets $`W_\mu =\varphi _\mu ^1(V_\mu V_\nu )`$ and $`W_\nu =\varphi _\nu ^1(V_\mu V_\nu )`$. We want to describe these two sets as submodels of $`\stackrel{~}{V}_\mu /\mathrm{\Gamma }_\mu `$ and $`\stackrel{~}{V}_\nu /\mathrm{\Gamma }_\nu `$ respectively, and show that $`g_{\mu \nu }=\varphi _\nu ^1\varphi _\mu `$ is a biholomorphism of models. Consider $`\stackrel{~}{W}_\mu =p_\mu ^1(\varphi _\mu ^1(V_\mu V_\nu ))\stackrel{~}{V}_\mu `$ and $`\stackrel{~}{W}_\nu =p_\nu ^1(\varphi _\nu ^1(V_\mu V_\nu ))\stackrel{~}{V}_\nu `$. Now, $$\stackrel{~}{W}_\mu =\left(\underset{jIJ}{}\text{}e_j\right)\times \left(\underset{jIIJ}{}\text{}^{}e_j\right)$$ and $$\stackrel{~}{W}_\nu =\left(\underset{jIJ}{}\text{}e_j\right)\times \left(\underset{jJIJ}{}\text{}^{}e_j\right).$$ Consider the universal covers $`W_\mu ^\mathrm{\#}`$, $`W_\nu ^\mathrm{\#}`$ of $`\stackrel{~}{W}_\mu `$, $`\stackrel{~}{W}_\nu `$ respectively. Notice that we have that $$W_\mu ^\mathrm{\#}=\left(\underset{jIJ}{}\text{}e_j\right)\times \left(\underset{jIIJ}{}\text{}e_j\right)\underset{jI}{}\text{}e_j=\stackrel{~}{V}_\mu $$ and $$W_\nu ^\mathrm{\#}=(\underset{jIJ}{}\text{}e_j=)\times \left(\underset{jJIJ}{}\text{}e_j\right)\underset{jJ}{}\text{}e_j=\stackrel{~}{V}_\nu ,$$ with projection maps $$\begin{array}{ccc}W_\mu ^\mathrm{\#}& & \stackrel{~}{W}_\mu \\ (\underset{¯}{z},\underset{¯}{\zeta })& & (\underset{¯}{z},\mathrm{exp}\underset{¯}{\zeta })\end{array}$$ and $$\begin{array}{ccc}W_\nu ^\mathrm{\#}& & \stackrel{~}{W}_\nu \\ (\underset{¯}{u},\underset{¯}{\eta })& & (\underset{¯}{u},\mathrm{exp}\underset{¯}{\eta }).\end{array}$$ Consider the discrete groups $$\mathrm{\Lambda }_\mu =\left\{(\mathrm{exp}Z,W)\right\|Z\underset{jIJ}{}\text{}e_j,W\underset{jIIJ}{}\text{}e_j,\pi _{\text{}}(Z+W)Q\}$$ and $$\mathrm{\Lambda }_\nu =\left\{(\mathrm{exp}U,V)\right\|U\underset{jIJ}{}\text{}e_j,V\underset{jJIJ}{}\text{}e_j,\pi _{\text{}}(U+V)Q\}$$ acting respectively on $`W_\mu ^\mathrm{\#}`$, $`W_\nu ^\mathrm{\#}`$ as follows $$\begin{array}{ccccc}(\mathrm{\Lambda }_\mu \hfill & ,& W_\mu ^\mathrm{\#})& & W_\mu ^\mathrm{\#}\\ ((\mathrm{exp}Z,W)\hfill & ,& (\underset{¯}{z},\underset{¯}{\zeta }))& & (\mathrm{exp}Z\underset{¯}{z},W+\underset{¯}{\zeta })\end{array}$$ $$\begin{array}{ccccc}(\mathrm{\Lambda }_\nu \hfill & ,& W_\nu ^\mathrm{\#})& & W_\nu ^\mathrm{\#}\\ ((\mathrm{exp}U,V)\hfill & ,& (\underset{¯}{u},\underset{¯}{\eta }))& & (\mathrm{exp}U\underset{¯}{u},V+\underset{¯}{\eta }).\end{array}$$ Notice that the projections induce homeomorphisms $`W_\mu ^\mathrm{\#}/\mathrm{\Lambda }_\mu W_\mu `$ and $`W_\nu ^\mathrm{\#}/\mathrm{\Lambda }_\nu W_\nu `$. Now we want to show that there is an equivariant biholomorphism $`g_{\mu \nu }^\mathrm{\#}:W_\mu ^\mathrm{\#}W_\nu ^\mathrm{\#}`$ that projects down to the mapping $`g_{\mu \nu }`$. Consider the isomorphisms $`\pi _\mu `$ and $`\pi _\nu `$ defined by (8). Notice that $`\pi _\nu ^1\pi _\mu `$ defines a biholomorphism from $`W_\mu ^\mathrm{\#}`$ to $`W_\nu ^\mathrm{\#}`$ that is equal to the identity on $`_{jIJ}\text{}e_j`$. Let moreover $`\rho `$, respectively $`\sigma `$, denote the projection of $`W_\nu ^\mathrm{\#}`$ onto the factor $`_{jIJ}\text{}e_j`$, respectively $`_{jJIJ}\text{}e_j`$. Define $$\begin{array}{cccc}g_{\mu \nu }^\mathrm{\#}:& W_\mu ^\mathrm{\#}& & W_\nu ^\mathrm{\#}\\ & (\underset{¯}{z},\underset{¯}{\zeta })& & (\mathrm{exp}\left(\rho \left(\pi _\nu ^1\pi _\mu (\underset{¯}{\zeta })\right)\right)\underset{¯}{z},\sigma \left(\pi _\nu ^1\pi _\mu (\underset{¯}{\zeta })\right)).\end{array}$$ It is straightforward to check that $`g_{\mu \nu }^\mathrm{\#}`$ is an equivariant biholomorphism that projects down to $`g_{\mu \nu }`$. Now complete this collection of charts with all the other compatible charts. This concludes the proof that $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$ is a complex quasifold of dimension $`n`$. The standard holomorphic action of $`T_{\text{}}^d`$ on $`\text{}_\mathrm{\Delta }^d`$ induces a holomorphic action of $`T_{\text{}}^d/N_{\text{}}`$ on $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$. The isomorphism (3) allows us to define a holomorphic action $`\tau :(\text{d}_{\text{}}/Q)\times (\text{}_\mathrm{\Delta }^d/N_{\text{}})\text{}_\mathrm{\Delta }^d/N_{\text{}}`$. To check that $`\tau `$ is holomorphic, consider, for each vertex $`\mu `$ of the polytope, the chart $`V_\mu `$, and the bijection $`\pi _\mu `$ given in (8). We then have that the following diagram commutes $$\begin{array}{ccc}\text{d}_{\text{}}\times \stackrel{~}{V}_\mu & \stackrel{\stackrel{~}{\tau }_\mu }{}& \stackrel{~}{V}_\mu \\ (W,\underset{¯}{z})& & \mathrm{exp}(\pi _\mu ^1(W))\underset{¯}{z}\\ \stackrel{}{}\stackrel{}{}& & \stackrel{}{}\stackrel{}{}\\ (\text{d}_{\text{}}/Q)\times (\stackrel{~}{V}_\mu /\mathrm{\Gamma }_\mu )& \stackrel{\tau _\mu }{}& \stackrel{~}{V}_\mu /\mathrm{\Gamma }_\mu \\ ([W],[\underset{¯}{z}])& & [\mathrm{exp}(\pi _\mu ^1(W))\underset{¯}{z}]\\ \stackrel{}{}\stackrel{}{}& & \stackrel{}{}\stackrel{}{}\\ (\text{d}_{\text{}}/Q)\times V_\mu & \stackrel{\tau }{}& V_\mu \\ ([W],[\underset{¯}{z}+\underset{¯}{z}_\mu ])& & [\mathrm{exp}(\pi _\mu ^1(W))(\underset{¯}{z}+\underset{¯}{z}_\mu )]=[\mathrm{exp}(\pi _\mu ^1(W))\underset{¯}{z}+\underset{¯}{z}_\mu ]\end{array}$$ and that $`\stackrel{~}{\tau }_\mu `$ is a holomorphic mapping. Notice finally that the dense open orbit for this action is given by $`\text{}_F^d/N_{\text{}}`$, where $`F=\text{Int}(\mathrm{\Delta })`$. $`\mathrm{}`$ ###### Remark 2.4 Suppose that the polytope $`\mathrm{\Delta }`$ is rational with respect to a lattice $`L\text{d}`$. Choose inward-pointing normal vectors to the facets of $`\mathrm{\Delta }`$ that are primitive in $`L`$, and take $`Q=L`$. Then the corresponding quotient $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$ is a complex orbifold and contains a dense open orbit for the action of the honest torus $`T_{\text{}}=\text{d}_{\text{}}/L`$; it is the usual toric variety associated to the polytope $`\mathrm{\Delta }`$. ###### Remark 2.5 Notice that, just like for honest toric varieties, the quasitorus $`\text{d}_{\text{}}/Q`$ is contained in $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$ as a dense open subset, and the action of $`\text{d}_{\text{}}/Q`$ on $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$ extends the standard action of $`\text{d}_{\text{}}/Q`$ on itself. Let us test this construction on three examples. Only the third is an example of a truly non-rational polytope. ###### Example 2.6 (The unit interval) Let us consider the polytope $`[0,1]\text{}^{}`$ and let us apply the above construction to the choice of vectors $`X_1=s`$, $`X_2=t`$, with $`s,t>0`$, $`s/t\text{}`$, and to the choice of quasilattice $`Q=s\text{}+t\text{}`$. It is easy to check that $`\text{}_\mathrm{\Delta }^2=\text{}^2\{0\}`$ and that $`N_{\text{}}=\{(e^{2\pi iZ},e^{2\pi i\frac{s}{t}Z})\|Z\text{}\}`$. We cover $`\text{}_\mathrm{\Delta }^2/N_{\text{}}`$ with two charts, one for each of its vertices, $`\mathrm{s}=\{0\}`$ and $`\mathrm{n}=\{1\}`$: $`V_\mathrm{s}=\{[z_1:z_2]\text{}_\mathrm{\Delta }^2/N_{\text{}}\|z_20\}`$ and $`V_\mathrm{n}=\{[z_1:z_2]\text{}_\mathrm{\Delta }^2/N_{\text{}}\|z_10\}`$. Consider now the the discrete group $`\mathrm{\Gamma }_\mathrm{s}=\text{}`$ acting on the set $`\stackrel{~}{V}_\mathrm{s}=\text{}`$ according to the rule $`(k,z)e^{2\pi i\frac{t}{s}k}z`$; this action is holomorphic and free on the connected, dense subset $`\stackrel{~}{V}_\mathrm{s}\{0\}`$ and the mapping $$\begin{array}{cccc}\varphi _\mathrm{s}:& \stackrel{~}{V}_\mathrm{s}/\mathrm{\Gamma }_\mathrm{s}& & V_\mathrm{s}\\ & [z]& & [z:1]\end{array}$$ is a homeomorphism. Similarly the group $`\mathrm{\Gamma }_\mathrm{n}=\text{}`$ acts on $`\stackrel{~}{V}_\mathrm{n}=\text{}`$ by the rule $`(m,z)e^{2\pi i\frac{s}{t}m}w`$; this action is holomorphic and free on the connected, dense subset $`\stackrel{~}{V}_\mathrm{n}\{0\}`$ and the mapping $$\begin{array}{cccc}\varphi _\mathrm{n}:& \stackrel{~}{V}_\mathrm{n}/\mathrm{\Gamma }_\mathrm{n}& & V_\mathrm{n}\\ & [w]& & [1:w]\end{array}$$ is a homeomorphism. Let us show that these two charts are compatible. The set $`\varphi _\mathrm{s}^1(V_\mathrm{s}V_\mathrm{n})`$ is a submodel of $`\stackrel{~}{V}_\mathrm{s}/\mathrm{\Gamma }_\mathrm{s}`$: it is the quotient of the space $`W_\mathrm{s}^\mathrm{\#}=\text{}`$ by the action of $`\mathrm{\Lambda }_\mathrm{s}=\text{}^2`$ given by $`(h,k)\zeta =\zeta +h+\frac{t}{s}k`$. Similarly the set $`\varphi _\mathrm{n}^1(V_\mathrm{s}V_\mathrm{n})`$ is the quotient of the space $`W_\mathrm{n}^\mathrm{\#}=\text{}`$ by the action of $`\mathrm{\Lambda }_\mathrm{n}=\text{}^2`$ given by $`(l,m)\eta =\eta +l+\frac{s}{t}m`$. The bijective mapping $`g_{\mathrm{sn}}=\varphi _\mathrm{n}^1\varphi _\mathrm{s}:\varphi _\mathrm{s}^1(V_\mathrm{s}V_\mathrm{n})`$ $``$ $`\varphi _\mathrm{n}^1(V_\mathrm{s}V_\mathrm{n})`$ $`\left[z\right]`$ $``$ $`\left[w=z^{\frac{s}{t}}\right]`$ is a biholomorphism of models: its lift is given by $`\zeta \eta =\frac{s}{t}\zeta `$. Now complete this collection with all the other compatible charts. ###### Remark 2.7 We can see from the previous example that complex quasifolds corresponding to the same polytope $`\mathrm{\Delta }`$ are not in general biholomorphic, in fact they are not even diffeomorphic. This was also visible in the rational case: the same construction applied to $`s`$, $`t`$ relatively prime integers yields in fact either ordinary or weighted projective space, which have different complex orbifold structures. This may appear to be in contradiction with Theorem 9.4 of Lerman-Tolman \[LT\], which implies that these spaces are biholomorphic - in reality their notion of biholomorphism is algebraic and, unlike ours, does not keep track of the orbifold structure. ###### Example 2.8 (The right triangle) Consider the right triangle in $`(\text{}^2)^{}`$ of vertices $`(0,0)`$, $`(s,0)`$ and $`(0,t)`$, where $`s,t`$ are two positive real numbers such that $`s/t\text{}`$. We apply the construction to the choice of inward-pointing normals $`X_1=(1,0)`$, $`X_2=(0,1)`$, $`X_3=(t,s)`$ and to the quasilattice $`Q=X_1\text{}+X_2\text{}+X_3\text{}`$. Then we have that $`\text{}_\mathrm{\Delta }^3=\text{}^3\{0\}`$ and $`N_{\text{}}=\{(e^{2\pi itZ},e^{2\pi isZ},e^{2\pi iZ})\|Z\text{}\}`$. We cover $`\text{}_\mathrm{\Delta }^3/N_{\text{}}`$ with three charts, one for each of its vertices: $`V_1=\{[z_1:z_2:z_3]\text{}_\mathrm{\Delta }^3/N_{\text{}}\|z_10\}`$, $`V_2=\{[z_1:z_2:z_3]\text{}_\mathrm{\Delta }^3/N_{\text{}}\|z_20\}`$ and $`V_3=\{[z_1:z_2:z_3]\text{}_\mathrm{\Delta }^3/N_{\text{}}\|z_30\}`$. The discrete group $`\mathrm{\Gamma }_1=\text{}`$ acts on $`\stackrel{~}{V}_1=\text{}^2`$ according to the rule $`h(z_2,z_3)=(e^{2\pi i\frac{s}{t}h}z_2,e^{2\pi i\frac{1}{t}h}z_3)`$ and the mapping $$\begin{array}{cccc}\varphi _1:& \stackrel{~}{V}_1/\mathrm{\Gamma }_1& & V_1\\ & [z_2,z_3]& & [1:z_2:z_3]\end{array}$$ is a homeomorphism. Similarly the group $`\mathrm{\Gamma }_2=\text{}`$ acts on $`\stackrel{~}{V}_2=\text{}^2`$ by the rule $`k(z_1,z_3)=(e^{2\pi i\frac{t}{s}k}z_1,e^{2\pi i\frac{1}{s}k}z_3)`$ and the mapping $$\begin{array}{cccc}\varphi _2:& \stackrel{~}{V}_2/\mathrm{\Gamma }_2& & V_2\\ & [z_1,z_3]& & [z_1:1:z_3]\end{array}$$ is a homeomorphism. Finally the group $`\mathrm{\Gamma }_3=\text{}`$ acts on $`\stackrel{~}{V}_3=\text{}^2`$ according to the rule $`l(z_1,z_2)=(e^{2\pi itl}z_1,e^{2\pi isl}z_2)`$ and the mapping $$\begin{array}{cccc}\varphi _3:& \stackrel{~}{V}_3/\mathrm{\Gamma }_3& & V_3\\ & [z_1,z_2]& & [z_1:z_2:1]\end{array}$$ is a homeomorphism. The changes of charts work as in the previous example. ###### Example 2.9 (The regular pentagon) Let us take the regular pentagon in $`(\text{}^2)^{}`$ of vertices $`(1,0,)`$, $`(a,b)`$, $`(c,d)`$, $`(c,d)`$ and $`(a,b)`$, where $`a=\mathrm{cos}\frac{2\pi }{5}`$, $`b=\mathrm{sin}\frac{2\pi }{5}`$, $`c=\mathrm{cos}\frac{4\pi }{5}`$, $`d=\mathrm{sin}\frac{4\pi }{5}`$. There exists no lattice $`L`$ with respect to which this simple convex polytope is rational. Let us apply the above construction to the choice of inward-pointing normal vectors $`X_1=(1,0)`$, $`X_2=(a,b)`$, $`X_3=(c,d)`$, $`X_4=(c,d)`$, $`X_5=(a,b)`$ and to the choice of quasilattice $`Q=_{j=1}^5X_j\text{}`$. Then we have that $$N_{\text{}}=\left\{(e^{2\pi iZ_1},e^{2\pi iZ_2},e^{2\pi iZ_3},e^{2\pi i[2a(Z_2Z_3)+Z_1)]},e^{2\pi i[2a(Z_2Z_1)+Z_3)]})\right\|(Z_1,Z_2,Z_3)\text{}^3\}.$$ We cover the quasifold $`\text{}_\mathrm{\Delta }^5/N_{\text{}}`$ with five charts: $$V_1=\{[z_1:z_2:z_3:z_4:z_5]\text{}_\mathrm{\Delta }^5/N_{\text{}}\|z_1,z_2,z_30\},$$ $$V_2=\{[z_1:z_2:z_3:z_4:z_5]\text{}_\mathrm{\Delta }^5/N_{\text{}}\|z_2,z_3,z_40\},$$ $$V_3=\{[z_1:z_2:z_3:z_4:z_5]\text{}_\mathrm{\Delta }^5/N_{\text{}}\|z_3,z_4,z_50\},$$ $$V_4=\{[z_1:z_2:z_3:z_4:z_5]\text{}_\mathrm{\Delta }^5/N_{\text{}}\|z_1,z_4,z_50\}$$ and $$V_5=\{[z_1:z_2:z_3:z_4:z_5]\text{}_\mathrm{\Delta }^5/N_{\text{}}\|z_1,z_2,z_50\}.$$ The discrete group $`\mathrm{\Gamma }_1=\text{}^3`$ acts on the set $`\stackrel{~}{V}_1=\text{}^2`$ according to the rule $`(h,k,l)(z_4,z_5)=(e^{4\pi ia(kl)}z_4,e^{4\pi ia(kh)}z_5)`$ and the mapping $$\begin{array}{cccc}\varphi _1:& \stackrel{~}{V}_1/\mathrm{\Gamma }_1& & V_1\\ & [z_4:z_5]& & [1:1:1:z_4:z_5]\end{array}$$ is a homeomorphism. We leave the discussion of the other four charts, and their mutual compatibility, to the reader. ## 3 Kähler structures In the previous section we associated to any simple convex polytope $`\mathrm{\Delta }\text{d}^{}`$, together with a choice of normal vectors $`X_j`$, $`j=1,\mathrm{},d`$, and of a quasilattice $`Q\text{d}`$ containing the $`X_j`$’s, a complex quasifold of dimension $`n`$. Starting with the same data it is possible to construct a symplectic quasifold of dimension $`2n`$, endowed with an effective Hamiltonian action of the quasitorus $`\text{d}/Q`$, such that the image of the corresponding moment mapping is exactly $`\mathrm{\Delta }`$ (see \[P, Theorem 3.3\]). The complex and symplectic quasifolds are both described as orbit spaces, the first is the quotient of $`\text{}_\mathrm{\Delta }^d`$ by $`N_{\text{}}`$, the second is a symplectic quotient with respect to the action of $`N`$ on $`\text{}^d`$. We want to prove that the complex and symplectic quotient can be identified, according to a general principle initiated by Kempf-Ness \[KN\] and later developed by Kirwan \[K\] and Ness \[N\]. More precisely, we shall prove that the two quotients are diffeomorphic and that the complex and symplectic structure are compatible, and thus define the structure of a Kähler quasifold (see Definition 1.11). To begin with, let us briefly recall from \[P\] the construction of the symplectic quotient that we are interested in. Consider the mapping $`J(\underset{¯}{z})=_{j=1}^d(\|z_j\|^2+\lambda _j)e_j^{}`$, where the $`\lambda _j`$’s are given in (2) and are uniquely determined by our choice of normal vectors. The mapping $`J`$ is a moment mapping for the standard action of $`T^d`$ on $`\text{}^d`$. Consider now the subgroup $`NT^d`$ and the corresponding inclusion of Lie algebras $`\iota :\text{n}\text{}^d`$. The mapping $`\mathrm{\Psi }:\text{}^d\text{n}^{}`$ given by $`\mathrm{\Psi }=\iota ^{}J`$ is a moment mapping for the induced action of $`N`$ on $`\text{}^d`$. Then the quotient space $`\mathrm{\Psi }^1(0)/N`$ is a compact symplectic quasifold of dimension $`2n`$; the quasitorus $`\text{d}/Q`$ acts on $`\mathrm{\Psi }^1(0)/N`$ in an effective, Hamiltonian fashion, and the image of the corresponding moment map is the polytope $`\mathrm{\Delta }`$. Let us now define a mapping between the symplectic and complex quotient. To define this mapping and to show that it is bijective we adapt to our setting the method described by Guillemin in \[G, Appendix 1\] for the smooth case. Before we can define this mapping, we need a preliminary lemma. ###### Lemma 3.1 The zero set $`\mathrm{\Psi }^1(0)`$ is contained in $`\text{}_\mathrm{\Delta }^d`$. Moreover, for any face $`F`$ of the polytope $`\mathrm{\Delta }`$, the orbit $`\text{}_F^d`$ intersects $`\mathrm{\Psi }^1(0)`$ in at least one point. Proof. We summarize the proof given in \[G, page 115\], which goes through without modification. Consider the exact sequence $$0\text{d}^{}\stackrel{\pi ^{}}{}(\text{}^d)^{}\stackrel{\iota ^{}}{}(\text{n})^{}0,$$ (10) where $`\pi `$ is the projection defined in (2). Notice that $`\underset{¯}{z}\mathrm{\Psi }^1(0)`$ if, and only if $`\iota ^{}(J(\underset{¯}{z}))=0`$. By (10) we have that, for a given $`\underset{¯}{z}\mathrm{\Psi }^1(0)`$, there exists a unique $`\zeta \text{d}^{}`$ such that $`J(\underset{¯}{z})=\pi ^{}(\zeta )`$. By making use of the explicit expression of $`J`$ we find $$\|z_j\|^2=\zeta ,X_j\lambda _j,j=1,\mathrm{},d.$$ (11) This implies that, given $`\underset{¯}{z}\mathrm{\Psi }^1(0)`$, the corresponding $`\zeta `$ lies in $`\mathrm{\Delta }`$ and that, again by (11), $`\underset{¯}{z}`$ is in $`\text{}_F^d`$, where $`F`$ is the face of $`\mathrm{\Delta }`$ containing $`\zeta `$ in its interior. Similarly, given any face $`F`$ of the polytope, we can always find a point $`\underset{¯}{z}\mathrm{\Psi }^1(0)\text{}_F^d`$. $`\mathrm{}`$ By Lemma 3.1, there is an injection $$:\mathrm{\Psi }^1(0)\text{}_\mathrm{\Delta }^d,$$ (12) which induces the mapping $$\chi :\mathrm{\Psi }^1(0)/N\text{}_\mathrm{\Delta }^d/N_{\text{}}$$ (13) that sends an $`N`$-orbit to the corresponding $`N_{\text{}}`$-orbit. This mapping is equivariant with respect to the actions of the quasitori $`\text{d}/Q`$ and $`\text{d}_{\text{}}/Q`$. We are now ready to state the main result of this section. ###### Theorem 3.2 Let d be a vector space of dimension $`n`$, and let $`\mathrm{\Delta }\text{d}^{}`$ be a simple convex polytope. Choose inward-pointing normals to the facets of $`\mathrm{\Delta }`$, $`X_1,\mathrm{},\text{X}_d\text{d}`$, and let $`Q`$ be a quasilattice containing these vectors. Then the mapping $$\chi :\mathrm{\Psi }^1(0)/N\text{}_\mathrm{\Delta }^d/N_{\text{}}$$ is an equivariant diffeomorphism of quasifolds. Moreover the induced symplectic form on the complex quasifold $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$ is Kähler. For the definitions of symplectic quasifold, diffeomorphism of quasifolds, and pullback of a differential form, we refer the reader to \[P\]. Before we can proceed with the proof of this theorem we need a number of remarks and lemmas. ###### Remark 3.3 The group $`A=\mathrm{exp}(i\text{n})`$ moves any point $`\underset{¯}{z}\mathrm{\Psi }^1(0)`$ out of $`\mathrm{\Psi }^1(0)`$. To see this, remark that, for any vector $`Y\text{n}`$, the corresponding function $`\mathrm{\Psi }_Y`$ satisfies $`d\mathrm{\Psi }_Y=ı_{\stackrel{~}{Y}}\omega _0`$, where $`\omega _0=\frac{1}{2\pi i}_{j=1}^ddz_jd\overline{z}_j`$ is the standard symplectic form of $`\text{}^d`$ and $`\stackrel{~}{Y}`$ is the vector field generated by $`Y`$. This implies that the gradient of the function $`\mathrm{\Psi }_Y`$ is the vector field generated by $`iY`$, which by (6) does not vanish in $`\text{}_\mathrm{\Delta }^d`$. The function $`\mathrm{\Psi }_Y`$ is therefore strictly increasing along its gradient flow, hence $`\mathrm{exp}(tiY)\underset{¯}{z}\mathrm{\Psi }^1(0)`$, for all $`Y\text{n}`$ and all $`t0`$. ###### Lemma 3.4 The following facts are equivalent: (i) the mapping $`\chi `$ is bijective; (ii) every $`N_{\text{}}`$-orbit in $`\text{}_\mathrm{\Delta }^d`$ intersects $`\mathrm{\Psi }^1(0)`$ in an $`N`$-orbit; (iii) every $`A`$-orbit in $`\text{}_\mathrm{\Delta }^d`$ intersects $`\mathrm{\Psi }^1(0)`$ in at least one point. Proof. It is obvious that (i) is equivalent to (ii). Before going on with the proof let us re-state the second and third point: (ii) for each $`N_{\text{}}`$-orbit through a point $`\underset{¯}{z}\text{}_\mathrm{\Delta }^d`$ there exists a point $`\underset{¯}{u}\mathrm{\Psi }^1(0)`$ such that $`(N_{\text{}}\underset{¯}{z})\mathrm{\Psi }^1(0)=N\underset{¯}{u}`$; (iii) for each $`A`$-orbit through a point $`\underset{¯}{z}\text{}_\mathrm{\Delta }^d`$ there exists at least one point $`\underset{¯}{w}(A\underset{¯}{z})\mathrm{\Psi }^1(0)`$. The proof is based on three key facts: the polar decomposition for the group $`N_{\text{}}`$ (see Remark 2.1), the invariance of $`\mathrm{\Psi }^1(0)`$ under the action of $`N`$, and Remark 3.3. To see that (ii) implies (iii), consider an $`A`$-orbit through $`\underset{¯}{z}\text{}_\mathrm{\Delta }^d`$. Then, by (ii), there exists $`\underset{¯}{u}\mathrm{\Psi }^1(0)`$ such that $`(N_{\text{}}\underset{¯}{z})\mathrm{\Psi }^1(0)=N\underset{¯}{u}`$. Since $`N_{\text{}}=NA`$, there exists $`\underset{¯}{w}N\underset{¯}{u}\mathrm{\Psi }^1(0)`$ such that $`\underset{¯}{w}A\underset{¯}{z}`$. To see that (iii) implies (ii), consider any $`N_{\text{}}`$-orbit through $`\underset{¯}{z}`$. By assumption, the orbit $`A\underset{¯}{z}`$ intersects $`\mathrm{\Psi }^1(0)`$ in at least one point $`\underset{¯}{w}`$. Since $`N_{\text{}}=NA`$ we have that $`N_{\text{}}\underset{¯}{z}=N_{\text{}}\underset{¯}{w}`$, that $`N\underset{¯}{w}(N_{\text{}}\underset{¯}{z})\mathrm{\Psi }^1(0)`$, and that this inclusion is in fact an equality by Remark 3.3. $`\mathrm{}`$ ###### Remark 3.5 If there exists a point $`\underset{¯}{w}(A\underset{¯}{z})\mathrm{\Psi }^1(0)`$, then $`\underset{¯}{w}`$ is uniquely determined by this property. If $`\underset{¯}{w}`$ and $`\underset{¯}{w}^{}`$ are in $`(A\underset{¯}{z})\mathrm{\Psi }^1(0)`$, then there exists an $`aA`$ such that $`a\underset{¯}{w}=\underset{¯}{w}^{}`$. By Remark 3.3 we have that $`a`$ is the identity and that $`\underset{¯}{w}=\underset{¯}{w}^{}`$. The following result will be essential to proving that the mapping $`\chi `$ is a bijection. ###### Lemma 3.6 The moment mapping $`\mathrm{\Psi }`$ maps any $`A`$-orbit diffeomorphically onto an open convex cone in $`\text{n}^{}`$; moreover, if two $`A`$-orbits lie in the same $`T_{\text{}}^d`$-orbit $`\text{}_F^d\text{}_\mathrm{\Delta }^d`$, then their images with respect to $`\mathrm{\Psi }`$ are identical. Proof. The proof given in \[G, Appendix 1: Theorem 2.1 and Theorem 2.2\] applies, the only delicate point here is to notice that, by (6), $`A`$ still acts freely on $`\text{}_\mathrm{\Delta }^d`$. Notice also that, by (7), this is certainly not true for the action of $`N`$. We now want to study the image by $`\mathrm{\Psi }`$ of an $`A`$-orbit through a point $`\underset{¯}{z}\text{}_\mathrm{\Delta }^d`$, namely $$\mathrm{\Psi }(A\underset{¯}{z})=\left\{\mathrm{\Psi }(e^{2\pi \alpha _1(X)}z_1,\mathrm{},e^{2\pi \alpha _d(X)}z_d)\right\|X\text{n}\},$$ where $`\alpha _j=\iota ^{}(e_j^{})`$, $`j=1,\mathrm{},d`$. Since $`A`$ acts freely on $`\text{}_\mathrm{\Delta }^d`$, the exponential mapping defines a diffeomorphism between n and $`Az`$, and we can identify the $`A`$-orbit with n. Therefore the set we are interested in is the image of the mapping $$\begin{array}{cccc}f:& \text{n}& & \text{n}^{}\\ & X& & \mathrm{\Psi }(e^{2\pi \alpha _1(X)}z_1,\mathrm{},e^{2\pi \alpha _d(X)}z_d).\end{array}$$ The point $`\underset{¯}{z}`$ lies in $`\text{}_F^d`$ for a face $`F`$ of the polytope $`\mathrm{\Delta }`$. Let $`I`$ be the corresponding set of indices, as defined in Section 2, then $$f(X)=\underset{jI}{}e^{4\pi \alpha _j(X)}\|z_j\|^2\alpha _j+\lambda _j\alpha _j.$$ Now we only have to check that $$\text{span}\{\alpha _j\text{n}^{}\|jI\}=\text{n}^{};$$ (14) this is equivalent to saying that there does not exist a non zero $`X\text{n}`$ such that $`\alpha _j(X)=0`$ for all $`jI`$. Indeed, if there was such an $`X`$, it would lie in $`\text{n}\text{s}_{\text{}}^F`$ which is $`\{0\}`$. Using (14) one can prove, following \[G\], that $`f`$ maps n diffeomorphically into an open convex cone of $`\text{n}^{}`$, and also that the image of $`f`$ depends only on $`I`$. $`\mathrm{}`$ We are now ready to proceed with the proof of the main result of this section. Proof of Theorem 3.2. The proof is divided into several steps. We first prove that the mapping $`\chi `$ is bijective and continuous. Then we lift $`\chi `$ locally to prove that it is a diffeomorphism, and finally we show that the pull-back via $`\chi ^1`$ of the symplectic form of $`\mathrm{\Psi }^1(0)/N`$, is Kähler on the complex quasifold $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$. In order to prove that $`\chi `$ is bijective we will show that Lemma 3.6 implies that, for a given $`A`$-orbit through a point $`\underset{¯}{z}\text{}_\mathrm{\Delta }^d`$, there exists a point $`\underset{¯}{w}\mathrm{\Psi }^1(0)`$ such that $`(A\underset{¯}{z})\mathrm{\Psi }^1(0)=\{\underset{¯}{w}\}`$. Then we will apply Lemma 3.4 to conclude. To prove the existence of $`\underset{¯}{w}`$, consider the $`T_{\text{}}^d`$-orbit $`\text{}_F^d`$ containing $`A\underset{¯}{z}`$. By Lemma 3.1 this $`\text{}_F^d`$ intersects $`\mathrm{\Psi }^1(0)`$ in at least one point, $`\underset{¯}{u}`$. On the other hand $`\mathrm{\Psi }(A\underset{¯}{z})=\mathrm{\Psi }(A\underset{¯}{u})\text{n}^{}`$ by Lemma 3.6, and the common image contains $`0=\mathrm{\Psi }(\underset{¯}{u})`$. So there exists $`aA`$ such that $`\mathrm{\Psi }(a\underset{¯}{z})=0`$; then $`\underset{¯}{w}=a\underset{¯}{z}(A\underset{¯}{z})\mathrm{\Psi }^1(0)`$. Notice that, by Remark 3.5, the point $`\underset{¯}{w}`$ is uniquely determined. Continuity of the mapping $`\chi `$ is implied by that of the immersion $``$. For the next step we need a covering of the symplectic quotient by a collection of its charts. We refer the reader to \[P, Theorem 3.1\] for the construction of the full atlas. Let $`\mu `$ be a vertex of the polytope $`\mathrm{\Delta }`$, and let $`I=\{r_1,\mathrm{},r_n\}\{1,\mathrm{},d\}`$ denote the corresponding subset of indices. For each vertex of $`\mathrm{\Delta }`$ we want to construct a corresponding chart. Let $`(a_{jh})M_{(n,d)}(\text{})`$ denote the matrix of the mapping $`\pi :\text{}^d\text{d}`$ with respect to the basis $`\{X_j\|jI\}`$ of d and the standard basis of $`\text{}^d`$. Consider the open subset of $`\stackrel{~}{V}_\mu `$ defined as follows $$\stackrel{~}{U}_\mu =\left\{\underset{¯}{w}\stackrel{~}{V}_\mu \right\|\underset{j=1}{\overset{n}{}}a_{jh}\left(\|w_{r_j}\|^2+\lambda _{r_j}\right)\lambda _h>0,hI\}.$$ The group $`\mathrm{\Gamma }_\mu `$ defined in Lemma 2.3 acts on $`\stackrel{~}{U}_\mu `$. Consider, for a given element $`\underset{¯}{w}\stackrel{~}{U}_\mu `$, the element $`\underset{¯}{w}^\mu =(w_1^\mu ,\mathrm{},w_d^\mu )\text{}^d`$ defined as follows $$\{\begin{array}{ccc}w_h^\mu =0\hfill & \text{if}& hI\\ w_h^\mu =\sqrt{_{j=1}^na_{jh}\left(\|w_{r_j}\|^2+\lambda _{r_j}\right)\lambda _h}\hfill & \text{if}& hI.\end{array}$$ Notice that $`\underset{¯}{w}+\underset{¯}{w}^\mu `$ belongs to $`\mathrm{\Psi }^1(0)`$: define $`\nu \text{d}^{}`$ so that $`<\nu ,X_j>=\|w_j\|^2+\lambda _j`$, for all $`jI`$; then it is easy to check that $`J(\underset{¯}{w}+\underset{¯}{w}^\mu )=\pi ^{}(\nu )`$ and $`\nu \mathrm{\Delta }`$. Consider now the open sets $`\widehat{U}_\mu =\widehat{V}_\mu \mathrm{\Psi }^1(0)\mathrm{\Psi }^1(0)`$ and $`U_\mu =\widehat{U}_\mu /N\mathrm{\Psi }^1(0)/N`$. Then the surjective mapping $$\begin{array}{cccc}q_\mu :& \stackrel{~}{U}_\mu & & U_\mu \\ & \underset{¯}{w}& & [\underset{¯}{w}^\mu +\underset{¯}{w}]\end{array}$$ induces a homeomorphism $$\begin{array}{cccc}\psi _\mu :& \stackrel{~}{U}_\mu /\mathrm{\Gamma }_\mu & & U_\mu \\ & [\underset{¯}{w}]& & [\underset{¯}{w}^\mu +\underset{¯}{w}].\end{array}$$ The above data define a chart, and the union of the $`U_\mu `$’s, for $`\mu `$ ranging over all the vertices of $`\mathrm{\Delta }`$, cover the symplectic quotient. Let us now show that the mapping $`\chi `$ lifts to a diffeomorphism $`\stackrel{~}{\chi }_\mu :\stackrel{~}{U}_\mu \stackrel{~}{V}_\mu `$ for each vertex $`\mu `$. For each $`\underset{¯}{w}\stackrel{~}{U}_\mu `$, there exists a unique element $`a(\underset{¯}{w})A`$ such that $`a(\underset{¯}{w})(\underset{¯}{w}+\underset{¯}{w}^\mu )`$ is of the form $`\underset{¯}{z}+\underset{¯}{z}^\mu `$, where $`\underset{¯}{z}\stackrel{~}{V}_\mu `$, and $`\underset{¯}{z}^\mu `$ is given by (9). Then we define $`\stackrel{~}{\chi }_\mu (\underset{¯}{w})=\underset{¯}{z}`$. We compute $`\underset{¯}{z}`$ explicitly following Lemma 2.3: $`\underset{¯}{z}=\stackrel{~}{\chi }_\mu (\underset{¯}{w})=\mathrm{exp}(i\pi _\mu ^1(\pi _{\text{}}(C(\underset{¯}{w})))\underset{¯}{w}`$, where the $`h`$-component of $`C(w)_{hJ}\text{}e_h`$ is given by $`\frac{1}{4\pi }\mathrm{log}A_h`$, with $`A_h=_{j=1}^na_{jh}(\|w_{r_j}\|^2+\lambda _{r_j})\lambda _h`$. Notice that the mapping $`\stackrel{~}{\chi }_\mu `$ thus defined is equivariant with respect to the action of $`\mathrm{\Gamma }_\mu `$ on $`\stackrel{~}{U}_\mu `$ and $`\stackrel{~}{V}_\mu `$, and it is a lift of $`\chi `$. It is now straightforward to deduce the bijectivity of $`\stackrel{~}{\chi }_\mu `$ from that of $`\chi `$, using the explicit expression for $`\stackrel{~}{\chi }_\mu `$. In order to prove that $`\stackrel{~}{\chi }_\mu `$ is a diffeomorphism we apply the inverse function theorem. Therefore we only have to check that its Jacobian matrix $`D(\stackrel{~}{\chi }_\mu )`$ is non-degenerate on $`\stackrel{~}{U}_\mu `$. Let $`(\underset{¯}{x},\underset{¯}{y})`$ be real coordinates in $`\stackrel{~}{U}_\mu `$, it turns out that the Jacobian matrix has the following form $$D(\stackrel{~}{\chi }_\mu )(\underset{¯}{x},\underset{¯}{y})=e^r(I_{2n}+\left(\begin{array}{c}M\\ N\end{array}\right)(^tM,^tN)),$$ where $`r`$ is the function $`_{hI}a_{kh}\frac{1}{2}\mathrm{log}A_h`$, and where $`M,N`$ are $`(n,dn)`$ matrices with entries $`u_ka_{kh}\frac{1}{\sqrt{A_h}}`$ and $`v_ka_{kh}\frac{1}{\sqrt{A_h}}`$ respectively. This implies that $`D(\stackrel{~}{\chi }_\mu )(w)`$ is positive definite for every $`\underset{¯}{w}\stackrel{~}{U}_\mu `$. To conclude that $`\chi `$ is a diffeomorphism observe that the continuity of the equivariant mapping $`\stackrel{~}{\chi }_\mu ^1`$ for each vertex $`\mu `$ implies that $`\chi ^1`$ is continuous, since $`\varphi _\mu `$ and $`\psi _\mu `$ are homeomorphisms. Now, having proved that $`\chi `$ is a diffeomorphism, we can consider the complex quotient endowed with a symplectic and a complex structure. We want to prove that the symplectic form is Kähler. This can be checked pointwise: let $`\underset{¯}{\overset{^}{w}}=\underset{¯}{w}+\underset{¯}{w}^\mu `$ be a point in $`\widehat{U}_\mu `$. Then the $`N`$-orbit through $`\underset{¯}{\overset{^}{w}}`$ is contained in $`\mathrm{\Psi }^1(0)`$ and the $`A`$-orbit through $`\underset{¯}{\overset{^}{w}}`$ is orthogonal to $`\mathrm{\Psi }^1(0)`$. This gives the isomorphism $$T_{\underset{¯}{\overset{^}{w}}}(\mathrm{\Psi }^1(0))/T_{\underset{¯}{\overset{^}{w}}}(N\underset{¯}{\overset{^}{w}})\text{}^d/T_{\underset{¯}{\overset{^}{w}}}(N_{\text{}}\underset{¯}{\overset{^}{w}}).$$ To conclude the proof we only have to remark that the symplectic structure on $`\mathrm{\Psi }^1(0)/N`$ and the complex structure on $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$, read on $`T_{\underset{¯}{\overset{^}{w}}}(\mathrm{\Psi }^1(0))/T_{\underset{¯}{\overset{^}{w}}}(N\underset{¯}{\overset{^}{w}})`$, are exactly the ones induced by the standard complex and symplectic structures of $`\text{}^d`$. $`\mathrm{}`$ ###### Remark 3.7 Consider the complex quasifold $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$ constructed in Theorem 2.2. Notice that, by varying the coefficients $`\lambda _j`$ in (2), we can produce many simple convex polytopes, each allowing the same choice of inward-pointing normals $`X_j`$ and of quasilattice $`Q`$ that we had made for $`\mathrm{\Delta }`$ (for example we can “inflate” $`\mathrm{\Delta }`$). The corresponding complex quasifolds are therefore exactly the same, but on the symplectic side in general one obtains non-equivalent symplectic structures. In other terms we obtain many non-isometric Kähler structures on the same complex quasifold $`\text{}_\mathrm{\Delta }^d/N_{\text{}}`$. We conclude this section with the discussion of an example. The unit interval provides a very useful model for a thorough understanding of the diffeomorphism $`\chi `$ and a neat example of Kähler quasifold. ###### Example 3.8 Consider the polytope $`[0,1]\text{}^{}`$, with the same choice of vectors and quasilattice made in Example 2.6 (we refer to that example for the notation). We can associate to these data the complex quasifold $`\text{}_\mathrm{\Delta }^2/N_{\text{}}`$ and also a symplectic quasifold, the quasisphere of \[P, Examples 1.13, 1.19, 2.10 and 3.5\], of which we recall only some features. Consider the symplectic quasifold $`\mathrm{\Psi }^1(0)/N`$, where $$\mathrm{\Psi }^1(0)=\{(z,w)\text{}^2\|t\|z_1\|^2+s\|z_2\|^2=st\},$$ and $$N=\{(e^{2\pi iX},e^{2\pi i\frac{s}{t}X})\|X\text{}\}.$$ The space $`\mathrm{\Psi }^1(0)/N`$ is covered by two charts $$U_\mathrm{s}=\{[z_1:z_2]\mathrm{\Psi }^1(0)/N\|z_20\}\text{and}U_\mathrm{n}=\{[z_1:z_2]\mathrm{\Psi }^1(0)/N\|z_10\}.$$ The corresponding local models are defined by $`\stackrel{~}{U}_\mathrm{s}=\{z\text{}\|\|z\|<\sqrt{s}\}\stackrel{~}{V}_\mathrm{s}`$ acted on by $`\mathrm{\Gamma }_\mathrm{s}`$ and by $`\stackrel{~}{U}_\mathrm{n}=\{z\text{}\|\|z\|<\sqrt{t}\}\stackrel{~}{V}_\mathrm{n}`$ acted on by $`\mathrm{\Gamma }_\mathrm{n}`$. To complete the picture we need the homeomorphisms: $$\begin{array}{cccc}\mathrm{\Phi }_\mathrm{s}:& \stackrel{~}{U}_\mathrm{s}/\mathrm{\Gamma }_\mathrm{s}& & U_\mathrm{s}\\ & [z]& & \left[z:\sqrt{t\frac{t}{s}\|z\|^2}\right]\end{array},\begin{array}{cccc}\mathrm{\Phi }_\mathrm{n}:& \stackrel{~}{U}_\mathrm{n}/\mathrm{\Gamma }_\mathrm{n}& & U_\mathrm{n}\\ & [w]& & \left[\sqrt{s\frac{s}{t}\|w\|^2}:w\right]\end{array}.$$ Now we are ready to write the local lifts $`\stackrel{~}{\chi }_\mathrm{s}`$ and $`\stackrel{~}{\chi }_\mathrm{n}`$ of the diffeomorphism $`\chi `$. Notice that $$(\chi \mathrm{\Phi }_\mathrm{s})([z])=[z:\sqrt{t\frac{t}{s}\|z\|^2}]_N_{\text{}}=[z(t\frac{t}{s}\|z\|^2)^{\frac{s}{2t}}:1]_N_{\text{}}.$$ Therefore a local lift $`\stackrel{~}{\chi }_\mathrm{s}:\stackrel{~}{U}_\mathrm{s}\stackrel{~}{V}_\mathrm{s}`$ is given by the equivariant diffeomorphism $$\stackrel{~}{\chi }_\mathrm{s}(z)=z\left(t\frac{t}{s}\|z\|^2\right)^{\frac{s}{2t}}$$ Analogously a local lift $`\stackrel{~}{\chi }_\mathrm{n}:\stackrel{~}{U}_\mathrm{n}\stackrel{~}{V}_\mathrm{n}`$ is given by the equivariant diffeomorphism $$\stackrel{~}{\chi }_\mathrm{n}(w)=w\left(s\frac{s}{t}\|w\|^2\right)^{\frac{t}{2s}}.$$ We exhibit now the complex quotient as a Kähler quasifold. We define the Kähler form $`\omega `$ by giving its local lifts on $`\stackrel{~}{V}_\mathrm{s}`$ and $`\stackrel{~}{V}_\mathrm{n}`$: $$\stackrel{~}{\omega }_\mathrm{s}=\frac{1}{2\pi i}\frac{s}{t}\frac{1}{\left(\frac{1}{s}+\frac{1}{t}\|z\|^{(1+\frac{s}{t})}\right)^2}dzd\overline{z},\stackrel{~}{\omega }_\mathrm{n}=\frac{1}{2\pi i}\frac{t}{s}\frac{1}{\left(\frac{1}{t}+\frac{1}{s}\|w\|^{(1+\frac{t}{s})}\right)^2}dwd\overline{w}.$$ (15) They are invariant under the action of $`\mathrm{\Gamma }_\mathrm{s}`$ and $`\mathrm{\Gamma }_\mathrm{n}`$ respectively and it is a straightforward computation to check that $`\stackrel{~}{\chi }_\mathrm{s}^{}\stackrel{~}{\omega }_\mathrm{s}=\frac{1}{2\pi i}dzd\overline{z}`$ and $`\stackrel{~}{\chi }_\mathrm{n}^{}\stackrel{~}{\omega }_\mathrm{n}=\frac{1}{2\pi i}dwd\overline{w}`$. On the other hand these are precisely the local lifts of the symplectic form on $`\stackrel{~}{U}_\mathrm{s}`$ and $`\stackrel{~}{U}_\mathrm{n}`$ respectively, so $`\omega `$ is the pullback via $`\chi ^1`$ of the symplectic form on $`\mathrm{\Psi }^1(0)/N`$ and hence defines a Kähler form on the complex quotient. To conclude the discussion of this significant example, it is worthwhile to check directly that the local forms given in (15) fulfill the definition of differential form. We have to show that they behave correctly under the change of charts. The first move is to pull them back to $`W_\mathrm{s}^\mathrm{\#}`$ and to $`W_\mathrm{n}^\mathrm{\#}`$, which are both equal to . We obtain $$\omega _\mathrm{s}^\mathrm{\#}=\frac{1}{2\pi i}\frac{s}{t}\frac{e^{(\zeta +\overline{\zeta })}}{\left(\frac{1}{s}+\frac{1}{t}e^{(\zeta +\overline{\zeta })\frac{1}{2}(1+\frac{s}{t})}\right)^2}d\zeta d\overline{\zeta },\omega _\mathrm{n}^\mathrm{\#}=\frac{1}{2\pi i}\frac{t}{s}\frac{1}{\left(\frac{1}{t}+\frac{1}{s}e^{(\eta +\overline{\eta })\frac{1}{2}(1+\frac{t}{s})}\right)^2}d\eta d\overline{\eta }$$ which are invariant under the respective action of $`\mathrm{\Lambda }_\mathrm{s}`$, $`\mathrm{\Lambda }_\mathrm{n}`$. It is easy to check that the mapping $`\zeta \eta =\frac{s}{t}\zeta `$ pulls $`\omega _\mathrm{n}^\mathrm{\#}`$ back to $`\omega _\mathrm{s}^\mathrm{\#}`$. Dipartimento di Matematica Applicata ”G. Sansone”, Via S. Marta 3, 50139 Firenze, ITALY, mailto:fiamma@dma.unifi.it and Laboratoire Dieudonné, Université de Nice, Parc Valrose, 06108 Nice Cedex 2, FRANCE, mailto:elisa@alum.mit.edu
warning/0004/cond-mat0004009.html	ar5iv	text	# Quantum transport in the cylindrical nanosize silicon-based MOSFET ## I Introduction During the last decade significant progress has been achieved in the scaling of the metal-oxide-semiconductor field-effect transistor (MOSFET) down to semiconductor devices with nanometer sizes. In Ref. , the fabrication of silicon quantum wires with lengths and widths of about 60 nm and $``$ 20 nm respectively is reported. The conductance of those quantum wires was measured for a wide range of temperatures, from 25 to 160 K. The fabrication and the investigation of a 40 nm gate length n-MOSFET are reported in Ref. . The resulting nanosize n-MOSFET operates rather normally at room temperature. Using nanoimprint lithography, a field effect transistor (FET) with a 100 nm wire channel was fabricated and the characteristics of this FET at room temperature were investigated. However, the small size of nanoscale MOSFET with a wide Si substrate negatively influences the device characteristics due to the floating body effect. Short channel effects together with random effects in the silicon substrate are very well known to cause a degradation of the threshold voltage and the appearance of uncontrollable charge and current in regions far from the gate electrode. Therefore in a nanoscale conventional MOSFET, the controlling ability of the gate electrode is substantially weakened. Recently, considerable attention has been paid to SOI (Si-on-insulator) MOSFETs, which are prospective for creating new nanosize devices. In Refs. , a MOSFET with very thin SOI was theoretically investigated on the basis of a 2D analytical model, while in Ref. , a 1D model was used. Drain-induced barrier lowering was considered and the physical mechanisms which determine the subthreshold slope (S-factor) were analyzed . As a result, a substantial reduction of the short-channel effect in the SOI MOSFET as compared to that in the bulk devices was established. As was shown theoretically in Ref. , the use of a lightly doped source and drain leads to an increase of the effective channel length what allows one to weaken the drain-induced barrier lowering. In an SOI MOSFET with a Si-Ge source , an improved drain-to-source breakdown voltage is achieved due to the absorption of excess holes in the channel region. In a transistor device with a channel sandwiched between oxide layers (dual-gate MOSFET), the floating body effects are significantly suppressed . A theoretical model of a dual-gate device is described in Ref. . In Ref. we have investigated the thermal equilibrium state of a nanoscale cylindrical silicon-based MOSFET device with a close gate electrode (MOSFETCGE). An advantage of the latter is the complete suppression of the floating body effect caused by external influences. Moreover, the short-channel effect in these devices can be even weaker than that in a dual-gate structure. The main goal of the present work is the investigation of quantum transport in a nanosize MOSFETCGE device. We have developed a flexible 2D model which optimally combines analytical and numerical methods and describes the main features of the MOSFETCGE device. The theoretical modeling of the quantum transport features involves the use of the Wigner distribution function formalism . The paper is organized as follows. In Section II, a description of the system is presented in terms of a one-electron Hamiltonian. In Section III the quantum Liouville equation satisfied by the electron density matrix is transformed into a set of one dimensional equations for partial Wigner distribution functions. A one-dimensional collision term is derived in Section IV. In Section V we describe a numerical model to solve the equations which have been derived for the partial Wigner distribution function. In Section VI the results of the numerical calculations are discussed. Finally, in Section VII we give a summary of our results and conclusions about the influence of the scattering processes in nanosize MOSFETs. ## II The Hamiltonian of the system We consider a cylindrical nanosize MOSFETCGE structure (Fig. 1) described by cylindrical coordinates $`(r,\varphi ,z)`$, where the $`z`$-axis is chosen to be the symmetry axis. In the semiconductor pillar the electron motion is determined by the following Hamiltonian $$\widehat{H}_j=\frac{\mathrm{}^2}{2m_j^{}}\frac{^2}{^2𝐫_{}}\frac{\mathrm{}^2}{2m_j^{}}\frac{^2}{^2z}+V(𝐫),$$ (1) where $`V(𝐫)=V_b(𝐫)+V_e(𝐫)`$ is the potential energy associated with the energy barrier and the electrostatic field, respectively; $`m_j^{}`$ and $`m_j^{}`$ are the effective masses of the transverse (in ($`x,y`$)-plane) and longitudinal (along $`z`$-axis) motion of an electron of the $`j`$-th valley. The electrostatic potential energy $`V_e(𝐫)`$ satisfies Poisson’s equation $$\mathrm{\Delta }V_e(𝐫)=\frac{e^2}{\epsilon _0\epsilon _i}\left(n(𝐫)+N_D(𝐫)N_A(𝐫)\right),i=1,2,$$ (2) where $`\epsilon _1`$ and $`\epsilon _2`$ are the dielectric constants of the semiconductor and oxide layers, respectively; $`n(𝐫)`$, $`N_D(𝐫)`$and $`N_A(𝐫)`$ are the concentrations of electrons, donors and acceptors respectively. In our calculations we assume that the source electrode is grounded whereas the potentials at the drain and gate electrodes are equal $`V_{ds}`$ and $`V_G`$, respectively. The study of the charge distribution in the cylindrical nanosize MOSFETCGE structure in the state of the thermodynamical equilibrium (see Ref. ) has shown that the concentration of holes is much lower than that of electrons so that electron transport is found to provide the main contribution to the current flowing through the MOSFET. For that reason, holes are neglected in the present transport calculations. ## III The Liouville equation In this section we consider ballistic transport of electrons. Neglecting scattering processes and inter-valley transitions in the conduction band, the one-electron density matrix can be written as $$\rho (𝐫,𝐫^{})=\underset{j}{}\rho _j(𝐫,𝐫^{}),$$ (3) where $`\rho _j(𝐫,𝐫^{})`$ is the density matrix of electrons residing in the $`j`$-th valley satisfying Liouville’s equation $$i\mathrm{}\frac{\rho _j}{t}=[H_j,\rho _j].$$ (4) In order to impose reasonable boundary conditions for the density matrix in the electrodes, it is convenient to describe the quantum transport along the $`z`$-axis in a phase-space representation. In particular, we rewrite Eq. (4) in terms of $`\zeta =(z+z^{})/2`$ and $`\eta =zz^{}`$ coordinates and express the density matrix $`\rho _j`$ as $$\rho _j(𝐫,𝐫)=\underset{ms,m^{}s^{}}{}\frac{1}{2\pi }_{\mathrm{}}^+\mathrm{}𝑑ke^{ik\eta }f_{jmsm^{}s^{}}(\zeta ,k)\mathrm{\Psi }_{jms}(𝐫_{},z)\mathrm{\Psi }_{jm^{}s^{}}^{}(𝐫_{}^{},z^{}),$$ (5) with a complete set of orthonormal functions $`\mathrm{\Psi }_{jms}(𝐫_{},z)`$. According to the cylindrical symmetry of the system, these functions take the following form: $$\mathrm{\Psi }_{jms}(𝐫_{},z)=\frac{1}{\sqrt{2\pi }}\psi _{jms}(r,z)e^{im\varphi }.$$ (6) The functions $`\psi _{jms}(r,z)`$ are chosen to satisfy the equation $$\frac{\mathrm{}^2}{2m_j^{}}\left[\frac{1}{r}\frac{}{r}\left(r\frac{}{r}\right)\frac{m^2}{r^2}\right]\psi _{jms}(r,z)+V(r,z)\psi _{jms}(r,z)=_{jms}(z)\psi _{jms}(r,z),$$ (7) which describes the radial motion of an electron. Here $`_{jms}(z)`$ are the eigenvalues of Eq. (7) for a given value of the $`z`$-coordinate which appears as a parameter. It will be shown, that $`_{jms}(z)`$ plays the role of an effective potential in the channel, and that $`\mathrm{\Psi }_{jms}(𝐫_{},z)`$ is the corresponding wavefunction of the transverse motion at fixed $`z`$. Substituting the expansion (5) into Eq. (4), and using Eq. (7), we arrive at an equation for $`f_{jmsm^{}s^{}}(\zeta ,k)`$ : $`{\displaystyle \frac{f_{jmsm^{}s^{}}(\zeta ,k)}{t}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}k}{m_j^{}}}{\displaystyle \frac{}{\zeta }}f_{jmsm^{}s^{}}(\zeta ,k)+{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}W_{jmsm^{}s^{}}(\zeta ,kk^{})f_{jmsm^{}s^{}}(\zeta ,k^{})𝑑k^{}`$ (9) $`{\displaystyle \underset{s_1,s_1^{}}{}}{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}\widehat{M}_{jmsm^{}s^{}}^{s_1s_1^{}}(\zeta ,k,k^{})f_{jms_1m^{}s_1^{}}(\zeta ,k^{})𝑑k^{},`$ where $$W_{jmsm^{}s^{}}(\zeta ,kk^{})=\frac{1}{2\pi i}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\left(_{jms}(\zeta +\eta /2)_{jm^{}s^{}}(\zeta \eta /2)\right)e^{i(k^{}k)\eta }𝑑\eta ,$$ (10) $$\widehat{M}_{jmsm^{}s^{}}^{s_1s_1^{}}(\zeta ,k,k^{})=\frac{1}{2\pi }\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\left[\delta _{s^{}s_1^{}}\widehat{\mathrm{\Gamma }}_{mss_1}(\zeta +\eta /2,k^{})+\delta _{ss_1}\widehat{\mathrm{\Gamma }}_{ms^{}s_1^{}}^{}(\zeta \eta /2,k^{})\right]e^{i(k^{}k)\eta }𝑑\eta ,$$ (11) $$\widehat{\mathrm{\Gamma }}_{mss_1}(z,k^{})=\frac{\mathrm{}}{2m_j^{}i}b_{jmss_1}(z)+\frac{\mathrm{}}{2m_j^{}}c_{jmss_1}(z)\left(i\frac{}{\zeta }+2k^{}\right),$$ (12) and $`b_{jmss_1}(z)`$ $`=`$ $`{\displaystyle \psi _{jms}^{}(r,z)\frac{^2}{z^2}\psi _{jms_1}(r,z)r𝑑r},`$ (13) $`c_{jmss_1}(z)`$ $`=`$ $`{\displaystyle \psi _{jms}^{}(r,z)\frac{}{z}\psi _{jms_1}(r,z)r𝑑r}.`$ (14) Note that Eq. (9) is similar to the Liouville equation for the Wigner distribution function, which is derived to model quantum transport in tunneling diodes (see Ref. ). The first drift term in the right-hand side of Eq. (9) is derived from the kinetic-energy operator of the longitudinal motion. It is exactly the same as the corresponding term of the Boltzmann equation. The second component plays the same role as the force term does in the Boltzmann equation. The last term in the right-hand side of Eq. (9) contains the operator $`\widehat{M}_{jmsm^{}s^{}}^{s_1s_1^{}}(\zeta ,k,k^{})`$, which mixes the functions $`f_{jmsm^{}s^{}}`$ with different indexes $`s,s^{}`$. It appears because $`\psi _{jms}(r,z)`$ are not eigenfunctions of the Hamiltonian (1). The physical meaning of the operator $`\widehat{M}_{jmsm^{}s^{}}^{s_1s_1^{}}(\zeta ,k,k^{})`$ will be discussed below. In order to solve Eq. (9), we need to specify boundary conditions for the functions $`f_{jmsm^{}s^{}}(\zeta ,k)`$. For a weak current, electrons incoming from both the source and the drain electrodes, are assumed to be maintained in thermal equilibrium. Comparing Eq. (5) with the corresponding expansion of the density matrix in the equilibrium state, one obtains the following boundary conditions $`f_{jmsm^{}s^{}}(0,k)`$ $`=`$ $`\delta _{ss^{}}\delta _{mm^{}}2\left[\mathrm{exp}\left(E_{jsmk}\beta E_{FS}\beta \right)+1\right]^1,k>0,`$ (15) $`f_{jmsm^{}s^{}}(L,k)`$ $`=`$ $`\delta _{ss^{}}\delta _{mm^{}}2\left[\mathrm{exp}\left(E_{jsmk}\beta E_{FD}\beta \right)+1\right]^1,k<0,`$ (16) where the total energy is $`E_{jsmk}=\mathrm{}^2k^2/2m_j^{}+_{jsm}(0)`$ for an electron entering from the source electrode $`(k>0)`$ and $`E_{jsmk}=\mathrm{}^2k^2/2m_j^{}+_{jsm}(L)`$ for an electron entering from the drain electrode $`(k<0)`$. $`E_{FS}`$ and $`E_{FD}`$ are the Fermi energy levels in the source and in the drain, respectively. Note, that Eq. (16) meets the requirement of imposing only one boundary condition on the function $`f_{jmsm^{}s^{}}(\zeta ,k)`$ at a fixed value of $`k`$ as Eq. (9) is a first order differential equation with respect to $`\zeta `$. Generally speaking, the solution of Eq. (9) with the conditions (16) depends on the distance between the boundary position and the active device region. Let us estimate how far the boundary must be from the active device region in order to avoid this dependence. It is easy to show that the density matrix of the equilibrium state is a decaying function of $`\eta =zz^{}`$. The decay length is of the order of the coherence length $`\lambda _T=\sqrt{\frac{\mathrm{}^2}{m_j^{}k_BT}}`$ at high temperature and of the inverse Fermi wavenumber $`k_F^1=\sqrt{\mathrm{}^2/2m_j^{}E_F}`$ at low temperature. So, it is obvious, that the distance between the boundary and the channel must exceed the coherence length or the inverse Fermi wavenumber, i. e. $`L\lambda _T`$ or $`Lk_F^1`$. For example, at $`T=300`$ K the coherence length $`\lambda _T3`$ nm is much less than the source or drain lengths. The functions $`f_{jmsm^{}s^{}}(\zeta ,k)`$, which are introduced in Eq. (5), are used in calculations of the current and the electron density. The expression for the electron density follows directly from the density matrix as $`n(𝐫)=\rho (𝐫,𝐫)`$. In terms of the functions $`f_{jmsm^{}s^{}}(\zeta ,k)`$, the electron density can be written as follows: $$n(𝐫)=\frac{1}{2\pi }\underset{jmsm^{}s^{}}{}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}f_{jmsm^{}s^{}}(z,k)𝑑k\mathrm{\Psi }_{jms}(𝐫_{},z)\mathrm{\Psi }_{jm^{}s^{}}^{}(𝐫_{},z).$$ (17) It is well-known , that the current density can be expressed in terms of the density matrix $$𝐣(𝐫,t)=\underset{j}{}\frac{e\mathrm{}}{2m_ji}\left(\frac{}{𝐫}\frac{}{𝐫^{}}\right)\rho _j(𝐫,𝐫^{},t)\|_{𝐫=𝐫^{}}.$$ (18) The total current, which flows through the cross-section of the structure at a point $`z`$, can be obtained by an integration over the transverse coordinates. Substituting the expansion (5) into Eq. (18) and integrating over $`r`$ and $`\phi `$, we find $$J=e\underset{j,m,s}{}\frac{1}{2\pi }\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑k\frac{\mathrm{}k}{m_j^{}}f_{jmsms}(z,k)\frac{2e\mathrm{}}{m_j^{}}\underset{\genfrac{}{}{0pt}{}{j,m,s,s^{}}{s^{}>s}}{}c_{jmss^{}}(z)\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑k\mathrm{Im}f_{jmsms^{}}(z,k),$$ (19) where $`\mathrm{Im}f`$ is the imaginary part of $`f`$. The first term in the right-hand side of Eq. (19) is similar to the expression for a current of the classical theory . The second term, which depends on the non-diagonal functions $`f_{jmsms^{}}`$ only, takes into account the effects of intermixing between different states of the transverse motion. The last term in the right-hand side of Eq. (9) takes into consideration the variation of the wavefunctions $`\psi _{jms}(r,z)`$ along the $`z`$-axis. In the source and drain regions, the electrostatic potential is essentially constant due to the high density of electrons. In these parts of the structure, the wavefunctions of the transverse motion are very weakly dependent on $`z`$, and consequently, the operator $`\widehat{M}_{jmsm^{}s^{}}^{s_1s_1^{}}`$ has negligible effect. Inside the channel, electrons are strongly localized at the Si/SiO<sub>2</sub> interface as the positive gate voltage is applied. Earlier calculations, which we made for the case of equilibrium , have shown that in the channel the dependence of $`\psi _{jms}(r,z)`$ on $`z`$ is weak, too. Therefore, in the channel the effect of the operator $`\widehat{M}_{jmsm^{}s^{}}^{s_1s_1^{}}`$ is negligible. In the intermediate regions (the source–channel and the drain–channel), an increase of the contribution of the third term in the right-hand side of Eq. (9) is expected due to a sharp variation of $`\psi _{jms}(r,z)`$. Since $`\widehat{M}_{jmsm^{}s^{}}^{s_1s_1^{}}`$ couples functions $`f_{jmsm^{}s^{}}(\zeta ,k)`$ with different quantum numbers ($`jms`$), it can be interpreted as a collision operator, which describes transitions of electrons between different quantum states of the transverse motion. Thus, the third term in the right-hand side of Eq. (9) is significant only in the close vicinity of the p-n junctions. Therefore, this term is assumed to give a small contribution to the charge and current densities. Under the above assumption, we have treated the last term in the right-hand side of Eq. (steady state of the system in a zeroth order approximation with respect to the operator $`\widehat{M}_{jmsm^{}s^{}}^{s_1s_1^{}}`$. Neglecting the latter, one finds that, due to the boundary conditions (16), all non-diagonal functions $`f_{jmsm^{}s^{}}(\zeta ,k)`$ ($`mm^{}`$ or $`ss^{}`$) need to be zero. In the channel, the energy of the transverse motion can be approximately written in the form $$_{jms}(z)=_{js}(z)+\frac{\mathrm{}^2m^2}{2m_{js}^{}R_{js}^2},$$ (20) where $`_{js}(z)`$ is the energy associated with the radial size quantization and $`\mathrm{}^2m^2/2m_{js}^{}R_{js}^2`$ is the energy of the angular motion with averaged radius $`R_{js}^{}`$. Hence, in Eq. 10 for the diagonal functions $`f_{jmsms}(\zeta ,k)`$ the difference $`_{jms}(\zeta +\eta /2)_{jms}(\zeta \eta /2)`$ can be substituted by $`_{js}(\zeta +\eta /2)_{js}(\zeta \eta /2)`$. Furthermore, summation over $`m`$ in Eq. (9) gives $$\frac{\mathrm{}k}{m_j^{}}\frac{}{\zeta }f_{js}(\zeta ,k)\frac{1}{\mathrm{}}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}W_{js}(\zeta ,kk^{})f_{js}(\zeta ,k^{})𝑑k^{}=0$$ (21) with $$f_{js}(\zeta ,k)=\frac{1}{2\pi }\underset{m}{}f_{jmsms}(\zeta ,k).$$ (22) In Eq. (21) the following notation is used $$W_{js}(\zeta ,k)=\frac{1}{2\pi }\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\left(_{js}(\zeta +\eta /2)_{js}(\zeta \eta /2)\right)\mathrm{sin}(k\eta )𝑑\eta .$$ (23) The effective potential $`_{js}(z)`$ can be interpreted as the bottom of the subband $`(j,s)`$ in the channel. The function $`f_{js}(\zeta ,k)`$ is referred to as a partial Wigner distribution function describing electrons which are travelling through the channel in the inversion layer subband $`(j,s)`$. ## IV Electron scattering In this section we consider the electron scattering from phonons and impurities. For this purpose we introduce a Boltzmann-like single collision term , which in the present case has the following form $$Stf_{jsmk}=\underset{j^{}s^{}m^{}k^{}}{}\left(P_{jsmk,j^{}s^{}m^{}k^{}}f_{j^{}s^{}m^{}k^{}}P_{jsmk,j^{}s^{}m^{}k^{}}f_{jsmk}\right).$$ (24) As was noted above, we have neglected all transitions between quantum states with different sets of quantum numbers $`j`$ and $`s`$. In the source and drain contacts the distribution of electrons over the quantum states of the angular motion corresponds to equilibrium. Consequently, due to the cylindrical symmetry of the system, we may fairly assume that across the whole structure the electron distribution is given by $$f_{jsmk}(z)=f_{js}(z,k)w_{jsm},$$ (25) where $$w_{jsm}=\sqrt{\frac{\mathrm{}^2\beta }{2m_j^{}R_{js}^2\pi }}\mathrm{exp}\left(\frac{\beta \mathrm{}^2m^2}{2m_j^{}R_{js}^2}\right)$$ (26) is the normalized Maxwellian distribution function with respect to the angular momentum $`m`$. The integration of the both sides of Eq. (24) over the angular momentum gives the one-dimensional collision term $$Stf_{js}(z,k)=\underset{k^{}}{}\left(P_{js}(k,k^{})f_{js}(z,k^{})P_{js}(k^{},k)f_{js}(z,k)\right),$$ (27) where $$P_{js}(k,k^{})=\underset{mm^{}}{}P_{jsmk,jsm^{}k^{}}w_{jsm^{}}.$$ (28) This collision term is directly incorporated into the one-dimensional Liouville equation (21) as $$\stackrel{~}{W}_{js}(z,k,k^{})=W_{js}(z,kk^{})+P_{js}(z,k,k^{})\delta _{k,k^{}}\underset{k^{}}{}P_{js}(z,k^{},k),$$ (29) where $`\stackrel{~}{W}_{js}(z,k,k^{})`$ is the modified force term in Eq. (21). In this work we consider scattering by acceptor impurities and acoustic phonons described by a deformation potential. The scattering rates are evaluated according to Fermi’s golden rule $$P_{jsmk,jsm^{}k^{}}=\frac{2\pi }{\mathrm{}}\left\|jsm^{}k^{}\left\|\widehat{H}_{int}\right\|jsmk\right\|^2\delta \left(E_{jsm^{}k^{}}E_{jsmk}\right),$$ (30) where $`\widehat{H}_{int}`$ is the Hamiltonian of the electron-phonon or the electron-impurity interaction. Hereafter, we model the potential of an ionized acceptor as $`U(𝐫)=4\pi e^2R_s^2/\epsilon _1\delta (𝐫)`$, where $`R_s`$ determines a cross-section for scattering by an impurity. Consequently, the absolute value of the matrix element is $$\left\|jsm^{}k^{}\left\|U(𝐫𝐫_i)\right\|jsmk\right\|=4\pi e^2R_s^2/\epsilon _1\psi _{js}^2(r_i,z_i).$$ (31) Averaging this over a uniform distribution of acceptors results in the following scattering rate $$P_{jsmk,jsm^{}k^{}}^i=C_i\underset{0}{\overset{R}{}}\psi _{js}^4(r,z)\delta \left(E_{jsm^{}k^{}}E_{jsmkk}\right)r𝑑r,$$ (32) where $`C_i=N_a\left(4\pi e^2R_s^2/\epsilon _1\right)^2/\mathrm{}`$ and $`N_a`$ is the acceptor concentration. At room temperature the rate of the scattering by acoustic phonons has the same form. Indeed, for $`T=300`$ K the thermal energy $`k_BT\mathrm{}\omega _𝐪`$, therefore the acoustic deformation potential scattering is approximately elastic, and the emission and absorption rates are equal to each other. For low energies we can approximate the phonon number as $`N_𝐪k_BT/\mathrm{}\omega _𝐪1`$ and the phonon frequency $`\omega _𝐪=v_0q`$, where $`v_0`$ is the sound velocity. Assuming equipartition of energy in the acoustic modes, the scattering rate is $$P_{jsmk,jsm^{}k^{}}^{ph}=\frac{2\pi }{V}C_{ph}\underset{𝐪}{}\left\|jsm^{}k^{}\left\|e^{i𝐪𝐫}\right\|jsmk\right\|^2\delta \left(E_{jsm^{}k^{}}E_{jsmk}\right),$$ (33) where the parameter $`C_{ph}=4\mathrm{\Sigma }^2k_BT/9\pi \rho v_0^2\mathrm{}`$. Integrating over $`𝐪`$ yields the scattering rate $`P_{jsmk,jsm^{}k^{}}^{ph}`$ in the form (32) with $`C_{ph}`$ instead of $`C_i`$. The full scattering rate $`P_{jsmk,jsm^{}k^{}}=P_{jsmk,jsm^{}k^{}}^i+P_{jsmk,jsm^{}k^{}}^{ph}`$ is then inserted into Eq. (28) in order to obtain the one-dimensional scattering rate $$P_{js}(z,k,k^{})=\left(C_i+C_{ph}\right)a(z)F\left(\frac{\mathrm{}^2k^2}{2m_j^{}}\frac{\mathrm{}^2k^2}{2m_j^{}}\right),$$ (34) where $`a(z)=\sqrt{{\displaystyle \frac{2m_j^{}}{\mathrm{}^2\pi \beta }}}R_{js}(z){\displaystyle \underset{0}{\overset{R}{}}}\psi _{js}^4(r,z)r𝑑r,F(x)=e^{x/2}K_0(\|x\|/2),`$where $`K_0(x)`$ is a McDonald function . In calculations of the scattering by acoustic phonons the following values of parameters for Si are used: $`\mathrm{\Sigma }=`$9.2 eV, $`\rho =2.328310^3`$ kg/m<sup>3</sup>, $`v_0=8.4310^5`$ cm/s ## V Numerical model The system under consideration consists of regions with high (the source and drain) and low (the channel) concentrations of electrons. The corresponding electron distribution difference would produce a considerable inaccuracy if we would have attempted to directly construct a finite-difference analog of Eq. (21). It is worth mentioning that, in the quasi-classical limit, i. e. $`_{js}(\zeta +\eta /2)_{js}(\zeta \eta /2)\frac{_{js}(\zeta )}{\zeta }\eta `$, Eq. (21) leads to the Boltzmann equation with an effective potential which has the following exact solution in the equilibrium state: $$f_{js}^{eq}(\zeta ,k)=\frac{1}{\pi }\underset{m}{}\left[\mathrm{exp}\left(\frac{\mathrm{}^2k^2}{2m_j^{}}\beta +_{j,s}(\zeta )\beta +\frac{\mathrm{}^2m^2\beta }{2m_{js}^{}R_{js}^2(\zeta )}E_F\beta \right)+1\right]^1$$ (35) For numerical calculations it is useful to write down the partial Wigner distribution function as $`f_{js}(\zeta ,k)=f_{js}^{eq}(\zeta ,k)+f_{js}^d(\zeta ,k)`$. Inserting this into Eq. (21), one obtains the following equation for $`f_{js}^d(\zeta ,k)`$: $$\frac{\mathrm{}k}{m_j^{}}\frac{}{\zeta }f_{js}^d(\zeta ,k)\frac{1}{\mathrm{}}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}W_{js}(\zeta ,kk^{})f_{js}^d(\zeta ,k^{})𝑑k^{}=B_{js}(\zeta ,k),$$ (36) where $$B_{js}(\zeta ,k)=\frac{1}{2\pi }\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑k^{}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑\eta \left(_{js}(\zeta +\frac{\eta }{2})_{js}(\zeta \frac{\eta }{2})\frac{_{js}(\zeta )}{\zeta }\eta \right)\mathrm{sin}\left[(kk^{})\eta \right]f_{js}^{eq}(\zeta ,k^{}).$$ (37) The unknown function $`f_{js}^d(\zeta ,k)`$ takes values of the same order throughout the whole system, and therefore is suitable for numerical computations. In the present work, we have used the finite-difference model, which is described in Ref. . The position variable takes the set of discrete values $`\zeta _i=\mathrm{\Delta }\zeta i`$ for $`\{i=0,\mathrm{},N_\zeta \}`$. The values of $`k`$ are also restricted to the discrete set $`k_p=(2pN_k1)\mathrm{\Delta }k/2`$ for $`\{p=1,\mathrm{},N_k\}`$. On a discrete mesh, the first derivative $`\frac{f_{js}}{\zeta }(\zeta _i,k_p)`$ is approximated by the left-hand difference for $`k_p>0`$ and the right-hand difference for $`k_p<0`$. It was shown In Ref. , that such a choice of the finite-difference representation for the derivatives leads to a stable discrete model. Projecting the equation (36) onto the finite-difference basis gives a matrix equation $`𝐋𝐟=𝐛`$. In the matrix $`𝐋`$, only the diagonal blocks and one upper and one lower co-diagonal blocks are nonzero: $$𝐋=\left(\begin{array}{ccccc}A_1& E& 0& \mathrm{}& 0\\ V& A_2& E& \mathrm{}& 0\\ 0& V& A_3& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& A_{N_\zeta 1}\end{array}\right).$$ (38) Here, the $`N_k\times N_k`$ matrices $`A_i`$, $`E`$, and $`V`$ are $`\left[A_i\right]_{pp^{}}=\delta _{pp^{}}{\displaystyle \frac{2m_j^{}\mathrm{\Delta }\zeta }{\mathrm{}^2(2pN_k1)\mathrm{\Delta }k}}W_{js}(\zeta _i,k_pk_p^{}),`$ $$\left[E\right]_{pp^{}}=\delta _{pp^{}}\theta \left\{\frac{N_k+1}{2}p\right\},\left[V\right]_{pp^{}}=\delta _{pp^{}}\theta \left\{p\frac{N_k+1}{2}\right\},$$ (39) and the vectors are $$[f_i]_p=f_{js}(\zeta _i,k_p),\mathrm{and}[b_i]_p=B_{js}(\zeta _i,k_p),i=1,N_{\zeta 1},i=1,N_k.$$ (40) A recursive algorithm is used to solve the matrix equation $`𝐋𝐟=𝐛`$. Invoking downward elimination, we are dealing with $`B_i=\left(A_iVB_{i1}\right)^1E`$ and $`N_i=\left(A_iVB_{i1}\right)^1\left(b_i+VN_{i1}\right)`$ $`(i=1,\mathrm{},N_\zeta )`$ as relevant matrices and vectors. Then, upward elimination eventually yields the solution $`f_i=B_if_{i+1}+N_i`$ ($`i=N_\zeta 1,\mathrm{},1`$). If an index of a matrix or a vector is smaller than 1 or larger than $`N_\zeta 1`$, the corresponding term is supposed to vanish. In the channel, the difference between effective potentials $`_{js}(\zeta )`$ with different ($`j,s`$) is of the order of or larger than the thermal energy $`k_BT`$. Therefore, in the channel only a few lowest inversion subbands must be taken into account. In the source and drain, however, many quantum states $`(j,s)`$ of the radial motion are strongly populated by electrons. Therefore, we should account for all of them in order to calculate the charge distribution. Here, we can use the fact that, according to our approximation, the current flows only through the lowest subbands in the channel. Hence, only for these subbands the partial Wigner distribution function of electrons is non-equilibrium. In other subbands electrons are maintained in the state of equilibrium, even when a bias is applied. So, in Eq. (17) for the electron density, we can substitute functions $`f_{jmsms}(z,k)`$ of higher subbands by corresponding equilibrium functions. Formally, adding and subtracting the equilibrium functions for the lowest subbands in Eq. (17), we arrive at the following equation for the electron density $$n(𝐫)=n_{eq}(𝐫)+\frac{1}{2\pi }\underset{js}{}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑k\left[f_{js}(z,k)\left\|\psi _{js}(r,z)\right\|^2f_{js}^{eq}(z,k)\left\|\psi _{js}^{eq}(r,z)\right\|^2\right],$$ (41) where $`n_{eq}(𝐫)`$ and $`\psi _{js}^{eq}(r,z)`$ are the electron density and the wavefunction of the radial motion in the state of equilibrium, respectively. The summation on the right-hand side of Eq. (41) is performed only over the lowest subbands. Since the electrostatic potential does not penetrate into the source and drain, we suppose that the equilibrium electron density in these regions is well described by the Thomas-Fermi approximation: $$n_{eq}(𝐫)=N_C\frac{2}{\sqrt{\pi }}F_{1/2}\left(\beta \left(eV(r,z)+E_FE_C\right)\right),$$ (42) where the Fermi integral is $$F_{1/2}(x)=_0^{\mathrm{}}\frac{\sqrt{t}dt}{\mathrm{exp}(tx)+1}.$$ (43) Here $`N_C`$ is the effective density of states in the conduction band and $`E_F`$ is the Fermi level of the system in the state of equilibrium. ## VI Numerical results During the device simulation three equations are solved self-consistently: (i) the equation for the wavefunction of the radial motion (7), (ii) the equation for the partial Wigner distribution function (21) and (iii) the Poisson equation (2). The methods of numerical solution of Eqs. (7) and (2) are the same as for the equilibrium state . The numerical model for Eq. (21) was described in the previous section. In the present calculations the four lowest subbands $`(j=1,2`$ $`s=1,2)`$ are taken into account. The electron density in the channel is obtained from Eq. (17), whereas in the source and drain regions it is determined from Eq. (41). The calculations are performed for structures with a channel of radius $`R=50`$ nm and for various values of the length: $`L_{ch}=`$40, 60, 70 and 80 nm. The width of the oxide layer is taken to be 4 nm. All calculations are carried out with $`N_\zeta =100`$ and $`N_k=100`$. The partial Wigner distribution functions, which are obtained as a result of the self-consistent procedure, are then used to calculate the current according to Eq. (18). We investigate two cases: ballistic transport and quantum transport. The scattering of the electrons is taken into account. The distribution of the electrostatic potential is represented in Fig. 2 for $`V_{ds}=0.3`$ V and $`V_G=1`$ V. This picture is typical for the MOSFET structure, which is considered here. The cross-sections of the electrostatic potentials for $`r=0,30,40,45,48,50`$ nm are shown in Fig. 3. The main part of the applied gate voltage falls in the insulator (50 nm $`<r<`$ 54 nm). Along the cylinder axis in the channel, the electrostatic potential barrier for the electron increases up to about 0.4 eV. Since the potential along the cylinder axis is always high, the current mainly flows in a thin layer near the semiconductor-oxide interface. This feature provides a way of controlling $`I_{ds}`$ through the gate voltage. Varying $`V_{ds}`$ and $`V_G`$ mainly changes the shape of this narrow path, and, as a consequence, influences the form of the effective potential $`_{js}(z)`$. As follows from Figs. 2 and 3, the radius of the pillar can be taken shorter without causing barrier degradation. At the p-n-junctions (source–channel and drain–channel) the electrons meet barriers across the whole semiconductor. These barriers are found to persist even for high values of the applied source–drain voltage and prevent an electron flood from the side of the strongly doped source. The pattern of the electrostatic potentials differs mainly near the semiconductor-oxide interface, where the inversion layer is formed. In Fig. 4 the effective potential for the lowest inversion subband $`(j=1,s=1)`$ is plotted as a function of $`z`$ for different applied bias $`V_{ds}=0,\mathrm{},0.5`$ V, $`V_G=1`$ V, $`L_{ch}=60`$ nm. It is seen that the effective potential reproduces the distribution of the electrostatic potential near the semiconductor-oxide interface. In the case of ballistic transport (dashed curve), the applied drain-source voltage sharply drops near the drain-channel junction (Figs. 3 and 4). The scattering of electrons (solid curve) smoothes out the applied voltage, which is now varying linearly along the whole channel. Note, that the potential obtained by taking into account scattering is always higher than that of the ballistic case. The explanation is clear from Fig. 5, where the linear electron density is plotted for $`V_G=1`$ V and $`L_{ch}=60`$ nm. It is seen that, due to scattering, the electron density in the channel (solid curve) rises and smoothes out. Hence, the applied gate voltage is screened more effectively, and as a result, the potential exceeds that of the ballistic case. It should be noted that at equilibrium ($`V_{ds}=0`$) the linear density and the effective potential for both cases (with and without scattering) are equal to each other. This result follows from the principle of detailed balance. The current-voltage characteristics (the current density $`I=J/2\pi R`$ vs. the source–drain voltage $`V_{ds}`$) are shown in Figs. 6 and 7 for the structures with channel lengths $`L_{ch}=40,\mathrm{},80`$ nm. At a threshold voltage $`V_{ds}0.2`$ V a kink in the $`I`$$`V`$ characteristics of the device is seen. At subthreshold voltages $`V_{ds}<0.2`$ V the derivative $`dV_{ds}/dJ`$ gives the resistance of the structure. It is natural, that scattering enhances the resistance of the structure (solid curve) compared to the ballistic transport (dashed curve). Scattering is also found to smear the kink in the $`I`$$`V`$ characteristic. At a voltage $`V_{ds}>0.2`$ V a saturation regime is reached. In this part of $`IV`$ characteristics, the current through the structure increases more slowly than it does at a subthreshold voltage. The slope of the $`IV`$ curve in the saturation regime rises when the length of the channel decreases. This effect is explained by a reduction of the p-n junction barrier potential as the length of the channel becomes shorter than the p-n junction width. In Fig. 7, one can see that, when the transistor is switched off ($`V_G<0.5`$ V), the influence of the scattering on the current is weak. This fact is due to a low concentration of electrons, resulting in a low amplitude of the scattering processes. In Figs. 8a and 8b the contour plots of the partial Wigner distribution function ($`j=1,s=1`$) are given for both cases (a – without and b – with scattering). The lighter regions in these plots indicate the higher density of electrons. Far from the p-n-junction, where the effective potential varies almost linearly, the partial Wigner distribution function can be interpreted as a distribution of electrons in the phase space. When electrons travel in the inversion layer without scattering, their velocity increases monotonously along the whole channel. Therefore, in the phase-space representation the distribution of ballistic electrons looks as a narrow stream in the channel (Fig. 8a). As it is expected, scattering washes out the electron jet in the channel (see Fig. 8b). It is worth mentioning that the electron stream in the channel does not disappear. It means that in this case the electron transport through the channel combines the features of both diffusive and ballistic motion. ## VII Summary We have developed a model for the detailed investigation of quantum transport in MOSFET devices. The model employs the Wigner distribution function formalism allowing us to account for electron scattering by impurities and phonons. Numerical simulation of a cylindrical nanosize MOSFET structure was performed. $`IV`$ characteristics for different values of the channel length were obtained. It is shown that the slope of the $`IV`$ characteristic in the saturation regime rises as the channel length increases. This is due to the decrease of the p-n junction barrier potential. Finally, we have demonstrated that the inclusion of a collision term in numerical simulation is important for low source–drain voltages. The calculations have shown that the scattering leads to an increase of the electron density in the channel and smoothes out the applied voltage along the entire channel. The analysis of the electron phase-space distribution in the channel has shown that, in spite of scattering, electrons are able to flow through the channel as a narrow stream although, to a certain extent, the scattering is seen to wash out this jet. Accordingly, features of both ballistic and diffusive transport are simultaneously encountered. ###### Acknowledgements. This work has been supported by the Interuniversitaire Attractiepolen — Belgische Staat, Diensten van de Eerste Minister – Wetenschappelijke, technische en culturele Aangelegenheden; PHANTOMS Research Network; F.W.O.-V. projects Nos. G.0287.95, 9.0193.97 and W.O.G. WO.025.99N (Belgium). FIGURE CAPTIONS Fig. 1. Scheme of the cylindrical nanosize MOSFET. Fig. 2. Distribution of the electrostatic potential in the MOSFET with $`L_{ch}=60`$ nm at $`V_G=1`$ V and $`V_{ds}=0.3`$ V. Fig. 3. Cross-sections of the electrostatic potentials without scattering (dashed curves) and with scattering from acceptor impurities and from an acoustic deformation potential (solid curves). Fig. 4. Effective potential as a function of $`z`$ for various $`V_{ds}`$, $`L_{ch}=60`$ nm. Fig. 5. Linear electron density in the channel as a function of $`z`$ for various $`V_{ds}`$, $`L_{ch}=60`$ nm. Fig. 6. Current-voltage characteristics at $`V_G=1`$ V for different channel lengths. Fig. 7. Current-voltage characteristics for MOSFET with $`L_{ch}=40`$ nm. Fig. 8. Contour plots of the partial Wigner distribution function $`f_{js}(z,k)`$ for the lowest subband ($`j=1,s=1`$) at $`V_G=1`$ V, $`V_{ds}=0.3`$ V, $`L_{ch}=60`$ nm: a – without scattering, b – with scattering.
warning/0004/math0004106.html	ar5iv	text	# A continuum of totally incomparable Hereditarily indecomposable Banach spaces ## 1. Introduction All Banach spaces considered in this paper are real, infinite dimensional. By a subspace of a Banach space we shall mean an infinite dimensional, closed linear subspace. A Banach space is said to be Hereditarily Indecomposable (H.I.) if for every pair $`Y`$, $`Z`$ of subspaces of $`X`$ with $`YZ=\{\mathrm{𝟎}\}`$, the subspace $`Y+Z`$ is not closed. The famous example of Gowers and Maurey of a Banach space without unconditional basic sequence, was observed by W. Johnson to be H.I. Since the appearance of the Gowers-Maurey space the study of H.I. spaces has been one of the most important research topics in modern Banach space theory. We refer to and for a detailed survey of results. It is proved in that every Banach space not containing an isomorph of $`\mathrm{}_1`$ has a subspace which is a quotient of an H.I. space. A recent result of S. Argyros states that a separable Banach space universal for the class of reflexive H.I. spaces, is also universal for the class of separable Banach spaces. Both results indicate the large variety of H.I. spaces. The aim of this paper is towards this direction. Our main result is the following: ###### Theorem 1.1. There exists a family of cardinality equal to the continuum whose members are totally incomparable, reflexive H.I. spaces. Recall that the Banach spaces $`X`$ and $`Y`$ are totally incomparable if no subspace of $`X`$ is isomorphic to a subspace of $`Y`$. The construction of H.I. spaces is not an easy task. The crucial step was Schlumprecht’s construction of an arbitrarily distortable Banach space . Recall that the Banach space $`(X,)`$ is arbitrarily distortable if for every $`\lambda >1`$, there exists an equivalent norm $`\|\|`$ on $`X`$ so that for every subspace $`Y`$ of $`X`$ there exist non-zero vectors $`x`$, $`y`$ in $`Y`$ such that $`x=y`$, yet $`\|x\|/\|y\|>\lambda `$. Schlumprecht’s space had an immense impact in the development of the theory because of its connection to the Gowers-Maurey construction, as well as to the solution of the distortion problem for $`\mathrm{}_p`$, $`1<p<\mathrm{}`$, . The first example of an arbitrarily distortable, asymptotic $`\mathrm{}_1`$ space was given in . They showed that there exist infinite subsets $`M=(m_i)`$, $`N=(n_i)`$ of $``$ so that the mixed Tsirelson space $`T(\frac{1}{m_i},S_{n_i})_{i=1}^{\mathrm{}}`$, is arbitrarily distortable. In the same paper this example was conditionalized to yield an asymptotic $`\mathrm{}_1`$ H.I. space. The proof of Theorem 1.1 is based on ideas from . However, our argument is considerably simpler. We shall next describe how this paper is organized. In Section 3 we introduce, for a given scalar $`d>1`$, the infinite subsets $`N`$ and $`P`$ of $``$ and the null scalar sequence $`𝐚`$, the $`(d,N,P,𝐚)`$ distortion property, Definition 3.1, of an asymptotic $`\mathrm{}_1`$ Banach space. This property will enable us to give a criterion, Theorem 3.2, for an an asymptotic $`\mathrm{}_1`$ Banach space to be arbitrarily distortable. We also show how to obtain totally incomparable arbitrarily distortable spaces. We apply Theorem 3.2 in Section 4 in order to give an alternative proof of the fact that certain mixed Tsirelson spaces are arbitrarily distortable , , . These spaces can be described as the completion of $`c_{00}`$, the space of all ultimately vanishing real sequences, under the norm given by $`x=sup\{_{i=1}^{\mathrm{}}\mu (\{i\})x(i):\mu \}`$, where $``$ is a suitable symmetric subset of the finitely supported signed measures on $``$ containing the point mass measures and closed under interval restrictions. The main difficulty in the study of mixed Tsirelson spaces is that the norming set $``$ is defined by means of an inductive procedure. We are able to by pass this difficulty by describing $``$ analytically and proving a decomposition result for its members, Lemma 4.3, which greatly simplifies the argument for the distortion of $`T(\frac{1}{m_i},S_{n_i})_{i=1}^{\mathrm{}}`$. In Section 5, we choose a subset $`𝒩`$ of $``$ which is maximal with respect to a Maurey-Rosenthal type of condition and show in Theorem 3.5 that the completion of $`c_{00}`$ under the norm induced by $`𝒩`$ is an H.I. space satisfying a $`(d,N,P,𝐚)`$ distortion property. Various choices of $`𝒩`$ give rise to totally incomparable H.I. spaces. In order to prove that a space $`X`$ is H.I., we employ Theorem 3.6 which loosely speaking asserts that if for every $`ϵ>0`$ there exist integers $`k<n`$ such that every block subspace $`Y`$ of $`X`$ contains a sufficiently large (in the Schreier sense) block basis $`z_1<\mathrm{}<z_p`$ with the property that $`_{i=1}^pa_iz_iϵ_{i=1}^pa_ie_i_n`$, whenever $`(a_i)_{i=1}^p^+`$, while $`_{i=1}^pa_iz_i_{i=1}^pa_ie_i_{Ck}`$, for every sequence $`(a_i)_{i=1}^p`$ in $``$, then $`X`$ contains no infinite unconditional sequence. In the above, $`(e_i)`$ is the natural unit vector basis of $`c_{00}`$ and $`_n`$, $`_{Ck}`$ denote the $`n`$-th Schreier and $`k`$-th conditional Schreier norms respectively. The precise statements for the results mentioned above are given in Section 3. The proof of Theorem 1.1, presented in Section 3, follows from Theorem 3.5 and Proposition 3.3 combined with two fundamental results of descriptive set theory, the infinite Ramsey theorem , and a theorem of Kuratowski . ## 2. Preliminaries We shall make use of standard Banach space facts and terminology as may be found in . If $`D`$ is any set, we let $`[D]`$ (resp. $`[D]^<\mathrm{}`$) denote the set of its infinite (resp. finite) subsets. Given $`M[]`$, the notation $`M=(m_i)`$ indicates that $`M=\{m_1<m_2<\mathrm{}\}`$. Let $`E`$ and $`F`$ be finite subsets of $``$. We write $`E<F`$ if $`\mathrm{max}E<\mathrm{min}F`$. Suppose now that $`X`$ is a Banach space with a Schauder basis $`(e_n)`$. A sequence $`(u_n)`$ in $`X`$ is a block basis of $`(e_n)`$ if there exist successive subsets $`F_1<F_2<\mathrm{}`$ of $``$ and a scalar sequence $`(a_n)`$ so that $`u_n=_{iF_n}a_ie_i`$, for every $`n`$. We adopt the notation $`u_1<u_2<\mathrm{}`$ to indicate that $`(u_n)`$ is a block basis of $`(e_n)`$. We let $`\mathrm{supp}u_n`$ denote the set $`\{iF_n:a_i0\}`$. The range $`r(u_n)`$ of $`u_n`$, is the smallest integer interval containing $`\mathrm{supp}u_n`$. The subspace of $`X`$ generated by a block basis of $`(e_n)`$ is called a block subspace. We next review two important hierarchies. The Schreier hierarchy $`\{S_\xi \}_{\xi <\omega _1}`$, and the repeated averages hierarchy, $`(\xi _n^M)_{n=1}^{\mathrm{}}`$, $`\xi <\omega _1`$, $`M[]`$, . Since we shall only be using the families $`\{S_\xi \}_{\xi <\omega }`$, and $`(\xi _n^M)_{n=1}^{\mathrm{}}`$, $`\xi <\omega `$, $`M[]`$, we confine the definitions to the finite ordinal case. The Schreier families. We let $`S_0=\{\{n\}:n\}\}\{\mathrm{}\}`$. Suppose $`S_\xi `$ has been defined, $`\xi <\omega `$. We set $$S_{\xi +1}=\{_{i=1}^nF_i:n,n\mathrm{min}F_1,F_1<\mathrm{}<F_n,F_iS_\xi (in)\}\{\mathrm{}\}.$$ An important property shared by the Schreier families is that they are hereditary: If $`FS_\xi `$ and $`GF`$, then $`GS_\xi `$. Another important property is that they are spreading: If $`\{p_1,\mathrm{},p_k\}S_\xi `$, $`p_1<\mathrm{}<p_k`$, and $`q_1<\mathrm{}<q_k`$ are so that $`p_iq_i`$ for all $`ik`$, then $`\{q_1,\mathrm{},q_k\}S_\xi `$. It is not hard to check that if $`F_1<\mathrm{}<F_n`$ are members of $`S_\alpha `$ such that $`\{\mathrm{min}F_i:in\}`$ belongs to $`S_\beta `$, then $`_{i=1}^nF_i`$ belongs to $`S_{\alpha +\beta }`$. The repeated averages hierarchy. We first let $`(e_n)`$ denote the unit vector basis of $`c_{00}`$. Given $`\xi <\omega `$ and $`M[]`$, we define by induction, a sequence $`(\xi _n^M)_{n=1}^{\mathrm{}}`$ of finitely supported probability measures on $``$ whose supports are successive subsets of $`M`$. If $`\xi =0`$, then $`\xi _n^M=e_{m_n}`$, for all $`n`$, where $`M=(m_n)`$. Assume that $`(\xi _n^M)_{n=1}^{\mathrm{}}`$ has been defined for all $`M[]`$. Set $$[\xi +1]_1^M=\frac{1}{m_1}\underset{i=1}{\overset{m_1}{}}\xi _i^M$$ where $`m_1=\mathrm{min}M`$. Suppose that $`[\xi +1]_1^M<\mathrm{}<[\xi +1]_n^M`$ have been defined. Let $$M_n=\{mM:m>\mathrm{max}\mathrm{supp}[\xi +1]_n^M\}\text{ and }k_n=\mathrm{min}M_n.$$ Set $$[\xi +1]_{n+1}^M=\frac{1}{k_n}\underset{i=1}{\overset{k_n}{}}\xi _i^{M_n}.$$ It follows that $`\mathrm{supp}\xi _n^M`$ belongs to $`S_\xi `$, and moreover it is a maximal (under inclusion) member of $`S_\xi `$. It can be easily shown, by induction, that if $`i`$ and $`j`$ belong to $`\mathrm{supp}\xi _n^M`$ and $`i<j`$, then $`\xi _n^M(\{i\})\xi _n^M(\{j\})`$. For a probability measure $`\mu `$ in $``$ and $`\xi <\omega `$, we set $`\mu _\xi =sup\{\mu (F):FS_\xi \}`$. It is proven in , that $`\xi _1^M_{\xi 1}\frac{\xi }{\mathrm{min}M}`$, for every $`\xi 1`$ and $`M[]`$. It follows that for every $`P[]`$, every $`\xi 1`$ and every $`ϵ>0`$, there exists $`M[P]`$ such that $`\xi _1^M_{\xi 1}<ϵ`$. This property of the repeated averages will be very useful in the sequel. For a detailed study of these hierarchies we refer to , , , , and . We continue by introducing some more terminology. A finite collection $``$ of finite subsets of $``$ is said to be $`rS_\xi `$-admissible, $`\xi <\omega `$, $`r`$, if there exists an enumeration $`\{I_k:kn\}`$ of $``$ such that $`I_1<\mathrm{}<I_n`$ and the set $`\{\mathrm{min}I_k:kn\}`$ is the union of $`r`$ members of $`S_\xi `$. In case $`\{\mathrm{min}I_k:kn\}`$ is a maximal (under inclusion) member of $`S_\xi `$, $``$ is called maximally $`S_\xi `$-admissible. A finite block basis $`u_1<\mathrm{}<u_n`$ in a Banach space with a basis is $`rS_\xi `$ (resp. maximally $`S_\xi `$)-admissible, if $`\{\mathrm{supp}u_i:in\}`$ is. In what follows, $`X`$ is a Banach space with a basis $`(e_n)`$. The support of every block basis of $`(e_n)`$ will always be taken with respect to $`(e_n)`$. ###### Definition 2.1. Let $`(u_n)`$ be a normalized block basis of $`(e_n)`$, $`ϵ>0`$ and $`1\xi <\omega `$. Set $`p_n=\mathrm{min}\mathrm{supp}u_n`$, $`n`$, and $`P=(p_n)`$. 1. A generic $`(ϵ,\xi )`$ average of $`(u_n)`$ is any vector of the form $`_{n=1}^{\mathrm{}}\xi _1^R(p_n)u_n`$, where $`R[P]`$ and $`\xi _1^R_{\xi 1}<ϵ`$. 2. An $`(ϵ,\xi )`$ average of $`(u_n)`$ is any generic $`(ϵ,\xi )`$ average of a normalized block basis of $`(u_n)`$. 3. A normalized $`(ϵ,\xi )`$ average of $`(u_n)`$ is any vector $`u`$ of the form $`u=\frac{v}{v}`$, where $`v`$ is an $`(ϵ,\xi )`$ average of $`(u_n)`$. In case $`v\frac{1}{2}`$, $`u`$ is a smoothly normalized $`(ϵ,\xi )`$ average of $`(u_n)`$. ###### Notation . Let $`E^{}`$ be a finite collection of successive intervals of $``$ and let $`u`$ be a finite linear combination of $`(e_n)`$. 1. We let $`I(u,E^{})`$ denote the number of elements of $`E^{}`$ which are intersected by $`\mathrm{supp}u`$. 2. Let $`D`$ be a finite block basis of $`(e_n)`$ such that the support of every member of $`D`$ intersects at least one member of $`E^{}`$. We set $`D(E^{},1)=\{uD:I(u,E^{})=1\}`$ and $`D(E^{},2)=\{uD:I(u,E^{})2\}`$. Before closing this section, we recall the definitions of the Schreier space, $`X^\xi `$, and conditional Schreier space, $`CX^\xi `$, $`\xi <\omega `$. $`X^\xi `$ is the completion of $`c_{00}`$ under the norm $`x_\xi =sup\{_{iF}\|x(i)\|:FS_\xi \}`$. $`X^0`$ is isometric to $`c_0`$. $`X^1`$ was introduced by Schreier in order to provide an example of a weakly null sequence without Cesaro summable subsequence. The generalized family of Schreier spaces $`\{X^\xi \}_{\xi <\omega _1}`$ was studied in , where it is shown that the natural Schauder basis $`(e_n)`$ of $`X^\xi `$ is $`1`$-unconditional and shrinking. For a detailed study of the spaces $`\{X^\xi \}_{\xi <\omega }`$ we refer to . The conditional Schreier spaces $`\{CX^\xi \}_{\xi <\omega }`$, were constructed by H. Rosenthal (unpublished). $`CX^\xi `$ is the completion of $`c_{00}`$ under the norm $$x_{C\xi }=sup\{\underset{k=1}{\overset{n}{}}\left\|\underset{iJ_k}{}x(i)\right\|:n,(J_k)_{k=1}^n\text{ are }S_\xi \text{ admissible intervals }\}.$$ The natural basis $`(e_n)`$ of $`CX^\xi `$ is of course, a conditional basis. When $`\xi =0`$, $`(e_n)`$ is equivalent to the summing basis of $`c_0`$. We also mention the following useful fact: Suppose $`(a_i)_{i=1}^n`$ is a non-increasing finite sequence of non-negative scalars. Then $`_{i=1}^n(1)^ia_ie_{t_i}_{C\xi }_{i=1}^na_ie_{t_i}_\xi `$, for every increasing sequence of integers $`(t_i)_{i=1}^n`$. ## 3. Main results We start this section by recalling that a normalized sequence $`(x_n)`$ in a Banach space is an $`ϵ`$-$`\mathrm{}_1^\xi `$ spreading model, $`ϵ>0`$, if $`_{iF}a_ix_iϵ_{iF}\|a_i\|`$, for every $`FS_\xi `$ and all choices of scalars $`(a_i)_{iF}`$. A Banach space $`X`$ with a basis $`(e_n)`$ is asymptotic $`ϵ`$-$`\mathrm{}_1^\xi `$, $`1\xi <\omega `$, if every normalized block basis of $`(e_n)`$ is an $`ϵ`$-$`\mathrm{}_1^\xi `$ spreading model. $`X`$ is asymptotic $`\mathrm{}_1`$, if it is asymptotic $`ϵ`$-$`\mathrm{}_1^1`$, for some $`ϵ>0`$ . For an asymptotic $`ϵ`$-$`\mathrm{}_1^\xi `$ space $`X`$ with a basis $`(e_n)`$ and $`\delta >0`$, we define $`\tau (X,\delta )=sup`$ $`\{\zeta <\omega :\text{ every normalized block basis of }(e_n)`$ $`\text{has a subsequence which is a }\delta \mathrm{}_1^\zeta \text{ spreading model }\}.`$ Evidently, $`\tau (X,ϵ)\xi `$. The modulus $`\tau (X,\delta )`$ is implicitly defined in and . Of course $`\tau (X,\delta )`$ depends on the choice of the basis $`(e_n)`$, but it will be clear from the context which basis is used. In case $`U`$ is a block subspace of $`X`$, $`\tau (U,\delta )`$ will be calculated with respect to the block basis that generates $`U`$. ###### Definition 3.1. Let $`X`$ be a Banach space with a basis $`(e_i)`$. Let $`N=(n_i)`$ and $`P=(p_i)`$ be infinite subsets of $``$ such that $`n_{i1}p_i<\frac{n_i}{2}`$, for every $`i`$. Let $`𝐚=(\delta _i)`$ be a decreasing null sequence of scalars, and let $`d>1`$. $`X`$ is said to satisfy the $`(d,N,P,𝐚)`$ distortion property if for every $`j`$, $`X`$ is an asymptotic $`\delta _j`$-$`\mathrm{}_1^{n_j}`$ space such that $`\tau (U,d\delta _j)<p_j`$, for every block subspace $`U`$ of $`X`$. ###### Theorem 3.2. Let $`(X,)`$ be a Banach space with a normalized, shrinking, bimonotone basis $`(e_i)`$. Suppose that there exist $`N`$, $`P`$ in $`[]`$, a scalar sequence $`𝐚=(\delta _i)`$ and $`d>1`$ so that $`X`$ satisfies the $`(d,N,P,𝐚)`$ distortion property. Then $`X`$ is arbitrarily distortable. ###### Proof. In the sequel the admissibility of every block basis of $`(e_i)`$ will always be considered with respect to $`(e_i)`$. Given $`j`$, we set $$𝒜_j=\{\delta _j\underset{i=1}{\overset{k}{}}x_i^{}:(x_i^{})_{i=1}^kB_X^{}\text{ is }S_{n_j}\text{admissible}\}.$$ In the above, the admissibility of $`(x_i^{})`$ is measured with respect to $`(e_i^{})`$, the sequence of functionals biorthogonal to $`(e_i)`$. Because $`\tau (X,\delta _j)n_j`$, we have that $`𝒜_jB_X^{}`$. Indeed, suppose that $`\delta _j_{i=1}^kx_i^{}𝒜_j`$ and let $`xX`$, $`x1`$. Put $`x_i=x\|\mathrm{ran}(x_i^{})`$, $`ik`$. Since $`(e_i)`$ is bimonotone, $`_{i=1}^kx_i1`$. Furthermore, $`(x_i)_{i=1}^k`$ is $`S_{n_j}`$ admissible. Hence, $`\delta _j_{i=1}^kx_i1`$ and the assertion follows. We define an equivalent norm $`_j`$ on $`X`$ in the following manner: $$x_j=\delta _jx+sup\{x^{}(x):x^{}𝒜_j\}.$$ Let $`(u_i)`$ be a normalized block basis of $`(e_i)`$, and let $`j_0`$. Let $`U`$ be the block subspace of $`X`$ generated by $`(u_i)`$. Since $`\tau (U,d\delta _{j_0})<p_{j_0}`$, there exists a normalized block basis $`(v_i)`$ of $`(e_i)`$ in $`U`$ having no subsequence which is a $`d\delta _{j_0}`$-$`\mathrm{}_1^{p_{j_0}}`$ spreading model. It follows, by the main result of combined with Corollary 3.6 of , that there exists a subsequence $`(v_i)_{iM}`$ of $`(v_i)`$ such that for every $`x^{}B_X^{}`$, the block basis $`V_x^{}=\{v_i:iM,\|x^{}(v_i)\|8d\delta _{j_0}\}`$, is $`S_{p_{j_0}}`$ admissible. We next choose $`v_0`$, a generic $`(\delta _{j_0},n_{j_0})`$ average of $`(v_i)_{iM}`$. It is easily seen that for some $`x_0^{}𝒜_{j_0}`$ we have that $`x_0^{}(v_0)\delta _{j_0}`$. Therefore, $`v_0_{j_0}\delta _{j_0}`$. On the other hand, $`V_x^{}`$ is $`S_{p_{j_0}}`$ admissible, for every $`x^{}B_X^{}`$ and $`p_{j_0}<n_{j_0}`$. It follows that $`v_0(8d+1)\delta _{j_0}`$. We let $`v=\frac{v_0}{v_0}`$ and observe that $`v_{j_0}\frac{1}{8d+1}`$. Let now $`j>j_0`$. Arguing similarly, we can find a normalized block basis $`(w_i)`$ of $`(u_i)`$ and a generic $`(\delta _{j}^{}{}_{}{}^{2},n_j)`$ average $`w_0`$ of $`(w_i)`$ such that $`v<w_0`$ and $`\delta _jw_0(8d+1)\delta _j`$. We let $`w=\frac{w_0}{w_0}`$. We are going to show that $`w_{j_0}(8d+5)\delta _{j_0}`$. Suppose that $`\delta _{j_0}_{i=1}^kx_i^{}𝒜_{j_0}`$, and let $`E^{}`$ denote the collection of the ranges of the $`x_i^{}`$’s. Let $`D=\{w_r:\|_{i=1}^kx_i^{}(w_r)\|8d\delta _j\}`$. Observe that by the choice of $`(w_i)`$ we have that $`D(E^{},1)`$ is $`2S_{n_{j_0}+p_j}`$ admissible. On the other hand $`D(E^{},2)`$ is $`2S_{n_{j_0}}`$ admissible and thus $`D`$ is $`4S_{2p_j}`$ admissible. Because $`2p_j<n_j`$, we obtain the estimate $`_{i=1}^kx_i^{}(w_0)(8d+4)\delta _j`$. Hence, $`w_{j_0}(8d+5)\delta _{j_0}`$, as claimed. Finally, $`\frac{v_{j_0}}{w_{j_0}}\frac{1}{(8d+1)(8d+5)\delta _{j_0}}`$. The proof is now complete since $`j_0`$ was arbitrary. ∎ ###### Proposition 3.3. Let $`X_r`$ have a shrinking basis $`(e_k^r)_{k=1}^{\mathrm{}}`$, $`r=1,2`$. Assume that $`X_r`$ satisfies the $`(d_r,N_r,P_r,𝐚)`$ distortion property, $`r=1,2`$, and that $`𝐚=(\delta _i)`$ satisfies $`lim_i\frac{\delta _{i+1}}{\delta _i}=0`$. Suppose that for every $`i_0`$ there exist $`i>j>i_0`$ such that $`n_i^1=n_j^2`$, where $`N_r=(n_k^r)_{k=1}^{\mathrm{}}`$, $`r=1,2`$. Then $`X_1`$ and $`X_2`$ are totally incomparable. ###### Proof. Suppose the assertion is false. A standard perturbation argument yields a normalized block basis $`(u_k)`$ of $`(e_k^1)`$ equivalent to a block basis $`(w_k)`$ of $`(e_k^2)`$. Let $`T`$ be an isomorphism from $`[(u_k)]`$ onto $`[(w_k)]`$ such that $`T(u_k)=w_k`$, for all $`k`$. We can choose $`i_0`$ such that $`\frac{\delta _{i+1}}{\delta _i}<\frac{1}{d_1TT^1}`$, for every $`ii_0`$. Our assumptions allow us to choose $`i>j>i_0`$ such that $`n_i^1=n_j^2`$. Let $`(v_k)`$ be a normalized block basis of $`(u_k)`$ having no subsequence which is a $`d_1\delta _i`$-$`\mathrm{}_1^{n_i^1}`$ spreading model. But since $`(T(v_k))`$ is a block basis of $`(w_k)`$, it follows that for every $`FS_{n_j^2}`$ and all choices of scalars $`(a_k)_{kF}`$ $$\underset{kF}{}a_kT(v_k)\frac{\delta _j}{T^1}\underset{kF}{}\|a_k\|.$$ Hence, $`_{kF}a_kv_k\frac{\delta _j}{TT^1}_{kF}\|a_k\|`$, for every $`FS_{n_j^2}`$ and all choices of scalars $`(a_k)_{kF}`$. However, $`\frac{\delta _i}{\delta _j}\frac{\delta _{j+1}}{\delta _j}`$, and therefore $`\frac{\delta _j}{TT^1}>d_1\delta _i`$. Thus, $`(v_k)`$ is a $`d_1\delta _i`$-$`\mathrm{}_1^{n_i^1}`$ spreading model contrary to our assumptions. ∎ ###### Definition 3.4. Let $`M=(m_i)[]`$ such that $`m_1>6`$ and $`m_i^2<m_{i+1}`$, for all $`i`$. Choose $`L[]`$, $`L=(l_i)`$ such that $`l_1>4`$ and $`2^{l_i}>m_i`$, for all $`i`$. The infinite subset $`N=(n_i)`$ of $``$ is said to be $`M`$-good, if $`l_j(f_j^N+1)<n_j`$, for all $`j`$. In the above, $`(f_j^N)`$ is the sequence given by $`f_1^N=1`$ while for $`j2`$, $$f_j^N=\mathrm{max}\left\{\underset{i<j}{}\rho _in_i:\rho _i\{0\}(i<j),\underset{i<j}{}m_i^{\rho _i}<m_j^3\right\}.$$ Note that $`f_j^N`$ is well defined because $`m_1>1`$. It is easy to see that for every $`P[]`$ there exists $`N[P]`$ which is $`M`$-good. The main result of Section 5 is the following ###### Theorem 3.5. Suppose $`N=(n_i)`$ is $`M`$-good. Set $`N^{(2)}=(n_{2i})`$, $`F^{(2)}=(f_{2i}^N+2)`$ and $`𝐚=(\frac{1}{m_{2i}})`$. Then there exists a reflexive H.I. space $`X(N)`$ satisfying the $`(6,N^{(2)},F^{(2)},𝐚)`$ distortion property. The proof is given in Section 5. We now pass to the ###### Proof of Theorem 1.1. We first choose $`N_0[]`$ such that every $`N[N_0]`$ is $`M`$-good. To see that such a $`N_0`$ exists, set $$𝒟=\{N[]:N\text{ is }M\text{good }\}.$$ We can easily verify that $`𝒟`$ is closed in the topology of pointwise convergence in $`[]`$, and therefore it is a Ramsey set. Because $`𝒟[R]\mathrm{}`$, for every $`R[]`$, the infinite Ramsey theorem yields $`N_0[]`$ such that $`[N_0]𝒟`$, as claimed. It is a well known fact that $`[N_0]`$ endowed with the topology of pointwise convergence is a perfect Polish space. We let $`[N_0]^2=[N_0]\times [N_0]`$ and set $$G=\{(N,R)[N_0]^2,N=(n_i),R=(r_i)\|i_0,i>j>i_0:n_{2i}=r_{2j}\}.$$ A straightforward application of the Baire category theorem yields that $`G`$ is a dense $`G_\delta `$ subset of $`[N_0]\times [N_0]`$. By a result of Kuratowski and Mycielski (cf. , p. 129, Theorem 19.1, or Proposition 3.6 of ), there exists $`C[N_0]`$ homeomorphic to the Cantor set such that $`(N_1,N_2)G`$, whenever $`N_1`$, $`N_2`$ are distinct elements of $`C`$. We can now apply Theorem 3.5 to obtain a family $`\{X(N):NC\}`$ of reflexive H.I. spaces such that for every $`NC`$, $`X(N)`$ satisfies the $`(6,N^{(2)},F^{(2)},𝐚)`$ distortion property, where $`N^{(2)}`$, $`F^{(2)}`$ and $`𝐚`$ are as in the statement of Theorem 3.5. Since $`(N_1,N_2)G`$ whenever $`N_1`$ and $`N_2`$ are distinct elements of $`C`$, Proposition 3.3 implies that $`X(N_1)`$ and $`X(N_2)`$ are totally incomparable. The proof of the theorem is now complete. ∎ To construct H.I. spaces we shall make use of the following ###### Theorem 3.6. Let $`X`$ be a Banach space with a basis $`(x_i)`$. Let $`(n_j)`$, $`(k_j)`$ be increasing sequences of positive integers such that $`k_j<n_j`$, for all $`j`$, and let $`(\delta _j)`$ be a null sequence of positive scalars. Assume that for every block subspace $`Y`$ of $`X`$ and every $`j`$ there exists a block basis $`z_1<\mathrm{}<z_p`$ of $`(x_i)`$ in $`Y`$ such that letting $`t_i=\mathrm{min}\mathrm{supp}z_i`$, $`ip`$, the following are satisfied: 1. $`\{t_i:ip\}`$ is a maximal $`S_{n_j}`$ set and $`_{i=1}^pa_iz_ic_1\delta _j_{i=1}^pa_ie_{t_i}_{n_j}`$, for every sequence $`(a_i)_{i=1}^p`$ in $`^+`$. 2. $`_{i=1}^pa_iz_ic_2_{i=1}^pa_ie_{t_i}_{Ck_j}+c_3\delta _{j}^{}{}_{}{}^{2}`$, for every sequence $`(a_i)_{i=1}^p`$ in $``$ with $`_{i=1}^p\|a_i\|1`$, where $`c_1`$, $`c_2`$ and $`c_3`$ are absolute positive constants. Then $`X`$ has no infinite unconditional sequence. If moreover, given $`Y`$, $`Z`$ block subspaces of $`X`$ and $`j`$, such a block basis $`(z_i)_{i=1}^p`$ can be found with the additional property that $`z_iY`$, if $`i`$ is odd, while $`z_iZ`$, if $`i`$ is even, then $`X`$ is H.I. ###### Proof. Let $`(u_i)`$ be an infinite block basis of $`(x_i)`$, and let $`j`$. Set $`P=\{p_i:i\}`$, where $`p_i=\mathrm{min}\mathrm{supp}u_i`$. We can find $`R[P]`$ such that $`[n_j]_1^L_{k_j}<\delta _{j}^{}{}_{}{}^{2}`$, for every $`L[R]`$. Let $`Y=[u_i:p_iR]`$. Choose $`z_1<\mathrm{}<z_p`$ in $`Y`$, according to the hypothesis. There exists $`L[R]`$ such that $`\{t_i:ip\}=\mathrm{supp}[n_j]_1^L`$. Put $`a_i=[n_j]_1^L(t_i)`$, $`ip`$, and note that $`(a_i)_{i=1}^p`$ is non-increasing. We now have that $$\underset{i=1}{\overset{p}{}}a_iz_ic_1\delta _j\underset{i=1}{\overset{p}{}}a_ie_{t_i}_{n_j}=c_1\delta _j.$$ On the other hand, $$\underset{i=1}{\overset{p}{}}(1)^ia_iz_ic_2\underset{i=1}{\overset{p}{}}a_ie_{t_i}_{k_j}+c_3\delta _{j}^{}{}_{}{}^{2},$$ as $`(a_i)_{i=1}^p`$ is non-increasing. Hence, $$\underset{i=1}{\overset{p}{}}(1)^ia_iz_i(c_2+c_3)\delta _{j}^{}{}_{}{}^{2}\frac{c_2+c_3}{c_1}\delta _j\underset{i=1}{\overset{p}{}}a_iz_i.$$ Since $`j`$ was arbitrary, $`(u_i)`$ is not unconditional. The moreover statement is immediate. ∎ ## 4. Mixed Tsirelson spaces Recall that if $``$ is a set of finitely supported signed measures on $``$ which satisfies the following: 1. $`e_n^{}`$, for all $`n`$, where $`e_n^{}`$ denotes the point mass measure at $`n`$. 2. $``$ is symmetric i.e., if $`\mu `$ then $`\mu `$, 3. $``$ is pointwise bounded, that is $`\mu (\{n\})1`$, for every $`\mu `$, 4. $``$ is closed under restriction to initial segments i.e., if $`\mu `$, then $`\mu \|\{1,\mathrm{},n\}`$, then one can define a norm $`_{}`$ on $`c_{00}`$ in the following manner: $$\underset{i=1}{\overset{\mathrm{}}{}}a_ie_i_{}=sup\{\underset{i=1}{\overset{\mathrm{}}{}}a_i\mu (\{i\}):\mu \},$$ for every finitely supported scalar sequence $`(a_i)`$. Of course, $`(e_i)`$ is the natural basis of $`c_{00}`$. Letting $`X_{}`$ denote the completion of $`(c_{00},_{})`$, we see that $`(e_n)`$ is a normalized, monotone basis for $`X_{}`$. In case $`\mu \|J`$, for every $`\mu `$ and $`J`$, then $`(e_n)`$ is $`1`$-unconditional and bimonotone. The main result of this section is ###### Theorem 4.1. Suppose $`N`$ is $`M`$-good. There exists $``$, a set of finitely supported signed measures on $``$ satisfying conditions 1-4, above, and such that the following properties are fulfilled: 1. $`(e_n)`$ is an $`1`$-unconditional, shrinking, bimonotone basis for $`X_{}`$. 2. $`X_{}`$ satisfies the $`(6,N,P,𝐚)`$ distortion property, where $`P=(f_i^N+2)`$ and $`𝐚=(\frac{1}{m_i})`$. We first give the construction of $``$ and prove a number of lemmas necessary for the proof of Theorem 4.1. Construction of $``$. Given $`M=(m_i)`$, $`N=(n_i)`$, with $`N`$ being $`M`$-good, we construct $``$, a set of signed measures on $``$ in the following manner: Let $`𝒟=\{(t_1,\mathrm{},t_{3n}):n,`$ $`t_{3i2}M(i<n),t_{3n2}=0,`$ $`t_{3i1}[]^<\mathrm{}\{\mathrm{}\},t_{3i}\{1,1\}(in)\}.`$ Given $`F𝒟^<\mathrm{}`$, $`F\mathrm{}`$, we let $`𝒯_F`$ denote the set of all tuples of length divisible by $`3`$ which are initial segments of elements of $`F`$. We can partially order the elements of $`𝒯_F`$ by initial segment inclusion and thus $`𝒯_F`$ becomes a finite tree with terminal nodes precisely the members of $`F`$. Given $`\alpha 𝒯_F`$ then $`mM`$ is an $`M`$-entry of $`\alpha `$, if $`m\alpha `$. We shall denote the last three entries of $`\alpha `$ by $`m_\alpha `$, $`I_\alpha `$ and $`ϵ_\alpha `$ respectively. A rooted tree $`𝒯=𝒯_F`$ (a tree is rooted if it has a unique root), is said to be appropriate provided the following properties hold: 1. If $`\alpha 𝒯`$ is terminal, then $`I_\alpha =\{p_\alpha \}`$, for some $`p_\alpha `$. 2. If $`\alpha 𝒯`$ is non-terminal and $`m_\alpha =m_j`$, for some $`j`$, then $`(I_\beta )_{\beta D_\alpha }`$ is $`S_{n_j}`$-admissible and $`I_\alpha =_{\beta D_\alpha }I_\beta `$. Here $`D_\alpha `$ stands for the set of the immediate successors of $`\alpha `$ in $`𝒯`$. We set $$𝒢=\{𝒯:𝒯\text{ is an appropriate tree }\}.$$ We make the convention that the empty tree belongs to $`𝒢`$. ###### Notation . Let $`𝒯𝒢`$ and $`\alpha 𝒯`$. 1. $`\alpha ^{}`$ stands for the predecessor of $`\alpha `$ in $`𝒯`$. In case $`\alpha `$ is the root of $`𝒯`$ we put $`\alpha ^{}=\mathrm{}`$. 2. $`\|\alpha \|`$ is the length of $`\alpha `$. Thus, $`\|\alpha \|=3n`$ if $`\alpha =(t_1,\mathrm{},t_{3n})`$. We now define $`o(𝒯)=\mathrm{max}\{\|\beta \|:\beta 𝒯\}`$, the height of the tree $`𝒯`$. 3. $`m(\alpha )=_{m_i\alpha ^{}}m_i`$. We set $`m(\alpha )=1`$ if $`\|\alpha \|=3`$. 4. $`n(\alpha )=_{m_i\alpha ^{}}n_i`$. We set $`n(\alpha )=0`$, if $`\|\alpha \|=3`$. Given $`𝒯𝒢`$, set $$\mu _𝒯=\underset{\alpha max𝒯}{}m(\alpha )^1ϵ(\alpha )ϵ_\alpha e_{p_\alpha }^{},$$ where $`max𝒯`$ is the set of terminal nodes of $`𝒯`$ and $`I_\alpha =\{p_\alpha \}`$ for $`\alpha max𝒯`$. We have also set $`ϵ(\alpha )=_{\beta <\alpha }ϵ_\alpha `$ for $`\alpha 𝒯`$. We make the convention $`ϵ(\alpha )=1`$, if $`\|\alpha \|=3`$. We also set $`\mu _{\mathrm{}}=0`$. Of course, $`\mu _𝒯`$ is a finitely supported signed measure on $``$ whose support is equal to $`I_{\alpha _0}`$, where $`\alpha _0`$ is the root of $`𝒯`$. We also observe that $`\|\mu _𝒯(\{n\})\|1`$, for all $`n`$. We finally set $`=\{\mu _𝒯:𝒯𝒢\}`$. Note that $`e_n^{}`$ as $`\left\{(0,\{n\},1)\right\}𝒢`$. We shall introduce some more notation in order to investigate properties of the set $``$. ###### Notation . Let $`𝒯𝒢`$ and let $`\alpha _0`$ denote its root. 1. Given $`\alpha 𝒯`$ set $`𝒯_\alpha =\{\beta \alpha ^{}:\beta 𝒯,\alpha \beta \}`$. Clearly, $`𝒯_\alpha 𝒢`$. 2. We let $`w(𝒯)=1`$, if $`\|𝒯\|=1`$. In case $`m_{\alpha _0}M`$, we set $`w(𝒯)=m_{\alpha _0}`$. 3. Let $`J`$. We let $`𝒯\|J`$ denote the tree resulting from $`𝒯`$ by keeping only those $`\alpha 𝒯`$ for which $`I_\alpha J\mathrm{}`$ and replacing $`I_\alpha `$ by $`I_\alpha J`$. It is easy to see that $`𝒯\|J𝒢`$. 4. We let $`𝒯`$ denote the tree resulting from $`𝒯`$ by changing $`ϵ_{\alpha _0}`$ to $`ϵ_{\alpha _0}`$. Clearly, $`𝒯𝒢`$ and moreover $`\mu _𝒯=\mu _𝒯`$. ###### Remark . Let $`𝒯𝒢`$. 1. If $`J`$, then $`\mu _{𝒯\|J}=\mu _𝒯\|J`$. 2. If $`\alpha 𝒯`$ then $`m(\alpha )ϵ(\alpha )\mu _𝒯\|I_\alpha =\mu _{𝒯_\alpha }`$. ###### Remark . Suppose $`𝒯_i𝒢`$, $`in`$. Let $`\alpha _i`$ be the root of $`𝒯_i`$, $`in`$. We shall say that $`\{𝒯_i:in\}`$ is $`S_\xi `$-admissible, $`\xi <\omega `$, if $`\{I_{\alpha _i}:in\}`$ is. We shall also write $`𝒯_1<\mathrm{}<𝒯_n`$ if $`I_{\alpha _1}<\mathrm{}<I_{\alpha _n}`$. It is easy to see that if $`𝒯_1<\mathrm{}<𝒯_n`$ is $`S_{n_j}`$-admissible then $`\frac{_{i=1}^n\mu _{𝒯_i}}{m_j}`$. It follows by our preceding remarks that $``$ is pointwise bounded, symmetric and closed under restriction to subsets of $``$. Hence $`(e_n)`$ is an $`1`$-unconditional, bimonotone basis for $`X_{}`$. It is not hard to check that $`X_{}`$ is isometric to $`T(\frac{1}{m_i},S_{n_i})_{i=1}^{\mathrm{}}`$. We also obtain by our preceding remarks that if $`(x_i)_{i=1}^k`$ is an $`S_{n_j}`$-admissible block basis of $`(e_n)`$ then $`_{i=1}^kx_i\frac{1}{m_j}_{i=1}^kx_i`$. Hence $`X_{}`$ is an asymptotic $`\frac{1}{m_j}`$-$`\mathrm{}_1^{n_j}`$ space. It follows that $`(e_n)`$ is boundedly complete. Let now $`\nu `$ be a $`w^{}`$-cluster point of $``$. Using the reflexivity argument of ( cf. also ), one obtains that for every $`ϵ>0`$ there exists $`k`$ such that $`\nu \|[e_i:ik]<ϵ`$. It follows from this that $`(e_n)`$ is shrinking and thus $`X_{}`$ is reflexive. ###### Remark . Suppose $`(u_n)`$ is a normalized block basis of $`(e_n)`$ and $`u`$ an $`(ϵ,n_j)`$ average of $`(u_n)`$. Then $`\frac{1}{m_j}u1`$. ###### Lemma 4.2. Let $`𝒯𝒢`$. Let $`F`$ be a subset of $`𝒯`$ consisting of pairwise incomparable nodes. Then $`\{I_\alpha :\alpha F\}`$ is $`S_p`$-admissible, where $`p=\mathrm{max}\{n(\alpha ):\alpha F\}`$. ###### Proof. By induction on $`o(𝒯)`$. If $`o(𝒯)=3`$ the assertion of the lemma is trivial. Assuming the assertion true when $`o(𝒯)<3k`$, $`k>1`$, let $`𝒯𝒢`$ with $`o(𝒯)=3k`$. If $`\|F\|=1`$ there is nothing to prove. So assume $`\|F\|2`$. Let $`\alpha _0`$ be the root of $`𝒯`$ and let $`w(𝒯)=m_i`$ for some $`i`$. We denote by $`D`$ the set of immediate successors of $`\alpha _0`$ in $`𝒯`$. Given $`\alpha D`$ let $`F_\alpha =\{\beta F:\alpha \beta \}`$. Because $`o(𝒯_\alpha )3k3`$ we can apply the induction hypothesis on $`𝒯_\alpha `$ and the set $`\{\beta \alpha ^{}:\beta F_\alpha \}`$ to deduce that the collection $`\{I_\beta :\beta F_\alpha \}`$ is $`S_{p_1}`$-admissible, where $`p_1=\mathrm{max}\{n(\beta \alpha ^{}):\beta F_\alpha \}`$. Since $`n(\beta \alpha ^{})=n(\beta )n(\alpha )`$ and $`n(\alpha )=n_i`$ whenever $`\alpha D`$, we obtain that $`\{I_\beta :\beta F_\alpha \}`$ is $`S_{pn_i}`$-admissible, for every $`\alpha D`$. But also, $`\{I_\alpha :\alpha D\}`$ is $`S_{n_i}`$-admissible whence $`\{I_\alpha :\alpha F\}`$ is $`S_p`$-admissible. ∎ To simplify our notation, we set $`f_j=f_j^N`$. We make the following observation: Let $`𝒯𝒢`$ and let $`\alpha 𝒯`$. Assume that $`m(\alpha )<m_j^3`$ and that all $`M`$-entries of $`\alpha ^{}`$ are smaller than $`m_j`$. Then $`n(\alpha )f_j`$. Our next lemma will be crucial for the proof of the main result. ###### Lemma 4.3 (Decomposition Lemma). Let $`𝒯_0𝒢`$. Let $`j`$ such that $`w(𝒯_0)<m_j`$. Then there exist an $`S_{f_j}`$-admissible subset $`𝒢_0`$ of $`𝒢`$ and a scalar sequence $`(\lambda _𝒯)_{𝒯𝒢_0}`$ in $`[1,1]`$ so that the following are satisfied: 1. $`\mu _{𝒯_0}=_{𝒯𝒢_0}\lambda _𝒯\mu _𝒯`$. 2. For each $`𝒯𝒢_0`$ at least one of the following hold: either $`w(𝒯)=1`$ (thus $`\mu _𝒯=\pm e_{𝒯(p)}^{}`$ for some $`𝒯(p)`$), or $`w(𝒯)m_j`$, or $`\|\lambda _𝒯\|\frac{1}{m_j^2}`$. ###### Proof. Let $`𝔅`$ denote the set of all branches of $`𝒯_0`$ (a branch is a maximal well ordered subset of $`𝒯_0`$). If $`w(𝒯_0)=1`$ the assertion is trivial. So assume that $`w(𝒯_0)=m_{i_0}`$ for some $`i_0<j`$. Given $`b𝔅`$ set $$\alpha ^1(b)=\mathrm{max}\{\beta b:m(\beta )<m_j^2\text{ and if }m_i\beta ^{}\text{ then }i<j\}.$$ Note that $`\alpha ^1(b)`$ is well defined and that $`(m_{i_0},I,ϵ)<\alpha ^1(b)`$ since $`i_0<j`$ ($`(m_{i_0},I,ϵ)`$ being the root of $`𝒯_0`$). Let us say that $`b𝔅`$ is of type $`1`$ if $`\alpha ^1(b)`$ is terminal in $`𝒯_0`$. If $`b`$ is not of type $`1`$ then it is of type $`2`$ (resp. $`3`$), if the last $`M`$-entry of $`\alpha ^1(b)`$ is greater than or equal (resp. smaller than) $`m_j`$. We then denote by $`\alpha ^2(b)`$ the immediate successor of $`\alpha ^1(b)`$ in $`b`$. We let $`A_1=\{\alpha ^1(b):b𝔅\text{ is of type }1\}`$, $`A_2=\{\alpha ^1(b):b𝔅\text{ of type }2\}`$ and $`A_3=\{\alpha ^2(b):b𝔅\text{ is of type }3\}`$. Observe that the following properties hold: 1. If $`\alpha A_3`$ then all $`M`$-entries of $`\alpha ^{}`$ are smaller than $`m_j`$, $`m(\alpha ^{})<m_j^2`$, yet $`m_j^2m(\alpha )<m_j^3`$. 2. If $`\alpha A_2`$, then $`\alpha `$ is non-terminal, all $`M`$-entries in $`\alpha ^{}`$ are smaller than $`m_j`$, the last $`M`$-entry of $`\alpha `$ is greater than or equal to $`m_j`$ and $`m(\alpha )<m_j^2`$. 3. If $`\alpha A_1`$ then $`\alpha `$ is terminal, all $`M`$-entries in $`\alpha ^{}`$ are smaller than $`m_j`$ and $`m(\alpha )<m_j^2`$. It is not hard to check now that $`A=_{t=1}^3A_t`$ consists of pairwise incomparable nodes of $`𝒯_0`$ and hence $`\{I_\alpha :\alpha A\}`$ consists of successive subsets of $``$. Moreover, $`I=\{I_\alpha :\alpha A\}`$. Because $`m(\alpha )<m_j^3`$ and all $`M`$-entries of $`\alpha ^{}`$ are smaller than $`m_j`$ whenever $`\alpha A`$, we obtain that $`n(\alpha )f_j`$ for all $`\alpha A`$. Lemma 4.2 now yields that $`\{I_\alpha :\alpha A\}`$ is $`S_{f_j}`$-admissible. Finally, we let $`𝒢_0=\{(𝒯_0)_\alpha :\alpha A\}`$. Since $`m(\alpha )ϵ(\alpha )\mu _{𝒯_0}\|I_\alpha =\mu _{(𝒯_0)_\alpha }`$, for all $`\alpha 𝒯`$, we set $`\lambda _{(𝒯_0)_\alpha }=\frac{1}{m(\alpha )ϵ(\alpha )}`$ for $`\alpha A`$. We can easily verify that the desired properties hold. ∎ In the sequel, we shall be using a variety of block bases of $`(e_n)`$. The support of each of them will always be taken with respect to $`(e_n)`$. ###### Lemma 4.4. Let $`(u_n)`$ be a normalized block basis of $`(e_n)`$. Let $`j`$, $`j2`$ and let $`u`$ be a generic $`(ϵ,f_j+1)`$ average of $`(u_n)`$ with $`ϵ<\frac{1}{2m_j}`$. Let $`i<j`$ and let $`𝒯_1<\mathrm{}<𝒯_t`$ in $`𝒢`$ be $`S_{n_i}`$-admissible. Then $`_{k=1}^t\mu _{𝒯_k}(u)2`$. In particular, $`\mu _𝒯(u)\frac{2}{w(𝒯)}`$, if $`w(𝒯)<m_j`$. ###### Proof. Observe that $`\frac{1}{m_j}_{k=1}^t\mu _{𝒯_k}`$ and hence $`_{k=1}^t\mu _{𝒯_k}(u_n)m_j`$, for all $`n`$. Let $`P=(p_n)`$, where $`p_n=\mathrm{min}\mathrm{supp}u_n`$. Set $`\xi =f_j+1`$ and suppose that $`u=_{n=1}^{\mathrm{}}\xi _1^R(p_n)u_n`$, for some $`R[P]`$. Let $`E^{}`$ denote the collection of the ranges of the $`\mu _{𝒯_k}`$’s, and let $`D`$ denote the collection of those $`u_n`$’s whose support intersects at least one member of $`E^{}`$. Put $`I_r=\{n:u_nD(E^{},r)\}`$, $`r=1,2`$. Because $`D(E^{},2)`$ is $`2S_{n_i}`$-admissible and $`n_if_j`$, we obtain that $`_{k=1}^t\mu _{𝒯_k}(_{nI_2}\xi _1^R(p_n)u_n)m_j2ϵ`$. On the other hand we clearly have that $`_{k=1}^t\mu _{𝒯_k}(_{nI_1}\xi _1^R(p_n)u_n)1`$. Thus, $`_{k=1}^t\mu _{𝒯_k}(u)2`$. ∎ ###### Lemma 4.5. Let $`(u_n)`$ be a normalized block basis of $`(e_n)`$. Let $`ϵ>0`$ and $`j`$. Then there exists a smoothly normalized $`(ϵ,f_j+1)`$ average of $`(u_n)`$. ###### Proof. Let $`P=(p_n)`$, where $`p_n=\mathrm{min}\mathrm{supp}u_n`$ for $`n`$. We can assume without loss of generality that $`\xi _1^R_{\xi 1}<ϵ`$ for every $`R[P]`$ where $`\xi =f_j+1`$. We are going to show that there exists a normalized block basis of $`(u_n)`$ admitting a generic $`(ϵ,\xi )`$ average of norm at least $`\frac{1}{2}`$. Suppose instead that this were false. Then it is easy to construct for every $`1rl_j`$, a block basis $`(u_i^r)`$ of $`(u_i)`$ so that letting $`p_i^r=\mathrm{min}\mathrm{supp}u_i^r`$ and $`P_r=(p_i^r)`$ the following are satisfied: 1. $`(u_i^r)`$ is a block basis of $`(u_i^{r1})`$. ($`u_i^0=u_i`$) 2. $`u_i^r=_{n=1}^{\mathrm{}}\xi _i^{P_{r1}}(p_n^{r1})\frac{u_n^{r1}}{u_n^{r1}}`$, for all $`i`$. ($`p_n^0=p_n`$) 3. $`u_i^r<\frac{1}{2}`$, for all $`i`$. 4. For every $`i`$, if $`u_i^r=_{nF_i^r}a_nu_n`$ with $`a_n>0`$ for $`nF_i^r`$, then $`_{nF_i^r}a_n2^{r1}`$ and $`(u_n)_{nF_i^r}`$ is $`S_{\xi r}`$-admissible. The construction is easily done by induction. Taking $`r=l_j`$ we see from 3. that $`u_i^{l_j}<\frac{1}{2}`$. On the other hand 4. implies that $`u_i^{l_j}\frac{2^{l_j1}}{m_j}`$ as $`\xi l_j<n_j`$. Thus, $`m_j>2^{l_j}`$ contradicting the choice of $`l_j`$. ∎ Our next lemma yields that $`X_{}`$ satisfies the $`(6,N,F,𝐚)`$ distortion property where $`F=(f_i+2)`$ and $`𝐚=(\frac{1}{m_i})`$. ###### Lemma 4.6. Let $`(u_j)`$ be a normalized block basis of $`(e_j)`$. Suppose that $`(y_j)`$ is a block basis of $`(u_j)`$ so that $`y_j`$ is a smoothly normalized $`(ϵ_j,f_j+1)`$ average of $`(u_j)`$ with $`ϵ_j<\frac{1}{2m_j}`$. Given $`j_0`$ and $`J_0[]`$, there exists $`J[J_0]`$ such that $`j_0<\mathrm{min}J`$ and for every $`𝒯𝒢`$, $`D_𝒯=\{y_j:jJ,\|\mu _𝒯(y_j)\|\frac{5}{m_{j_0}}\}`$ is $`S_{f_{j_0}+1}`$-admissible. ###### Proof. Note first that Lemma 4.5 guarantees the existence of the block basis $`(y_j)`$. Let $`P=(p_j)_{jJ_0}`$, where $`p_j=\mathrm{min}\mathrm{supp}y_j`$. By passing to a subsequence of $`(y_j)_{jJ_0}`$, if necessary, we can assume that the union of any $`4`$ $`S_{f_{j_0}}`$ subsets of $`P`$ belongs to $`S_{f_{j_0}+1}`$. Choose $`J[J_0]`$, $`J=(j_i)`$, such that $`j_0<j_1`$ and $`y_{j_i}_\mathrm{}_1<\frac{m_{j_{i+1}}}{m_{j_i}}`$, for every $`i2`$ (if $`v=_{i=1}^na_ie_i`$, then $`v_\mathrm{}_1=_{i=1}^n\|a_i\|`$). Let $`𝒯_0𝒢`$. Suppose first that $`w(𝒯_0)m_{j_0}`$. We show that in this case $`\|D_{𝒯_0}\|1`$. Indeed, suppose first that $`w(𝒯_0)<m_{j_1}`$. Lemma 4.4 yields that $`\|\mu _{𝒯_0}(y_j)\|\frac{4}{m_{j_0}}`$ for all $`jJ`$ whence $`D_{𝒯_0}=\mathrm{}`$. If $`w(𝒯_0)m_{j_1}`$ choose $`s2`$ so that $`m_{j_{s1}}w(𝒯_0)<m_{j_s}`$. Observe that if $`1i<s1`$ then $$\|\mu _{𝒯_0}(y_{j_i})\|\frac{1}{w(𝒯_0)}y_{j_i}_\mathrm{}_1<\frac{1}{w(𝒯_0)}\frac{m_{j_{s1}}}{m_{j_i}}<\frac{1}{m_{j_0}}.$$ When $`is`$, Lemma 4.4 yields $`\|\mu _{𝒯_0}(y_{j_i})\|\frac{4}{w(𝒯_0)}`$ $`<\frac{4}{m_{j_0}}`$. Hence $`D_{𝒯_0}\{y_{j_{s1}}\}`$ and so our claim holds. The final case to consider is that of $`w(𝒯_0)<m_{j_0}`$. Clearly, $`D_{𝒯_0}=\mathrm{}`$, if $`w(𝒯_0)=1`$. We employ the decomposition Lemma 4.3 to find an $`S_{f_{j_0}}`$ admissible subset $`𝒢_0`$ of $`𝒢`$ and scalars $`(\lambda _𝒯)_{𝒯𝒢_0}`$ satisfying the conclusion of Lemma 4.3. Let $`E^{}`$ denote the collection of the ranges of the $`\mu _𝒯`$’s ($`𝒯𝒢_0`$). Our previous work implies that $`D_{𝒯_0}(E^{},1)`$ is $`2S_{f_{j_0}}`$ admissible. But also, $`D_{𝒯_0}(E^{},2)`$ is $`2S_{f_{j_0}}`$ admissible since $`𝒢_0`$ is $`S_{f_{j_0}}`$ admissible. It follows that $`D_{𝒯_0}`$ is $`S_{f_{j_0}+1}`$ admissible. ∎ ###### Proof of Theorem 4.1. Let $`U`$ be a block subspace of $`X_{}`$ spanned by the normalized block basis $`(u_j)`$ of $`(e_j)`$. Let $`j_0`$ and choose a block basis $`(y_j)_{jJ}`$ of $`(u_j)`$ satisfying the conclusion of Lemma 4.6. Applying Corollary 3.4 of (cf. also Corollary 3.3 of ), we obtain that for every subsequence of $`(y_j)_{jJ}`$ which is a $`\delta `$-$`\mathrm{}_1^{f_{j_0}+2}`$ spreading model, it must be the case that $`\delta \frac{5}{m_{j_0}}`$ and thus $`\tau (U,\frac{6}{m_{j_0}})<f_{j_0}+2`$. ∎ Terminology. Let $`j_0`$ and $`(y_j)_{jJ}`$ satisfy the conclusion of Lemma 4.6. Every normalized $`(ϵ,n_{j_0})`$ average $`u`$ of $`(u_j)_{j=1}^{\mathrm{}}`$ of the form $`u=\frac{v}{v}`$, where $`v`$ is a generic $`(ϵ,n_{j_0})`$ average of $`(y_j)_{jJ}`$, will be called a normalized $`(ϵ,n_{j_0})`$ average of $`(u_j)_{j=1}^{\mathrm{}}`$ resulting from Lemma 4.6. Note that Lemmas 4.5 and 4.6 guarantee the existence of such averages for every block basis $`(u_j)_{j=1}^{\mathrm{}}`$. ###### Corollary 4.7. Let $`(y_j)`$ satisfy the assumptions of Lemma 4.6. Given $`j_0`$ and $`J_0[]`$, there exists $`J[J_0]`$ such that $`j_0<\mathrm{min}J`$ and for every $`𝒯_0𝒢`$, $`w(𝒯_0)m_{j_0}`$, $`D_{𝒯_0}=\{y_j:jJ,\|\mu _𝒯(y_j)\|\frac{5}{m_{j_0}m_e}\}`$ is $`S_{f_{j_0}+1}`$-admissible, where we have set $`m_e=\mathrm{min}\{m_{j_0},w(𝒯_0)\}`$. ###### Proof. We choose $`J_0`$ and $`j_0`$ as we did in the proof of Lemma 4.6. Suppose first that $`w(𝒯_0)>m_{j_0}`$. Because $`m_i^2<m_{i+1}`$, the argument in the proof of Lemma 4.6 shows that $`\|D_{𝒯_0}\|1`$. When $`w(𝒯_0)<m_{j_0}`$, we apply the decomposition Lemma 4.3 to find an $`S_{f_{j_0}}`$ admissible subset $`𝒢_0`$ of $`𝒢`$ and scalars $`(\lambda _𝒯)_{𝒯𝒢_0}`$ satisfying the conclusion of Lemma 4.3. Note that if $`𝒯𝒢_0`$ and $`w(𝒯)=m_{j_0}`$, then $`\|\lambda _𝒯\|\frac{1}{w(𝒯_0)}`$ and thus for all $`jJ`$, $`\|\lambda _𝒯\mu _𝒯(y_j)\|\frac{4}{m_{j_0}m_e}`$, by Lemma 4.4. Using the splitting argument of Lemma 4.6, we conclude that $`D_{𝒯_0}`$ is $`S_{f_{j_0}+1}`$-admissible. ∎ ###### Corollary 4.8. Let $`u`$ be a normalized $`(ϵ,n_{j_0})`$ average of $`(u_j)_{j=1}^{\mathrm{}}`$ resulting from Lemma 4.6 with $`ϵ\frac{1}{12m_{j_0}^2}`$. Let $`𝒢_0`$ be an $`S_{n_i}`$-admissible subset of $`𝒢`$, $`i<j_0`$, such that $`m_{j_0}\{w(𝒯):𝒯𝒢_0\}`$. Then, $`\|_{𝒯𝒢_0}\mu _𝒯(u)\|\frac{6}{m_e}`$, where $`m_e=\mathrm{min}\{w(𝒯):𝒯𝒢_0\}\{m_{j_0}\}`$. ###### Proof. Set $`\xi =n_{j_0}`$. Let $`p_j=\mathrm{min}\mathrm{supp}y_j`$, $`jJ`$ and $`P=\{p_j:jJ\}`$. There exists $`R[P]`$ so that $`u=\frac{v}{v}`$, where $`v=_{jJ}\xi _1^R(p_j)y_j`$ and $`\xi _1^R_{n_{j_0}1}<ϵ`$. Note that $`v\frac{1}{m_{j_0}}`$. Applying a splitting argument similar to that of Lemma 4.6 and taking in account Corollary 4.7, we obtain that $`\{y_j:jJ,\|_{𝒯𝒢_0}\mu _𝒯(y_j)\|\frac{5}{m_{j_0}m_e}\}`$ is $`3S_{2f_{j_0}+1}`$-admissible. The assertion follows from Lemma 4.4 and the fact that $`2f_{j_0}+1<n_{j_0}`$. ∎ ###### Remark . It is easy to see that in case $`w(𝒯)=1`$, for all $`𝒯𝒢_0`$, one obtains the estimate $`\|_{𝒯𝒢_0}\mu _𝒯(u)\|\frac{1}{m_{j_0}}`$. ## 5. Hereditarily indecomposable spaces This section is devoted to the proof of Theorem 3.5. Recall that $`X`$ is H.I. if and only if, for every pair of subspaces $`Y`$, $`Z`$ of $`X`$ and every $`ϵ>0`$, there exist $`yY`$ and $`zZ`$ so that $`yz`$ and $`yzϵy+z`$. Let $`M[]`$, $`M=(m_i)`$ and let $`N[]`$, $`N=(n_i)`$, which is $`M`$-good. Let $``$ be the set of measures constructed in the previous section by using the sets $`M`$ and $`N`$. We shall choose $`𝒩`$ so that the resulting space $`X_𝒩`$ is a reflexive H.I. space satisfying the conclusion of Theorem 3.5. We can find an injection $$\sigma :\{𝒯_1<\mathrm{}<𝒯_n:n,𝒯_i𝒢(in)\}\{m_{2j}:j\}$$ so that $`\sigma (𝒯_1,\mathrm{},𝒯_n)>w(𝒯_i)`$, for all $`in`$. ###### Definition 5.1. 1. An $`S_p`$-admissible sequence $`𝒯_1<\mathrm{}<𝒯_n`$ in $`𝒢`$ is said to be $`S_p`$-dependent, $`p`$, if $`w(𝒯_1)=m_{2j_1}`$, for some $`j_1>\frac{p}{2}`$, and $`\sigma (𝒯_1,\mathrm{},𝒯_{i1})=w(𝒯_i)`$, for all $`2in`$. 2. Let $`𝒯_1<\mathrm{}<𝒯_n`$ in $`𝒢`$, $`p`$ and $`𝒢_0𝒢`$. We shall say that $`𝒯_1<\mathrm{}<𝒯_n`$ admits an $`S_p`$-dependent extension in $`𝒢_0`$, if there exist $`l`$, $`k\{0\}`$ and an $`S_p`$-dependent sequence $`_1<\mathrm{}<_{n+k}`$ in $`𝒢_0`$ so that $`_{k+i}\|[l,\mathrm{})=𝒯_i`$, for all $`in`$. 3. A subset $`𝒢_0`$ of $`𝒢`$ is said to be self-dependent, if the following property is satisfied for every $`𝒯𝒢_0`$: Let $`\alpha 𝒯`$ so that its last $`M`$-entry equals $`m_{2j+1}`$ for some $`j`$. Let $`D_\alpha `$ denote the set of immediate successors of $`\alpha `$ in $`𝒯`$. Then $`\{𝒯_\beta :\beta D_\alpha \}`$ admits an $`S_{n_{2j+1}}`$-dependent extension in $`𝒢_0`$. ###### Definition 5.2. We let $`𝔇`$ denote the union of all non-empty, self-dependent, symmetric and closed under restriction to intervals, subsets of $`𝒢`$. Recall that $`𝒢_0𝒢`$ is symmetric if $`𝒯𝒢_0`$ whenever $`𝒯𝒢_0`$. $`𝒢_0`$ is closed under interval restrictions if $`𝒯\|J𝒢_0`$ whenever $`𝒯𝒢_0`$ and $`J`$ is an interval. Of course $`𝔇`$ is a maximal, under inclusion, subset of $`𝒢`$ with respect to the aforementioned properties. Set $`𝒩=\{\mu _𝒯:𝒯𝔇\}`$. We will show that $`X_𝒩`$ is H.I. ###### Remark . The maximality of $`𝔇`$ implies the following: 1. $`e_n^{}𝒩`$, for all $`n`$. 2. If $`𝒯𝔇`$, then $`𝒯_\alpha 𝔇`$, for all $`\alpha 𝒯`$ and so the decomposition Lemma 4.3 holds for $`𝔇`$. 3. If $`𝒯_1<\mathrm{}<𝒯_k`$ in $`𝔇`$ is $`S_{n_{2i}}`$-admissible, $`i`$, then $`\frac{\mu _{𝒯_1}+\mathrm{}+\mu _{𝒯_k}}{m_{2i}}𝒩`$. 4. If $`𝒯_1<\mathrm{}<𝒯_k`$ in $`𝔇`$ is $`S_{n_{2i+1}}`$-dependent, $`i`$, then $`\frac{\mu _{𝒯_1}+\mathrm{}+\mu _{𝒯_k}}{m_{2i+1}}𝒩`$. 5. Because of 3., all the results obtained in the previous section about $`(ϵ,\xi )`$ averages in $`X_{}`$, where $`\xi `$ is either $`n_j`$ or $`f_j+1`$ for some $`j`$, still hold in $`X_𝒩`$ provided $`j`$ is even. Note that $`X_𝒩`$ is reflexive by the same argument that showed $`X_{}`$ is reflexive. Thus $`(e_i)`$ is a shrinking basis for $`X_𝒩`$. ###### Proof of Theorem 3.5. It follows from Theorem 4.1 and our preceding remarks that $`X_𝒩`$ satisfies the $`(6,N^{(2)},F^{(2)},𝐚)`$ distortion property. We show that $`X_𝒩`$ is H.I. This is accomplished through Theorem 3.6. Let $`(u_n)`$ be a normalized block basis of $`(e_n)`$ and let $`j`$. Set $`P=(p_n)`$, where $`p_n=\mathrm{min}\mathrm{supp}u_n`$. We can assume that the union of any $`7`$ $`S_{f_{2j+1}}`$ subsets of $`P`$ belongs to $`S_{f_{2j+1}+1}`$. Successive applications of Corollary 4.8 yield a normalized block basis $`g_1<\mathrm{}<g_p`$ of $`(u_n)`$, $`𝒯_1<\mathrm{}<𝒯_p`$ in $`𝔇`$, and integers $`j_1<\mathrm{}<j_p`$ with $`2j+1<j_1`$, satisfying the following: 1. $`g_i`$ is a normalized $`(\frac{1}{12m_{2j_i}^2},n_{2j_i})`$ average of $`(u_n)`$ resulting from Lemma 4.6. 2. $`w(𝒯_i)=m_{2j_i}`$, $`\mathrm{supp}\mu _{𝒯_i}r(g_i)`$ and $`\mu _{𝒯_i}(g_i)>\frac{1}{2}`$, for all $`ip`$. 3. $`\sigma (𝒯_1,\mathrm{},𝒯_{i1})=w(𝒯_i)`$, for all $`ip`$. 4. $`\{g_i:ip\}`$ is maximally $`S_{n_{2j+1}}`$-admissible. Put $`\theta _i=(\mu _{𝒯_i}(g_i))^1`$, $`z_i=\theta _ig_i`$, and note that $`1\theta _i<2`$, $`ip`$. We’ll show that $`(z_i)_{i=1}^p`$ satisfies conditions 1. and 2. of Theorem 3.6, with $`\delta _j=\frac{1}{m_{2j+1}}`$, “$`n_j`$$`=n_{2j+1}`$ and $`k_j=f_{2j+1}+1`$. Condition 1. is immediate since $`𝒯_1<\mathrm{}<𝒯_p`$ is $`S_{n_{2j+1}}`$-dependent. Condition 2. is achieved by establishing the following Claim: Given $`𝒯𝔇`$, there exist intervals $`J_1<\mathrm{}<J_s`$ in $`\{1,\mathrm{},p\}`$ so that 1. $`\{z_{\mathrm{min}J_t}:ts\}`$ is $`S_{f_{2j+1}+1}`$-admissible. 2. $`\mu _𝒯\|\{z_i:iJ_t\}`$ is constant for all $`ts`$. 3. $`\|\mu _𝒯(z_i)\|<\frac{12}{m_{2j+1}^2}`$, for all $`i_{t=1}^sJ_t`$. To prove the claim suppose first that $`w(𝒯)>m_{2j+1}`$. Corollary 4.8 yields that $`\|\mu _𝒯(z_i)\|\frac{12}{m_{2j+1}^2}`$, for at most one $`ip`$, and thus the claim holds in this case. Next assume that $`w(𝒯)=m_{2j+1}`$. Without loss of generality, there exist an $`S_{n_{2j+1}}`$-dependent sequence $`_1<\mathrm{}<_l`$ in $`𝔇`$ and an interval $`J`$ so that $`\mu _𝒯=\frac{1}{m_{2j+1}}_{k=1}^l\mu __k\|J`$. Let $`i_0`$ be the largest $`i`$ for which $`w(𝒯_i)`$ is an element of $`\{w(_k):kl\}`$, or let $`i_0=0`$, if no such $`i`$ exists. The injectivity of $`\sigma `$ and Corollary 4.8 imply that if $`i_0\{0,1\}`$, or if $`w(𝒯_{i_0})=w(_1)`$, then $`\|\mu _𝒯(z_i)\|<\frac{12}{m_{2j+1}^2}`$, for all $`ii_0`$. If $`i_0>1`$, then the injectivity of $`\sigma `$ yields $`w(𝒯_{i_0})=w(_{i_0})`$, $`𝒯_i=_i`$ for $`i<i_0`$ yet $`𝒯_{i_0}_{i_0}`$. It follows now by Corollary 4.8, that $`\|\mu _𝒯(z_i)\|<\frac{12}{m_{2j+1}^2}`$, for all $`i>i_0`$. We also observe that there exists $`i_1<i_0`$ such that $`\mu _𝒯(z_i)=0`$, if $`i<i_1`$, while $`\mu _𝒯(z_i)=\frac{1}{m_{2j+1}}`$ if $`i_1<i<i_01`$. Concluding, there exist four intervals $`J_1<J_2<J_3<J_4`$ in $`\{1,\mathrm{},p\}`$, some of which may possibly be empty, such that $`\mu _𝒯\|\{z_i:iJ_t\}`$ is constant for every $`t4`$, while $`\|\mu _𝒯(z_i)\|<\frac{12}{m_{2j+1}^2}`$, for each $`i_{t=1}^4J_t`$. Finally, assume $`w(𝒯)<m_{2j+1}`$. If $`w(𝒯)=1`$, the claim trivially holds so suppose that $`w(𝒯)>1`$. Choose $`𝒢_0𝔇`$ $`S_{f_{2j+1}}`$-admissible and scalars $`(\lambda _{})_{𝒢_0}`$ according to the decomposition Lemma 4.3. By splitting the $`z_i`$’s into two sets, those whose support intersects at least two of the ranges of the $`\mu _{}`$’s, and those whose support intersects at most one, we deduce from our previous work that there exist intervals $`J_1<\mathrm{}<J_s`$ in $`\{1,\mathrm{},p\}`$ so that $`\{z_{\mathrm{min}J_t}:ts\}`$ is $`7S_{f_{2j+1}}`$-admissible, $`\mu _𝒯\|\{z_i:iJ_t\}`$ is constant for all $`ts`$, and $`\|\mu _𝒯(z_i)\|<\frac{12}{m_{2j+1}^2}`$, for all $`i_{t=1}^sJ_t`$. Thus the claim holds and the proof is complete. ∎
warning/0004/hep-th0004062.html	ar5iv	text	# M[any] Vacua of IIBWork supported by RFBR grant # 99-01-00190, INTAS grant # 950681, and Scientific School support grant 96-15-0628 ## Introduction IKKT, or IIB matrix model , is a statistical mechanic model which was introduced as a nonperturbative regularisation of IIB superstring model in the Schild gauge . It is the string analog of lattice gauge models and provides a tool for non-perturbative numerical study of the string theory (M-theory) . The picture one has in the IKKT model is closely related to the Connes’ approach to the noncommutative geometry . In this approach a manifold is described in terms of the so called Connes’ triple $`(𝒜,\mathrm{\Delta },)`$, rather than as a set of points. In this triple, $`𝒜`$ is the algebra of bounded operators (algebra of functions), $`\mathrm{\Delta }`$ is some elliptic differential operator (Laplace or Dirac operator, for simplicity we consider Euclidean signature) and $``$ is the Hilbert space where $`𝒜`$ and $`\mathrm{\Delta }`$ are represented. “Points” of such a manifold can be identified with the spectra of some position operators built from the triple, other local and global characteristics can as well be extracted from it. In the case of IKKT matrix model the rôle of $`𝒜`$ is played by Hermitian matrices with “smooth” eigenvalue distribution and bounded square trace in the limit $`N\mathrm{}`$, while the $`\mathrm{\Delta }`$ and $``$ are represented by the background solution and the adjoint representation of $`U(N)`$. Compactifications of this model to $`d`$-dimensional noncommutative tori was shown to result in noncommutative Yang–Mills (YM) models in respective dimensions . In particular, compactification to a circle yields the Banks–Fischler–Shenker–Susskind (BFSS) matrix model which was introduced as a non-perturbative regularisation of the light-cone membrane action in $`D=11`$, . The compactification in Ref. is introduced as a restriction of the background solution to satisfy some periodicity conditions modulo gauge transformations. The respective periods are identified with the periods of the torus on which the model is compactified. From the other hand, given a BPS background, i.e. a solution to equations of motion which preserves a part of the supersymmetry, one can map the space of $`N\times N`$ Hermitian matrices of IKKT model to the space of real or matrix valued functions on some non-commutative manifold. Under certain conditions this map is isomorphism of algebras where the matrix product is mapped into noncommutative or star product of functions due to the noncommutativity of the manifold. The properties of the manifold are exclusively determined by the respective BPS background solution. Explicitly the BPS solution is given by a set of matrices having scalar commutators. In the case when the matrix of commutators is nondegenerate (or for the subset on which it is nondegenerate) one can easily construct the map from the space of arbitrary Hermitian matrix fluctuations around respective BPS background to the space of functions on a noncommutative manifold . The limit of commutative manifold corresponds to infinite commutator of the solution. The case when the commutator is degenerate corresponds just to an opposite situation. One may argue, and this was the main reason for arriving at actual paper, that only commutative solutions give the global minimum of the IKKT actions. However, as we show in the Section 2., this is not exactly the case, since due to homogeneity of the IKKT action in bosonic fields, the respective action vanishes on purely bosonic solutions. Unfortunately, this regretful error committed in the pioneering work , was reproduced in the succeeding papers<sup>1</sup><sup>1</sup>1See e.g. most recent paper .. In spite of this disappointing discovery the study of the spectrum of fluctuations around a commutative solution still presents some interest, in particular due to the fact that, as it will be seen below, the degeneracy in the commutator of the solution corresponds to increasing the dimensionality of corresponding noncommutative YM model, which gives a “physically” different model (compare this with the situation with the first and second class constraints, ). To obtain a noncommutative $`U(1)`$ YM model from the matrix fluctuations one has to impose some irreducibility condition, while allowing certain degeneracy to the background solution one comes to description of a “vector fibre bundle” over the noncommutative manifold which leads to “nonabelian” noncommutative YM. Respective degeneracy can be interpreted as compactification on many coinciding branes, while the nondegenerate case corresponds to a single brane . The objective of the actual paper is to analyse the relation between the background commutative/degenerate vacuum solution and the properties of the resulting YM model. Although, the commutative solution is a particular case of a generic BPS solution, as we already mentioned, it is “physically” different, since it corresponds to a singular limit of the generic case. The plan of the paper is as follows. In the next section we give a brief account of IIB matrix model. After that we consider a BPS solution, show that it corresponds to a vanishing action and consider in more details the commutative case. We impose a set of conditions such a solution must respect, and build explicitly the map between the matrix fluctuations and functions of noncommutative manifold which appears to be a “noncommutative” product of two commutative manifolds dual to each other. This allow to find in the Section 4. a representation of IIB matrix model in terms of YM fields for both commutative and generic degenerate case. Finally, we discuss the results and consider the consequences and possible generalisations of the actual analysis. ## 1 The IIB Matrix Model The IKKT, or IIB matrix model, is described by the classical action: $$S=\frac{1}{g^2}tr\left(\frac{1}{4}[A_\mu ,A_\nu ]^2+\frac{1}{2}\overline{\psi }\mathrm{\Gamma }^\mu [A_\mu ,\psi ]\right),$$ (1) where $`A_\mu `$, $`\mu =1\mathrm{}D=10`$, $`\psi `$ and $`\overline{\psi }`$ are $`N\times N`$, $`N\mathrm{}`$ Hermitian matrices which act on $`N`$ dimensional Hilbert space $``$, with scalar product $`\eta ^{}\eta `$, $`\eta `$. Spinor matrices $`\overline{\psi }`$ and $`\psi `$ also carry the $`SO(10)`$ Majorana–Weyl spinor index, while bosonic ones $`A_\mu `$ carry vector one. We assume summation convention for repeated indices. Equations of motion corresponding to the action (1) look as follows $`[A_\mu ,[A_\mu ,A_\nu ]][\overline{\psi },\mathrm{\Gamma }_\mu \psi ]`$ $`=0,`$ (2) $`\mathrm{\Gamma }^\mu [A_\mu ,\psi ]=[\overline{\psi },A_\mu ]\mathrm{\Gamma }^\mu `$ $`=0.`$ (3) This model possesses a $`𝒩=2`$ supersymmetry corresponding to the transformations, $`\delta _{(1)}A^\mu `$ $`=i\overline{ϵ}\mathrm{\Gamma }^\mu \psi ,`$ (4) $`\delta _{(1)}\psi `$ $`={\displaystyle \frac{i}{2}}[A_\mu ,A_\nu ]\mathrm{\Gamma }^{\mu \nu }ϵ,`$ (5) and, $`\delta _{(2)}A_\mu =0,`$ (6) $`\delta _{(2)}\psi =\xi ,`$ (7) where $`ϵ`$ and $`\xi `$ are Majorana–Weyl spinors, as well as a SU(N) gauge symmetry, $`A_\nu U^1A_\nu U,`$ (8) $`\psi U^1\psi U,`$ (9) $`\overline{\psi }U^1\overline{\psi }U,`$ (10) and $`U`$ is $`N\times N`$ unitary matrix. Since only the adjoint representation is in use this symmetry is $`SU(N)/_N`$ rather that $`U(N)`$. This will not remain true if “matter” fields are included, in this case $`U(1)`$ and $`_N`$ will act nontrivially on the fields in the fundamental representation of $`U(N)`$ Another symmetry is already mentioned $`SO(10)`$ “covariance” of 10-dimensional vectors and spinors. ## 2 The Vacuum In the limit $`N\mathrm{}`$ equations (2,3) have an important class of purely bosonic ($`\psi =0`$) solutions which preserve a part of the supersymmetry (BPS solutions). These solutions are given by $`A_\mu =p_\mu `$, where matrices $`p_\mu `$ satisfy, $$[p_\mu ,p_\nu ]=iB_{\mu \nu },$$ (11) with $`B_{\mu \nu }`$ as $`u(N)`$ matrix proportional to the unity one, $`B_{\mu \nu }𝕀`$. For finite $`N`$ the commutator (11) cannot be satisfied but only mod some matrix vanishing in the sense of operator norm as $`N\mathrm{}`$. A solution of the described type can be realised e.g. by shift-and-clock operators on functions defined on a lattice . Naively, action (1) computed on the solution (11) should be equal to , $$S_{BPS}=\frac{N}{4g^2}(B_{\mu \nu })^2.$$ (12) As one can see this expression does not realise the minimum (even local one) unless $`B_{\mu \nu }=0`$, since one can devaluate the action (12) by a smooth variation of the solution, e.g. $`\delta p_\mu =\lambda p_\mu `$, which decreases the action if $`\lambda <0`$. The problem is accomplished by the fact that the variation of the action is first order in $`\lambda `$, which should not be the case if one varies a background solution. Indeed, this problem is solved if we see that for any finite $`N`$ the action computed on a purely bosonic solution of equations of motion vanishes. Thus, for any finite $`N`$ the bosonic part of the action (1) can be rewritten as, $$S=\frac{1}{4g^2}trA_\nu [A_\mu ,[A_\mu ,A_\nu ]],$$ (13) which vanishes identically if $`A_\mu `$ is a solution to (2) with zero $`\psi `$. The case with infinite $`N`$, however, is still indefinite since it depends how the limit $`N\mathrm{}`$ is achieved. One can obtain the limit either from $`p_\mu (N)`$ satisfying equations of motion or from ones not satisfying them, in the last case the value of action depends on how fast $`p_\mu (N)`$ approaches a solution as $`N`$ goes to infinity. Thus, for a configuration $`p_\mu (N)`$, satisfying, $$[p_\mu ,p_\nu ]=iB_{\mu \nu }𝕀(N),$$ (14) where $`𝕀(N)=𝕀+ϵ(N)`$ is a unity matrix approaching sequence. In order to approach (11), one has to require, $`trϵ(N)=N`$ (15) $`\underset{N\mathrm{}}{lim}ϵ(N)=0,`$ (16) where the limit is computed using the operator norm, $`ϵ=sup_{\eta ^{}\eta =1}\sqrt{\eta ^{}ϵ^{}ϵ\eta }`$, $`\eta `$. Eqs. (15) and (16) are contradictory unless we restrict the space of vectors $`\eta `$ on which the operator norm is computed e.g. by vectors having a fixed finite number of nonzero coordinates in the limit $`N\mathrm{}`$. In this case, the action (1) computed on on this sequence of configurations reads, $$S_{BPS}(N)=\frac{1}{4g^2}B^2tr𝕀(N)^2=\frac{1}{4g^2}B^2(trϵ(N)^2N),$$ (17) where the factor $`(trϵ(N)^2N)`$ is ambiguous. The additional requirement that for all $`N`$ the configuration $`p_\mu `$ to be a solution to (2) eliminates the ambiguity since for any $`N`$ one has $`S_{BPS}(N)=0`$, so we assume this requirement to be satisfied. The vanishing of the classical BPS action results in “equality” in the factor $`e^S`$ between solutions $`p_\mu `$ with different commutators (11), either BPS or non-BPS. The BPS vacua are, however, preferred since one expects to have no loop corrections to them . The above problems are absent in the case of commutative solutions, i.e. when $`B_{\mu \nu }=0`$. Before considering this case in more details, consider a generic solution (11). By a proper linear transformation $`B_{\mu \nu }`$ can be brought to the “standard” form having block diagonal form with two dimensional blocks of the form, $$\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),$$ and a $`r\times r`$ zero block, corresponding to zero modes of $`B_{\mu \nu }`$. The set of matrices $`(p_\mu )`$ is split in this case in three subsets $`(p_i,q^i,p_a)`$, where $`p_i`$ and $`q^i`$ form canonical conjugate pairs, and $`p_a=z_a^\mu p_\mu `$, $`a=1,\mathrm{},r=corankB`$, correspond to zero modes $`z_a^\mu `$ of $`B_{\mu \nu }`$, $`B_{\mu \nu }z_a^\nu =0`$, $`[p_i,q^j]=i\delta _i^j,`$ (18) $`[p,p]=[q,q]=0.`$ (19) Intuitively, at this stage one can see that $`p_\mu `$ corresponding to a nondegenerate part of $`B_{\mu \nu }`$ carries $`1`$ (one-particle Hamiltonian) degree of freedom while one corresponding to zero mode caries $`2`$ degrees. The second degree of freedom should be given by the canonical conjugate $`q^a`$, if it exists. We hope this statement will become more transparent below. For a nonzero $`B_{\mu \nu }`$, as well as for zero one it can be seen that the set of adjoint operators $`P_\mu =[p_\mu ,]`$, (which act on $`u(N)`$ matrices), is a commutative one, $$(P_\mu P_\nu P_\nu P_\mu )a=[p_\mu ,[p_\nu ,a]][p_\nu ,[p_\mu ,a]]=i[B_{\mu \nu },a]=0.$$ (20) This fact common with one that $`P_\mu `$ are selfadjoint with respect to the scalar product, $$(a,b)=tra^{}b,$$ (21) where $`a^{}`$ stands for the hermitian conjugate matrix, means that $`P_\mu `$ can be diagonalised even for noncommutative set of matrices $`p_\mu `$ satisfying eq. (11). One can, therefore, decompose any hermitian $`N\times N`$ matrix in the orthogonal basis of $`P_\mu `$ eigenmatrices. Having this decomposition at hand we will construct in the next section the map from the space of Hermitian matrices to the space of real functions on some noncommutative manifold. In the case when $`B_{\mu \nu }=0`$ one can diagonalise not only adjoint operators $`P_\mu `$ but also all matrices $`p_\mu `$. We say that $`p_\mu `$ is a nondegenerate commutative vacuum background solution if the following are satisfied: * All $`p_\mu `$ are commutative, $$[p_\mu ,p_\nu ]=0.$$ (22) * Matrices $`p_\mu `$ are (functionally) independent, this implies also linear independence, $$\alpha ^\mu p_\mu =0\text{all }\alpha ^\mu =0.$$ (23) * They satisfy, $`trp_\mu `$ $`=0,`$ (24) $`trp_\mu p_\nu `$ $`=0,\text{for }\mu \nu .`$ (25) * Eigenvalues of $`p_\mu `$ form a $`D`$-dimensional irregular lattice with sites symmetrically distributed with respect to the origin (the centre of the lattice), in the range $`\|\lambda \|\mathrm{\Lambda }`$<sup>2</sup><sup>2</sup>2One can alternatively chose eigenvalues to lie in other compact domain, e.g. a box, $`\mathrm{\Lambda }_\mu \lambda _\mu \mathrm{\Lambda }_\mu `$., where $`\|\lambda \|=\sqrt{\lambda _\mu ^2}`$. In the limit $`N\mathrm{}`$ the lattice becomes dense. This implies that one can rearrange the matrix labels (which is equivalent to a gauge transformation) in such a way that eigenvalues $`\lambda _\mu `$ are distributed according to “smooth” and “monotonic” functions of the matrix index $`s_\nu `$, $`\lambda _\mu (s_\nu )`$, which is the $`s_\nu `$-th eigenvalue, where $`s_\nu `$ form a linear $`D`$-dimensional lattice, $`det\left({\displaystyle \frac{\lambda _\mu (s)}{s_\nu }}\right)>0,`$ (26) $`\lambda _\mu (s)=\lambda _\mu (s),`$ (27) $`\|\lambda _\mu (s+1)\lambda _\mu (s)\|=\left\|{\displaystyle \frac{\lambda _\mu (s)}{s_\nu }}\right\|=o(N^1),`$ $`N\mathrm{},\mathrm{\Lambda }=\text{fixed}.`$ (28) The last gives the smooth distribution of eigenvalues in the limit $`N\mathrm{}`$. In this limit quantity $`\mathrm{\Lambda }`$ plays the role of UV cutoff. This implies that for a finite but large values of $`N`$ one is justified to manipulate with $`\lambda _\mu (s)`$ as with quasicontinuous quantities. The authentic continuum limit is achieved by sending $`\mathrm{\Lambda }`$ to infinity after the limit $`N\mathrm{}`$ is reached. The eigenvalue problem in the case of commuting $`p_\mu `$ is equivalent to the Cartan decomposition problem in Lie algebra $`u(N)`$, see e.g. . In this context the above enlisted conditions say that $`p_\mu `$ must form the orthogonal basis in a $`D`$-dimensional plane of Cartan subalgebra of $`su(N)`$ to which no root is orthogonal. ## 3 The Map As usual any matrix $`a`$ is a function of a pair of labels $`(s,s^{})`$, $`a=a(s,s^{})`$. Matrix product and trace are given by, respectively, $`(ab)(s^{},s^{\prime \prime })`$ $`={\displaystyle \underset{\{s_\mu \}}{}}a(s^{},s)b(s,s^{\prime \prime }),`$ (29) $`tra`$ $`={\displaystyle \underset{\{s_\mu \}}{}}a(s,s).`$ (30) In what follows we will use also quasicontinuous notations where the sums we will substitute by the integrals. Taking into account that the increment of $`s`$ in (29) and (30) is $`ds=\mathrm{\Delta }s=1`$ these sums can be written as following integrals, $`(ab)(s^{},s^{\prime \prime })`$ $`={\displaystyle 𝑑sa(s^{},s)b(s,s^{\prime \prime })},`$ (31) $`tra`$ $`={\displaystyle 𝑑sa(s,s)}.`$ (32) Where the range of integration is given by $`\frac{N_\mu }{2}s_\mu \frac{N_\mu }{2}`$ and $`_\mu N_\mu =N`$. The unity matrix is given by the $`\delta `$ symbol, $$𝕀(s,s^{})=\delta _{ss^{}}\delta (ss^{}).$$ (33) Note, that in these notations $`\delta (0)=1`$, this differs from the usual “continuous” $`\delta `$-function by a diverging factor of order $`N`$. The limit $`N\mathrm{}`$ is achieved in the following way. First, one introduce an UV cutoff $`L`$ and rescale $`s\frac{2L}{N}s`$. Thus the increment $`\mathrm{\Delta }s`$ becomes really small and one has the genuine integration in eqs. (29-32). The cutoff removing is obtained in the limit $`L\mathrm{}`$. Since we plan to identify $`p_\mu `$ with momentum operators the continuous $`\lambda _\mu `$ spectrum corresponds to noncompact manifolds, while to have the compact result one has to keep the discreteness of the spectrum and sums instead of the integrals. In this picture matrices $`p_\mu `$ look as follows, $$p_\mu (s,s^{})=\lambda _\mu (s)\delta (ss^{}),$$ (34) where $`\lambda _\mu (s)`$ scan the eigenvalue lattice of $`p_\mu `$, while an arbitrary diagonal matrix $`\varpi `$ looks like, $$\varpi (s,s^{})=\varpi (s)\delta (ss^{}).$$ (35) It is not difficult to compute the action of operator $`P_\mu `$ on an arbitrary matrix $`a(s,s^{})`$, $$\begin{array}{c}(P_\mu a)(s,s^{})=[p_\mu ,a](s,s^{})=\hfill \\ \hfill 𝑑s^{\prime \prime }(\lambda _\mu (s)\delta (ss^{\prime \prime })a(s^{\prime \prime },s^{})a(s,s^{\prime \prime })\lambda _\mu (s^{\prime \prime })\delta (s^{\prime \prime }s^{}))\\ \hfill =(\lambda _\mu (s)\lambda _\mu (s^{}))a(s,s^{}).\end{array}$$ (36) One immediately finds that the spectrum of $`P_\mu `$ is given by eigenvalues $`k_\mu =\lambda _\mu (t)\lambda _\mu (t^{})`$ for any $`t`$ and $`t^{}`$ taking values in the lattice $`\{s_\mu \}`$, and corresponding to eigenvectors, $`f_{tt^{}}(s,s^{})`$ $`=\delta (st)\delta (s^{}t^{}),`$ (37) $`P_\mu f_{tt^{}}`$ $`=\left(\lambda _\mu (t)\lambda _\mu (t^{})\right)f_{tt^{}}.`$ (38) For large values of $`N`$ eigenvalues $`k`$ become degenerate, the respective eigenfunctions maybe labelled by eigenvalue $`k_\mu `$ and a label $`t`$ counting the degenerate states, $$f_t(k)=𝑑t^{}\delta (t^{}s(\lambda (s)k))f_{tt^{}}=f_{ts(\lambda (s)k)},$$ (39) where the function $`s(\lambda )`$ is the inverse to the $`\lambda (s)`$: $`s(\lambda (s))=s`$, $`\lambda (s(\lambda ))=\lambda `$. Now, consider the matrix $`q^\mu `$ defined by, $$q^\mu =i\frac{}{k_\mu }dtf_t(k)\|_{k=0}=iA_\alpha {}_{}{}^{\mu }(s)\delta ^{}{}_{}{}^{\alpha }(s^{}s),$$ (40) where, $$A_\alpha {}_{}{}^{\mu }(s)=\frac{s_\alpha }{\lambda _\mu }\|_{\lambda =\lambda (s)}=\left(\left(\frac{\lambda }{s}\right)^1(s)\right){}_{\alpha }{}^{}{}_{}{}^{\mu },$$ (41) and $`\delta ^{}{}_{}{}^{\alpha }(s)=\frac{}{s_\alpha }\delta (s)`$. As one can immediately see the commutator of $`p_\nu `$ with $`q^\mu `$ is the canonical one, $$[p_\mu ,q^\nu ]=i\delta _\mu ^\nu 𝕀,$$ (42) i.e. $`p_\mu `$ and $`q^\nu `$ form a canonical conjugate pair and can be interpreted as, respectively, momentum and coordinate operators, as we anticipated earlier. Strictly speaking, such pairs of operators exist only in the limit $`N\mathrm{}`$, i.e. when the spectrum of either $`q`$ or $`p`$ is continuous. From $`q`$’s one can construct a matrix $`E(k)=e^{ik_\mu q^\mu }`$. Matrix $`E(k)`$ is nondegenerate and it is an eigenvector of $`P_\mu `$ corresponding to the value $`k`$, $$[p_\mu ,E(k)]=k_\mu E(k),$$ (43) having the squared norm, $$E(k)^2=trE^{}(k)E(k)=N.$$ (44) For a finite $`N`$ the dimensionality of $`k`$-eigenspace depends on $`k`$ and decrease as $`k`$ approaches to $`2\mathrm{\Lambda }`$ since matrix edges are approaching for large values of $`k`$. When $`N\mathrm{}`$ this difference disappears and one can regard the spaces with different $`k`$ as isomorphic. This approximation is accurate as soon as $`\|tt^{}\|N`$, where $`\lambda (t)\lambda (t^{})=k`$. Matrix $`E(k)`$ can be considered as the eigenvalue shift operator, $$E(k)f(k^{})f(k+k^{}).$$ (45) Due to isomorphism and nondegeneracy of $`E(k)`$, eigenvectors corresponding to a nonzero $`k`$ are given by the product of zero vectors and matrix $`E(k)`$ as follows, $$E_t(k)=f_t(0)E(k),$$ (46) where $`f_t(0)`$ is the basis in the space of diagonal matrices which form the zero space. It is transparent that this basis is orthogonal in the space of $`N\times N`$ if $`f_t(0)`$ form an orthonormal basis in the space of diagonal matrices, the norm of $`E_t(k)`$ coinciding with that of $`f_t(0)`$. Thus, given an orthonormal basis in the zero space of $`P_\mu `$ one can spread it out through all $`k`$-eigenspaces using eq. (46). Since eigenvectors of $`P_\mu `$ form a complete orthonormal set in the space of $`N\times N`$ matrices, any Hermitian matrix $`a`$ can be expanded in the basis of $`P_\mu `$ eigenvectors, $$a=\underset{t,k}{}\stackrel{~}{a}_t(k)f_t(k)=\underset{t,k}{}\stackrel{~}{a}_t(k)f_t(0)e^{ikq}.$$ (47) Orthonormality of the spectrum assures the inverse transformation, $$\stackrel{~}{a}_t(k)=\frac{1}{f_t(0)^2}tre^{ikq}f_t^{}(0)a.$$ (48) From the properties a vacuum solution one can trace that arbitrary matrix commuting with $`p_\mu `$ is a diagonal matrix (35), and it may be represented as a function of $`p_\mu `$, $$\varpi =\varpi (p)\varpi (s(p)),$$ (49) where $`\varpi (s(\lambda ))\|_{\lambda =p}`$ is the same function as in (35), but computed on diagonal matrices $`p_\mu `$. One may expand $`\varpi `$ in multiple ways e.g. in Fourier series as function of $`p_\mu `$, $$\varpi =\underset{z}{}\stackrel{~}{\varpi }(z)e^{ipz},$$ (50) where $`z_\mu `$ are points of the lattice of eigenvalues of operator $`Q^\mu =[q^\mu ,]`$, $$[q^\mu ,e^{ipz}]=z^\mu e^{ipz}.$$ (51) Since operators $`Q^\mu `$ are also commutative selfadjoint operators (they commute also with $`P_\mu `$) their spectrum satisfy, $$tre^{ip(zz^{})}=N\delta _{zz^{}},$$ (52) which gives the ground for inverse transformation, $$\stackrel{~}{\varpi }(z)=\frac{1}{N}tre^{ipz}\varpi .$$ (53) Using the above decomposition for the diagonal matrices one can rewrite the expansion eq. (47) for arbitrary Hermitian matrix $`a`$ as follows, $`a={\displaystyle \underset{z,k}{}}\stackrel{~}{a}(z,k)e^{ikz}e^{ipz}e^{ikq}={\displaystyle \underset{z,k}{}}\stackrel{~}{a}(z,k)e^{ipz+ikq},`$ (54) $`\stackrel{~}{a}(k,z)={\displaystyle \frac{1}{N}}trae^{ipzikq},`$ (55) $`\stackrel{~}{a}^{}(k,z)=\stackrel{~}{a}(k,z).`$ (56) At this point one can identify the space of Hermitian matrices with the space of real functions on the “noncommutative manifold” $`Spec(P)\times Spec(Q)`$, where $`Spec(P)`$ and $`Spec(Q)`$ are the varieties of eigenvalues of $`P`$, and respectively of $`Q`$. In fact when $`N\mathrm{}`$ with $`\mathrm{\Lambda }`$ fixed $`Spec(P)`$ form a compact manifold, while $`Spec(Q)`$ tends to a noncompact lattice. Continuum limit for this lattice is achieved in the limit $`\mathrm{\Lambda }\mathrm{}`$. This map is given by<sup>3</sup><sup>3</sup>3We use the same character to denote both the matrix and corresponding function the difference being that for functions we write explicitly its arguments, i. e. if $`a`$ denote some matrix we write $`a(z,l)`$ for the corresponding function and viceversa., $$aa(x,l)=\underset{z,k}{}\stackrel{~}{a}(z,k)e^{ilz+ikx},$$ (57) where $`\stackrel{~}{a}(z,k)`$ is defined by eq. (55). This map can be turned backward, therefore, it is one-to-one. The corresponding matrix is given by, $$a=\underset{k,z}{}\stackrel{~}{a}(k,z)e^{ipz+ikq},$$ (58) now, $`\stackrel{~}{a}(z,k)`$ is the Fourier transform of the function $`a(x,l)`$, $$\stackrel{~}{a}(z,k)=\frac{1}{N}\underset{x,l}{}a(x,l)e^{ilzikx}.$$ (59) Under the correspondence (5759) the matrix product is mapped into the star product given by, $$(ab)ab(x,l)=e^{\frac{i}{2}\left(\frac{^2}{x^{}l}\frac{^2}{xl^{}}\right)}a(x,l)b(x^{},l^{})\|_{\genfrac{}{}{0pt}{}{x^{}=x}{l^{}=l}},$$ (60) while the trace corresponds to the lattice integration over $`x`$ and $`l`$, $$tra\underset{\genfrac{}{}{0pt}{}{lSpec(P)}{xSpec(Q)}}{}a(x,l)𝑑x𝑑la(x,l).$$ (61) Commutators with $`p_\mu `$ and with $`q^\mu `$ give the lattice analogs of differentiation with respect to $`x^\mu `$ and respectively $`l_\mu `$, $`[p_\mu ,a]i{\displaystyle \frac{}{x^\mu }}a(x,l),`$ (62) $`[q^\mu ,a]i{\displaystyle \frac{}{l_\mu }}a(x,l).`$ (63) ## 4 The Spectrum Let us now return back to the Matrix Model action (1) and consider arbitrary fluctuations around the vacuum solution $`p_\mu `$ satisfying (2225), $$A_\mu =p_\mu +ga_\mu ,$$ (64) here $`a_\mu `$ is an arbitrary hermitian matrix. Perturbed action looks as follows, $$S=tr\left(\frac{1}{4}([p_\mu ,a_\nu ][p_\nu ,a_\mu ]+g[a_\mu ,a_\nu ])^2+\overline{\psi }\mathrm{\Gamma }^\mu [(p_\mu +ga_\mu ),\psi ]\right).$$ (65) Now, let use the correspondence between matrices and functions to map the fluctuation $`a_\mu `$ to function $`a_\mu (x,l)`$. Then, action (65) is rewritten as follows, $$S=d^Dxd^Dl\left(\frac{1}{4}_{\mu \nu }^2(x,l)+\overline{\psi }\mathrm{\Gamma }^\mu _\mu \psi (x,l)\right),$$ (66) where $`d^Dx`$ and $`d^Dl`$ are invariant measures on $`Spec(Q)`$ and $`Spec(P)`$ respectively. They are given by eigenvalue distribution densities for $`P`$ and $`Q`$. Also, $`_{\mu \nu }=i{\displaystyle \frac{}{x^\mu }}a_\nu (x,l)i{\displaystyle \frac{}{x^\nu }}a_\mu (x,l)g[a_\mu ,a_\nu ]_{}(x,l)`$ (67) $`_\mu \psi (x,l)=i{\displaystyle \frac{}{x^\mu }}\psi (x,l)g[a_\mu ,\psi ]_{}(x,l),`$ (68) where $`[,]_{}`$ stands for the star commutator defined as, $$[a,b]_{}=abba.$$ (69) Action (66) give the exact description of IIB matrix model in $`N\mathrm{}`$ limit in terms of functions on the manifold $`Spec(P)\times Spec(Q)`$. This action possesses a huge gauge symmetry given by the following transformations by a star-unitary function $`U(x,l)`$, $`a_\mu (x,l)`$ $`U^1a_\mu U(x,l){\displaystyle \frac{i}{g}}U^1{\displaystyle \frac{}{x^\mu }}U(x,l)`$ (70) $`\psi (x,l)`$ $`U^1\psi (x,l),`$ (71) where $`U^1(x,l)=U^{}(x,l)`$, and star conjugate function is given by the function corresponding to the Hermitian conjugate matrix and in this particular case coincides with the complex conjugate function, $$U^{}(x,l)=(U^{})(x,l)=U^{}(x,l).$$ (72) Global gauge group $`G`$ can be identified with “constant” transformations $`U`$, satisfying $`_\mu U=0`$. This means that $`U`$ may depend only on the momentum parameter $`l`$, i.e. $`U=U(l)`$. At first sight it seems that the gauge group is $`U(1)`$ group localised along $`Spec(P)`$, since gauge symmetry is realised by scalar functions and not matrix valued ones. Indeed, the “global” group is Abelian, $$U(l)U^{}(l)=U(l)U^{}(l)=U^{}(l)U(l),$$ (73) but this commutativity does not hold for local transformations. Consider the algebra of infinitesimal gauge transformations $`𝔊`$. This is algebra of real functions with star commutator, $$[f,g]_{}=ih,f,g,h𝔊.$$ (74) Natural basis for this algebra form following functions, $$L_{k,z}(x,l)=e^{ilz+ikx}.$$ (75) The commutator for generators $`L_{k,z}`$ looks as follows, $$[L_{k,z},L_{k^{},z^{}}]_{}=2i\mathrm{sin}\frac{1}{2}(k^{}zkz^{})L_{k+k^{},z+z^{}}.$$ (76) Thus, one may interpret the model (66) as an ordinary (commutative) Yang–Mills model with the local algebra of gauge transformations (76). It is worthwhile to note that although the model we got is defined in the terms the flat space, the gauge symmetry the model has contains the group of local $`x`$-reparameterisations as a subgroup, whose infinitesimal transformations are generated by, $`L_ϵ=ϵ^\mu (x)l_\mu `$, $$[L_ϵ,f(x)]=ϵ^\mu (x)_\mu f(x),$$ (77) but this is not all. It is not difficult to show that algebra (76) contains Virasoro subalgebra which may serve as an indication that action (66) also describe the string spectrum. Indeed, consider a functions $`\theta (x,l)`$ defined mod $`2\pi `$ and its canonical conjugate $`w(x,l)`$, $$[w,\theta ]_{}=i.$$ (78) Then, generators $`L_n=e^{in\theta }w`$ satisfy, $`L_n^{}=L_n`$ (79) $`[L_n,L_m]_{}=i(nm)L_{n+m},`$ (80) which is exactly the classical Virasoro algebra. The above results can be readily extended to the case of a generic background commutator $`B_{\mu \nu }`$. In this case we assume the subset $`(p_i,p_a)`$ to satisfy conditions (22)—(28). Operators $`q^i`$, and respective eigenfunctions are already known, one has only to find the counterparts to $`p_a`$. Repeating the derivations of the previous section for this particular case, one comes to the map from matrices to functions $`f(x^\mu ,l^a)`$, with star product defined as, $$(ab)(x,l)=e^{\frac{i}{2}\stackrel{~}{C}^{\mu \nu }\frac{^2}{x^\mu x^\nu }\frac{i}{2}\left(\frac{^2}{x^al_a}\frac{^2}{x^al_a^{}}\right)}a(x,l)b(x^{},l^{}),$$ (81) tensor $`\stackrel{~}{C}^{\mu \nu }`$ is defined by, $$\stackrel{~}{C}^{\mu \nu }B_{\nu \alpha }=B_{\alpha \nu }\stackrel{~}{C}^{\nu \mu }=\mathrm{\Pi }_\alpha ^\mu ,$$ (82) where $`\mathrm{\Pi }_\alpha ^\mu `$ is the projector to the space orthogonal to zero modes of $`B_{\mu \nu }`$. The YM action (66) keeps in this case the same form except the integration is done over $`d^Dxd^rl`$, and the star product is given by eq. (81). As we see, zero modes of $`B_{\mu \nu }`$ lead to extra integrations in (66) over $`dl`$, or extra gauge (and, respectively, physical) degrees of freedom. ## 5 Generalisations and Conclusions In analysing the spectrum of fluctuations around a vacuum solution of the classical action (1) we made several assumptions about the respective solutions. Namely, we required $`p_\mu `$ to form an independent and nondegenerate commutative set of matrices. This lead us to the action (66) for fields on the product of the manifold corresponding to the spectrum of coordinate operator and its dual given by the spectrum of momentum operator. This representation contains a broad symmetry including invariance with respect to local reparameterisations as well as $`2D`$ conformal symmetry, and can be interpreted as a Yang Mills model with gauge algebra given by (76). The above assumptions seem natural, but one may pose a question: what may happen if one gives up some of them? Consider first the case when all conditions are respected except the number of independent matrices is not equal to $`D`$ but is smaller. Let, in particular, matrices $`p_\mu `$ be expressed as linear combinations of the independent subset $`(p_\alpha )`$ $`\alpha =1,\mathrm{},p+1<D`$, $$p_\mu =\xi _\mu ^\alpha p_\alpha ,$$ (83) where $`rank\xi _\mu ^\alpha =p+1`$. In this case, one can find the position operators corresponding to independent $`p_\alpha `$, and a generic hermitian matrix will be expandable in the basis of the matrices $`e^{ip_\alpha z^\alpha +ik_\alpha q^\alpha }`$. Therefore, after manipulations like in previous section, hermitian matrices are now represented by the real functions defined on the spectrum of the subset of independent $`p_\alpha `$ and $`q^\alpha `$, $`\alpha =1,\mathrm{},p+1`$. As a result, one has action (65) mapped into a $`(p+p)`$-dimensional action which corresponds to the reduction of the action (66) of $`(D+D)`$-dimensional YM model to a $`(p+p+2)`$ dimensional plane given by $`\xi _\mu ^\alpha `$. One may describe the above situation as a localisation of IIB matrix model to a $`p`$-brane in contrast to $`(D1)`$-brane we had initially in (66). Consider now the opposite case, i.e. when all matrices $`p_\mu `$ together fail to possess nondegenerate set of eigenvalues. The last means that i) there exist eigenvalue sets $`(\lambda _\mu )`$ whose eigenspace are degenerate and, therefore, ii) there are matrices $`\pi `$ commuting with all $`p_\mu `$ but not being functions of $`p_\mu `$. It is not difficult to show that the set of these matrices generate at most central extended Lie algebra $`𝔏`$, $$[\pi _a,\pi _b]=ic_{ab}𝕀+f_{ab}^c\pi _c.$$ (84) Thus, one is able map the generic Hermitian matrix $`a`$ from $`u(N)`$ algebra to a $`𝔏`$-valued function $`a(x,l)`$. In particular, when $`𝔏=u(n)`$ this corresponds to a “nonabelian” generalisation of the action (66). The “physical” interpretation of this solution is the localisation of IIB matrix model to $`n`$ copies of coinciding branes, . All above says that the spectrum of the matrix model is in a strong dependence of the vacuum solution chosen. Due to the limit $`N\mathrm{}`$, the different background solutions $`p_\mu `$ lead to different continuum models (66). One may conjecture that the respective models may be classified by the rank $`r`$ of $`B_{\mu \nu }`$ and intersection of stabiliser groups of each $`p_\mu `$, or by the symmetry of the vacuum. In this paper we considered the solution with the smallest possible group of symmetry giving the $`D`$-dimensional action. In fact, this configuration has the “largest” measure (entropy factor), or better to say “the moduli space” among the commutative solutions. The opposite extreme is given by the “solution” where all $`p_\mu `$ are proportional to unity matrix $`p_\mu =\lambda _\mu 𝕀`$ which describes the $`u(N)`$ YM model localised on a point, i.e. the original matrix model. The “moduli space” of such solution is parameterised by just $`D`$ numbers $`\lambda _1,\mathrm{},\lambda _D`$. Thus, in IIB matrix model one has a plenty of vacua. For a finite $`N`$ these vacua are connected through the fluctuations corresponding to zero modes of the Hessian matrix $`S_{\mu \nu }\frac{^2S}{A_\mu A_\nu }(p)`$. As $`N\mathrm{}`$, we expect such fluctuations to fall out of the class of allowed functions, namely continuous $`L^2`$-integrable functions, and the vacua to become separated. ###### Acknowledgments. I am grateful to P. Pyatov, A. Nersesian, and T. Bakeev for useful discussions and critical remarks.
warning/0004/astro-ph0004084.html	ar5iv	text	# Magnetic Flux Expulsion in the Powerful Superbubble Explosions and the 𝛼⁢"-"⁢Ω Dynamo. ## 1 Introduction The magnetic field in our Galaxy and in other spiral galaxies is usually believed to have been amplified from a weak seed field by a hydromagnetic dynamo, which exists due to the presence of the large-scale differential rotation and small-scale cyclonic turbulence in the Galaxy (Parker 1970, 1971; Vainshtein $`\&`$ Ruzmaikin 1971, 1972; Moffat 1978). It has been suggested that any primordial magnetic field could be expelled from the Galaxy by the dynamic motions in less than a billion years (Parker 1971) so it would seem that some flux amplification is necessary to explain the Galactic field. The theory of such a dynamo has been formulated in a precise way through the mean field equations and the solutions to these equations indicate that the field would be amplified. The resulting field appears to correspond to the magnetic field patterns in our galaxy and others. On the other hand, a number of criticisms of this theory have emerged. One of them concerns the intense development of the small scale fields, which could damp the turbulence and stop the dynamo action unless they saturate at levels which do not interfere with mean field dynamo (Kulsrud $`\&`$ Anderson 1992; Parker 1992; Vainshtein $`\&`$ Cattaneo 1992). Another unresolved problem is the expulsion of flux from the galactic disc. A very important point is that the theory of the $`\alpha \text{-}\mathrm{\Omega }`$ dynamo predicts the amplification of some small preexisting magnetic field only if some magnetic diffusion is present in the Galaxy. However, too large a diffusion is destructive for the dynamo, because of its dissipative role in the process of generation and it seems attractive to suppose that a dynamo without diffusion at all will be the most effective. But Ruzmaikin, Shukurov, $`\&`$ Sokoloff (1988) showed that this is not possible. The physical reason for this is rooted in the very strong flux freezing of the galactic plasma, because magnetic lines cannot break and the number of field lines in the disc can be increased only by toroidal stretching, which is accomplished by the dynamo action. Any stretching creates field of both signs, to conserve the total flux, and, for net amplification to occur, those portions of field lines which are of the wrong sign must be expelled from the Galaxy. In the standard $`\alpha \text{-}\mathrm{\Omega }`$ dynamo theory this is done by magnetic diffusion. Ruzmaikin, Shukurov, $`\&`$ Sokoloff (1988) considered $`\alpha \text{-}\mathrm{\Omega }`$-dynamo in the case of the thin disc and demonstrated that temporal evolution of the magnetic flux is governed by the following set of equations: $$\frac{}{t}\underset{0}{\overset{1}{}}B_r(t,z)𝑑z=\beta \frac{B_r}{z}\|_0^1,$$ (1) $$\frac{}{t}\underset{0}{\overset{1}{}}B_\phi (t,z)𝑑z=\beta \frac{B_\phi }{z}\|_0^1+G\underset{0}{\overset{1}{}}B_r𝑑z,$$ (2) where $`G=r\mathrm{\Omega }/r`$-measure of the differential rotation, $`\beta `$-magnetic diffusion, and $`0`$ and $`1`$ correspond to the center plane and boundary of the Galactic disc. The $`\beta `$ terms represent the expulsion of flux. It is clear, that setting $`\beta =0`$ in (1) and (2) we immediately get that $$\underset{0}{\overset{1}{}}B_r𝑑z=const,$$ (3) and that $`\underset{0}{\overset{1}{}}B_\varphi 𝑑z`$ may grow only linearly due to the stretching of lines by the galactic differential rotation. Thus, there is no exponential field growth, essential to the dynamo. One can easily see that for the dynamo to operate, there must be some nonzero flux escape through the upper boundary of the disc, that is $`\beta 0`$. This is consistent with the topological constraint that the total number of lines of force including those negative lines expelled from the disc must be constant. Indeed, a simple estimate of the expulsion terms making use of the numerical results in Ruzmaikin, Shukurov, $`\&`$ Sokoloff (1988) shows that the negative flux expelled during one growth cycle is comparable to the positive flux in the disc at the beginning of the e-folding. Although the $`\alpha \text{-}\mathrm{\Omega }`$ dynamo theory is complicated, this physical intuition of flux expulsion can be considered in the absence of these complications. Further, the mechanism of expulsion need not be tied to the $`\beta `$ diffusion inside the disc. The main problem with the expulsion of flux is that this flux is loaded with matter so that it is related to the expulsion of matter against the strong gravity of the galactic disc. The most likely process to expel flux is the phenomenon of sequential supernova (SN) or superbubbles (SB), which sweep up matter into dense, radiatively cooled shells. Magnetic field, tied to the matter due to the strong flux freezing in the ISM, is also swept up and deposited in these shells. If some part of the shell leaves the Galaxy, it carries the frozen-in magnetic field with it, thus producing the flux expulsion. But most of these superbubbles are not powerful enough to expel matter out of the gravitational well of the disc. The only possibility for flux expulsion seems to be: as the bubble expands, the field lines in the shells of SBs are not horizontal but form arcs, along which matter can slide down, lowering the amount of matter on the top of the lines and allowing some flux to escape. There is further difficulty with the mechanism involving the SBs which is relevant to its application to the $`\alpha \text{-}\mathrm{\Omega }`$-dynamo. The $`\alpha \text{-}\mathrm{\Omega }`$-dynamo assumes small-scale turbulence while the cavities produced by SB explosions may reach $`500`$ pc or larger, which is greater than some of the length scales of the galactic disk. For the dynamo theory in its conventional form to be applicable it is important that turbulence be small scale, because it involves the expansion of the turbulent electromotive force $``$, which describes the effect of turbulent motions on the mean (or ensemble- averaged) magnetic field, in terms of the mean magnetic field itself and its spatial derivatives: $$_i=\alpha _{ij}<B_j>+\beta _{ijk}\frac{<B_j>}{x_k}$$ (4) (Moffat 1978) If the scale of the turbulence is too large, then the expansion is invalid and usual $`\alpha \text{-}\mathrm{\Omega }`$-dynamo theory must be modified. That’s why, for example, direct application of $`\alpha `$ and $`\beta `$ tensors calculated by Ferri$`\stackrel{`}{\mathrm{e}}`$re (1995, 1998) for SBs and SNs to the $`\alpha \text{-}\mathrm{\Omega }`$-dynamo theory can lead to an overly optimistic estimates of the rate of flux escape. Although $`\alpha \text{-}\mathrm{\Omega }`$-dynamo theory is not strictly applicable to the case of SBs, the actual operation of them in amplifying the field is clear from the work of Ferri$`\stackrel{`}{\mathrm{e}}`$re (1991, 1995, 1998). It is also clear that the rapid escape of the flux from the disk is essential. In this paper we show that because of the deep gravitational well of the Galaxy it is difficult for the matter and field lines to escape, and consequently for the mean field to grow. In Ferri$`\stackrel{`}{\mathrm{e}}`$re’s works she finds the lines of force rising with the SB but does not follow them long enough to see that they must fall back into the disk and inhibit the growth of the field. In this note we quantitatively examine the dynamics of the rising field lines carried by SBs and show that even with sliding the matter and flux are unlikely to escape. Thus, the requirement of the escape of flux provides a strong constraint for the $`\alpha \text{-}\mathrm{\Omega }`$ dynamo to overcome if it is to amplify the galactic magnetic field. ## 2 Sliding of matter from the top of SB; formulation of the problem Multiple supernovae from OB associations can carve out large cavities of hot gas, called superbubbles. McCray $`\&`$ Snow (1979) first described them. When the SB expands into surrounding medium it sweeps up interstellar matter, giving rise to a massive expanding shell. Inside the volume surrounded by this shell a hot rarefied low-density gas is contained which provides the pressure driving further expansion of the shell. The energetic source for sustaining this pressure is provided by continuous energy input from the SN explosions in the center of SB. Weaver et al. (1977) calculated the evolution of the bubble driven by the continuous wind from the central source and Mac Low $`\&`$ McCray (1988) applied this theory to the case of supernovae driven SBs. They show that the radius of such a SB, expanding in a uniform medium of density $`\rho _0`$ with continuous energy input in the center $`L_{SN}=L_{38}10^{38}\mathrm{ergs}\mathrm{s}^1`$ (the luminosity of SB conveniently expressed in units of $`10^{38}\mathrm{ergs}\mathrm{s}^1`$), is given by $$R(t)=\left(\frac{125}{154\pi }\right)^{1/5}L_{SN}^{1/5}\rho _0^{1/5}t^{3/5}=267\left(\frac{L_{38}t_7^3}{n_0}\right)^{1/5}\mathrm{pc},$$ (5) where $`t_7=t/10^7`$ yr, with the velocity of the envelope changing as $$u(t)=\dot{R}(t)15.7\left(\frac{L_{38}}{n_0t_7^2}\right)^{1/5}\mathrm{km}\mathrm{s}^1.$$ (6) The inner pressure in the volume bounded by the shell varies as $`P_{in}(t)={\displaystyle \frac{7}{(3850\pi )^{2/5}}}L_{SN}^{2/5}\rho _0^{3/5}t^{4/5}`$ $`=4.1\times 10^{12}\left({\displaystyle \frac{L_{38}^2n_0^3}{t_7^4}}\right)^{1/5}\mathrm{ergs}\mathrm{sm}^3,`$ (7) due to the work done on the expansion of the shell and the energy injection in the center of SB. When the shell expands in the real Galactic environment, there is also a gravitational force which tends to slow down the vertical expansion. Also, the distribution of medium is highly inhomogeneous with height $`z`$ over the Galactic plane. Various components of ISM have different length scales and characteristic densities, but they all have exponential or Gaussian decreasing profiles in $`z`$ and thus drop very rapidly with height. This effect is very noticeable for powerful superbubbles, for which the expansion radius may exceed the height scale of the matter distribution. As we will see this can lead to Rayleigh-Taylor instabilities. At the final stage of the SB expansion, its velocity becomes comparable to the sound velocity of the ambient ISM and shock wave will no longer exist. However, this does not prevent the SB from further expansion, because there is still a lot of momentum in its massive shell of swept up material and the ram pressure at high $`z`$ is negligible. This means that expansion actually continues until the velocity of the shell drops to zero under the action of Galactic gravity: $$u=0.$$ (8) This stopping condition is different from the one which Ferri$`\stackrel{`}{\mathrm{e}}`$re used (1998) in her investigation of the role of SBs in the $`\alpha \mathrm{\Omega }`$ dynamo. Her condition (expansion stops when the shell velocity is of the order of speed of sound in the ambient ISM, in other words when the pressure inside the SB becomes comparable to the ram pressure) is applicable only for massless shells without any inertia. In reality shell is quite massive and the momentum stored in it drives further expansion of the bubble against the ram pressure and the galactic gravity. Also when the SB expands in the presence of the gravity, the swept-up matter in its shell is likely to slide downward along the shell as mentioned before. This can influence the efficiency of flux removal by SBs. Obviously, in the absence of sliding, escaping matter takes with it all the frozen-in magnetic flux. Now, if we allow the matter in the shell to slide downward perpendicular to the field $`B`$ in the shell it will take the magnetic field lines away from the top of the SB, thus reducing the magnetic flux to be removed. But it can also have a significant effect on the dynamics of the SB itself, because as the matter slides from the top in any direction, the upper part of the shell becomes lighter and acceleration due to the inner pressure of the hot gas will gets larger in proportion to the surface density decrease, while the gravitational deceleration stays the same. This results in some additional acceleration of the top of the shell and, in principle, this can significantly change the dynamics of this part of SB if there is enough time before the shell stops. If this happens, and a substantial part of the mass slides down, the top of the shell may continue its expansion upward and drive the remaining matter and frozen-in flux to a further distance from the galactic plane and possibly expel the flux entirely from the Galaxy. For this reason it is very important to estimate the numerical value of this effect. ## 3 Basic equations Giuliani (1982) gave a general formulation of the thin-shell approximation for hypersonic, hydromagnetic flows, axisymmetric about the $`z`$ axis, including motions along the shell. We use his equations in our description of gravitational matter removal from the top of SB. All our further consideration is restricted to the case when the angular distance of the shell element from the shell top, $`\theta `$, is very small, $`\theta 1`$. We suppose the expansion of the SB is described by some law $`R=R(t)`$. We will also suppose for simplicity that the form of the shell near its top can be roughly approximated by a sphere. The validity of this assumption will be discussed later. We made some further simplifications, one of which is the neglect of the pressure gradient along the shell. It is clear enough that any such pressure gradient would only reduce the downsliding of the matter. This means that our estimate is only an upper limit of the sliding and the real sliding will actually be smaller. We also completely neglect the influence of the magnetic field on the dynamics of the sliding. This assumption is justified in the early periods of the Galaxy’s life, when the magnetic field was weak. At the present time this is not completely valid, because the magnetic field is strong and magnetic tension may play some dynamical role in the expansion process. SBs with strong magnetic field were considered analytically by Ferri$`\stackrel{`}{\mathrm{e}}`$re (1991) and numerically by Tomisaka (1992). Thus, the only external force in our analysis is the gravity due to the stars and ISM in the Galaxy. We take gravity as given, because the self-gravity of the bubble is negligible. With this in mind the system of the equations of Giuliani describing the gravitational fall of matter from the top of the shell reduce to $$\frac{}{t}(R^2\mathrm{sin}\theta \sigma )=\rho _0uR^2\mathrm{sin}\theta \frac{}{\theta }\left(R(t)\mathrm{sin}\theta \sigma v_{}\right),$$ (9) $$\frac{v_{}}{t}+\frac{\rho _0uv_{}}{\sigma }+\frac{v_{}}{R}\left(u+\frac{v_{}}{\theta }\right)g\mathrm{sin}\theta =0.$$ (10) Here $`v_{}`$ is the tangential velocity of the matter along the shell arising from the presence of gravity, $`u=\dot{R}(t)`$ is the velocity of the shell expansion (dot means time derivative), $`\sigma `$ is the surface density of the shell, $`\rho _0`$ is the density of the unperturbed gas in front of the shell, and $`g`$ is the gravitational acceleration. Now, at an early stage of SB expansion, when its radial velocity is very high and the transverse velocity due to the galactic gravity is not very large, the effect of sliding is negligible, because the total velocity of the shell element is directed almost radially. Thus, it is reasonable to suppose that the most noticeable effect of sliding will occur only during the later stages of the shell evolution, when it sufficiently slows down in the radial direction. But by the onset of this later stage the SB has expanded in the direction perpendicular to the galactic plane to a distance larger than the scale height of matter distribution. This reasoning allows us to neglect the second term proportional the ambient density in equation (10) . It makes equation (10) completely independent of equation (9) since it then contains no terms in $`\sigma `$. However, in equation (9), we must keep terms in $`\rho _0`$ because $`\sigma `$ still grows due to swept up matter. Also we can neglect the term $`v_{}/\theta `$ in (10) compared to the radial velocity $`u`$. We discuss this omission later. Thus equations (9) and (10) reduce to the following simplified system of equations: $`{\displaystyle \frac{}{t}}(R^2(t)\mathrm{sin}\theta \sigma )=\rho _0uR^2(t)\mathrm{sin}\theta `$ $`R(t)\sigma {\displaystyle \frac{}{\theta }}\left(\mathrm{sin}\theta v_{}\right)R(t)\mathrm{sin}\theta v_{}{\displaystyle \frac{\sigma }{\theta }},`$ (11) $$\frac{v_{}}{t}+v_{}\frac{u}{R(t)}=g\mathrm{sin}\theta .$$ (12) In general, we assume that gravitational acceleration, $`g`$, is a function of the $`z`$-coordinate in the Galaxy. Bearing in mind that $`u=\dot{R}(t)`$ and $`v_{}=0`$ at $`t=0`$, we can integrate equation (12) to get $`v_{}(t)={\displaystyle \frac{\mathrm{sin}\theta }{R(t)}}{\displaystyle \underset{0}{\overset{t}{}}}R(t^{})g(t^{})𝑑t^{}`$ $`={\displaystyle \frac{\mathrm{sin}\theta }{R(t)}}{\displaystyle \underset{0}{\overset{t}{}}}R(t^{})g(z_0+R(t^{}))𝑑t^{},`$ (13) where $`z_0`$ is the height in the Galactic plane where the explosion occured. We see that $`v_{}\mathrm{sin}\theta `$. This means that in equation (11) we may neglect the last term near the top of SB since it is proportional to $`\mathrm{sin}^2\theta `$. Then equation (11) further reduces to $`\dot{\sigma }+{\displaystyle \frac{1}{R}}\left[2\dot{R}+{\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{}{\theta }}\left(\mathrm{sin}\theta v_{}\right)\right]\sigma =\rho _0\dot{R}.`$ (14) This equation can be integrated with the initial condition $`\sigma =0`$ at $`t=0`$: $`\sigma (t)={\displaystyle \frac{1}{R^2(t)}}\mathrm{exp}\left({\displaystyle \underset{0}{\overset{t}{}}}\kappa (t^{})𝑑t^{}\right)`$ $`\times {\displaystyle \underset{0}{\overset{t}{}}}R^2\dot{R}\rho (z_0+R(t^{}))\mathrm{exp}\left({\displaystyle \underset{0}{\overset{t^{}}{}}}\kappa (t^{\prime \prime })dt^{\prime \prime }\right)dt^{},`$ (15) where $$\kappa =\frac{1}{\mathrm{sin}\theta }\frac{}{\theta }\left(\mathrm{sin}\theta v_{}\right).$$ (16) Substitution of expression (13) into (16) gives $`\kappa =2\mathrm{cos}\theta {\displaystyle \frac{1}{R^2(t)}}{\displaystyle \underset{0}{\overset{t}{}}}R(t^{})g(t^{})𝑑t^{}`$ $`{\displaystyle \frac{2}{R^2(t)}}{\displaystyle \underset{0}{\overset{t}{}}}R(t^{})g(z_0+R(t^{}))𝑑t^{},`$ (17) near the top of SB. From the formulae (15) and (16) it is easy to see that the importance of sliding near the top is determined by the quantity $$\zeta =e^\kappa =\mathrm{exp}\left(\frac{2}{R^2(t)}\underset{0}{\overset{t}{}}R(t^{})g(t^{})𝑑t^{}\right),$$ (18) which can be determined for any given $`R(t)`$. If $`\zeta 1`$, we can safely neglect the sliding of matter but if $`\zeta 1`$, sliding will play important role in the dynamics of late stages of SB expansion. ## 4 Influence of sliding on the shell expansion Let us examine the equation for $`R(t)`$ near the top of the bubble taking into account the effect of the unloading the matter from the shell’s top on the radial expansion of the SB itself. To do this we consider a solid angle $`d\mathrm{\Omega }`$ of the shell near its top. We can write the following equation for its motion: $$\frac{d}{dt}\left(\sigma \dot{R}R^2d\mathrm{\Omega }\right)=\left(P_{in}P_{out}\right)R^2d\mathrm{\Omega }R^2\sigma gd\mathrm{\Omega }.$$ (19) We can combine this equation with equation (15) to find $`R(t)`$ and $`\sigma (t)`$ near $`\theta =0`$ as a functions of time $`t`$ if we know how $`P_{in}`$ behaves. All other quantities, namely $`P_{out}`$ and $`g`$ are given empirically, as a function of $`z`$, by the position of the shell’s top. Equations (19) and (14) can be combined to give $$\ddot{R}+\dot{R}^2\frac{\rho _0}{\sigma }\frac{2\dot{R}}{R^2}\underset{0}{\overset{t}{}}R(t^{})g(t^{})𝑑t^{}=\frac{P_{in}P_{out}}{\sigma }g$$ (20) and together with equation (15) give $`R(t)`$. These equations are solved numerically to get the behavior of $`R`$. First we supposed that the inner pressure is governed by equation (7). This underestimates the pressure and the height reached. But if there is no escape of flux under this assumption there is certainly no escape in real conditions. We have taken the model of ISM from Ferri$`\stackrel{`}{\mathrm{e}}`$re (1998), that is we supposed that density and pressure of ISM are contributed by $`5`$ components, having different number densities, length scales and temperatures: neutral, cold, warm, ionized, and hot. We also use her approximation for the gravitational acceleration $`g`$. We suppose for simplicity that the luminosity of the SB is constant in time during 37 Myr, until the death of the $`8M_{}`$ stars (Ferri$`\stackrel{`}{\mathrm{e}}`$re 1995), and equal to $$L=10^{36}\times N\mathrm{erg}\mathrm{s}^1,$$ (21) where $`N`$ is the number of SNs in the star cluster. For this calculation we suppose that the inner pressure $`P_{in}`$ changes in accordance with equation (7) during the first $`37`$ Myr of the SB expansion, with $`\rho _0`$ the mass density at the site of explosion. After $`37`$ Myr the interior cools adiabatically because the inner pressure and inertia of the shell continue to drive the shell expansion so that the volume bounded by the shell increases. We also carry out calculations for the different pressure law. But we neglect radiative cooling of the hot gas in the interior of the bubble during the entire explosion, so that real SB expansion is always smaller than we obtain here. In Figure 1a we show the dependence of $`R`$ upon time $`t`$ for the case of a SB with $`N=75`$ SNs in it, going off near the Sun at an initial Galactic altitude $`100`$ pc. For comparison we also depict the curve without sliding for the same SB, which is obtained by setting $`g`$ in the integral in the left-hand side of equation (20) to zero (but not in the right-hand side!). If the sliding is taken into account, the time when expansion stops is $`t_s=35`$ Myr and $`R_{max}=1336`$ pc; neglecting sliding we get $`t_s=30`$ Myr and $`R_{max}=1094`$ pc. It is obvious that effect of sliding should be more pronounced in more powerful bubbles because of their longer lifetimes and stronger gravity at the heights to which they bring the matter. To illustrate this we plot in Figure 1b the maximum radii of SBs of various luminosities with and without sliding for the same conditions as the Figure 1a. It is clearly seen that sliding plays important role for powerful SBs, making their final radii dozens of percent larger than that without sliding. In Figure 2 we plot, for comparison, the expansion velocity of the $`N=75`$ SB top and the velocity of sliding along the shell for the same SB, divided by $`\mathrm{sin}\theta `$ (which virtually equals to $`v_{}/\theta `$ near the SB top). One can see that, up to the first $`20`$ Myr, $`u`$ is larger than $`v_{}/\theta `$, which justifies our neglect of corresponding term in equation (10). By that time most of the SB expansion has already occured, so the inclusion of the term with $`v_{}/\theta `$ into (12) does not change the final results significantly. Moreover, if we do include it, it would only suppress sliding, as can be seen from (10), so that our results for sliding and shell expansion can be considered to give upper bounds for the final height of the shell. ## 5 Rayleigh-Taylor instability It is tempting to expect that for even more powerful SBs, than $`N=100`$ SNs, the effect of sliding will be even more pronounced and the final size of the SB may reach scales comparable to the size of Galaxy, thus making possible the escape of the flux. But when the luminosity of the SB approaches $`10^{38}\mathrm{ergs}\mathrm{s}^1`$, corresponding to the number of SNs $`N100`$, another effect becomes important for the fate of SB. At such a large luminosity the shell starts accelerating some time before its expansion could be stopped by the gravity and it would accelerate, in principle, to a very high velocity if there is enough time for it. However, it was first noticed by Mac Low $`\&`$ McCray (1988) and then proven numerically by them (Mac Low $`\&`$ McCray 1989) that as soon as the shell starts to accelerate it becomes Rayleigh-Taylor unstable and eventually breaks up. They called this process a “blowout” of the shell into the halo. Indeed, the effective gravity in the moving shell in the case of acceleration is directed towards the center of SB, that is the dense cold gas in the shell is pushed by rarefied hot gas of the interior and this leads to the instability. The shell fragments into blobs of cold, dense gas, which continue moving with velocities they attained before fragmentation. There is no further significant acceleration of the shell, because the hot rarefied gas from the interior of the SB escapes into halo thus allowing the inner pressure to drop. After that time each blob will move ballistically, sweeping up some small mass, though this effect might be not so important at high altitudes, where the density of all the components of the ISM is very low. This means that if the speed of the blob at this time is less than the escape velocity from the Galaxy, matter can not leave but must return to the Galactic plane from halo and there is no contribution to the flux escape. For this instability the galactic gravity contributes to the effective gravity as well, thus making the shell unstable even when it is still decelerating. In our situation, when we have a semi-infinite hot rarefied medium and a dense slab of gas of finite thickness $`H`$ up on it, the fastest growing mode of the instability has a scale of the order of $`H`$, so that the growth rate of the instability is given by $$\gamma ^2=\left(\ddot{R}+g\right)/H.$$ (22) The increment of instability or the amount of e-foldings reduces to $$\gamma 𝑑t=𝑑t\sqrt{\left(g+\ddot{R}\right)/H}.$$ (23) For the fragmentation to proceed effectively we require $$\gamma 𝑑t>1.$$ (24) To get the shell thickness, let us note that the gas initially located between the heights $`z_0+z`$ and $`z_0+z+dz`$ is deposited into the shell between $`h+dh`$ and $`h`$ from the shell outer surface. The conservation of the number of particles accounting for the sliding gives that number density of particles in the shell at local thickness $`h`$ related to the number density of particles initially at the point $`z_0+z`$ as $$\frac{\sigma }{\sigma _0}n_0(z_0+z)z^2dz=n(h)R^2dh,$$ (25) where $$\sigma _0(t)=\frac{1}{R^2(t)}\underset{0}{\overset{t}{}}R^2\dot{R}\rho (t^{})𝑑t^{},$$ (26) is just $`\sigma `$ without sliding. The number density $`n`$ is given by $`n=P_{sh}(h)/kT_{sh}`$, where $`T_{sh}`$ is the temperature in the shell, which we assume to be equal to $`10^4`$ K, and $`P_{sh}`$ is the pressure at the local thickness $`h`$ in the shell. Here we have taken into account only the thermal pressure and neglected the cosmic ray pressure. This seems to be the reasonable assumption, because, as will be shown later, they have very high drift velocity and may easily escape from the compressed shell along the magnetic field lines deposited into it, since the lines themselves leave the shell. The pressure $`P_{sh}`$ is comparable to the inner pressure $`P_{in}`$, because inside the shell it has to drop from $`P_{in}`$ on the inner surface to $`P_{out}P_{in}`$ outside the shell. Combining all these considerations we obtain that $$H\frac{\sigma }{\sigma _0}\underset{0}{\overset{R}{}}\frac{n_0(z_0+z)kT_{sh}z^2}{P_{in}R^2}𝑑z.$$ (27) Numerical estimates show that during the first several e-foldings of the Rayleigh-Taylor instability, the shell thickness $`H`$ changes very little, less than $`10\%`$. During this time the size of the shell changes less than $`30\%`$, so that the geometrical effect of the stretching the scale of the perturbation mode in the expanding shell is quite moderate. For this reason the modes which were unstable at the very onset of the instability stay unstable during several e-foldings thus developing the nonlinear stage of the instability and disrupting the shell. During the development of the instability the velocity of the shell does not change drastically, it is larger only $`(35)\%`$ than the velocity at the very onset of instability. Thus, we may safely assume that after the Rayleigh-Taylor instability is fully developed in the shell, we get no further acceleration of the shell fragments. We consider the role of the Rayleigh-Taylor instability on the fate of SBs of various luminosities with inner pressure from (7) located at the altitude $`z_0=100`$ pc near the Sun. The results for maximum size of the shell and time when it either stops due to gravity or fragments because of Rayleigh-Taylor instability, are shown on Figure 3. Equation (24) is chosen to be the condition for Rayleigh-Taylor fragmentation of the shell, because it corresponds to the onset of the nonlinear stage of this instability when the shell is being disrupted. The two branches of Figure 3 correspond to the bubbles which were stopped (left branch) and to those which were disrupted (right branch). We see that for a chosen dependence of inner pressure upon time the transition to the Rayleigh-Taylor regime occurs for $`N=104`$ supernovae in the SB. The maximum possible size and the greatest lifetime of the SB are achieved just before this $`N`$ and are equal to $`R_{max}=2616`$ pc and $`t_{max}=52.6`$ Myr. After that size is reached the expansion drops rapidly due to the early onset of the Rayleigh-Taylor instability in the shell, and we see that for SBs with $`N1000`$ shell fragmentation occurs at a very early time $`56`$ Myr. It must be emphasized that at such an early stages the approximation of constant luminosity may not be justified and dynamics might be more complex. We believe that the result of the analysis is not very sensitive to the model assumptions. In Figure 4 we plot the dependence of the blob velocity upon the luminosity of the SB. One can see that it is sufficiently less than $`430\mathrm{km}\mathrm{s}^1`$, the estimated lower bound on the Galactic escape velocity (Leonard $`\&`$ Tremaine 1990), so that we may conclude that shells formed by SBs with number of SNs in them $`N10^3`$ do not give rise to the mass and flux outflow from the Galaxy. Even if we go to a SB with a luminosity an order of magnitude larger ($`N=10^4`$), the velocity at the moment of fragmentation is only $`v_{max}336\mathrm{km}\mathrm{s}^1`$, which is obviously not enough to leave the Galaxy. ## 6 Importance of the shape of the SB top In our treatment of the SB expansion we have considered the shape of the shell near its top to be spherical. This enables us to use a simple fact that in this case the projection of Galactic gravitational acceleration along the shell is just $`g_{}=g\mathrm{sin}\theta `$ which sufficiently simplifies the problem. In reality, of course, the surface density is not uniform on the top but depends upon the angle $`\theta `$. The nonuniformity grows with growing $`\theta `$. The inner pressure will accelerate parts of the shell closer to the top stronger than ones further from the top and it will distort the form of the shell. This distortion in its turn changes $`g_{}`$ which influences the sliding of the matter and thus leads to further changes of the shell’s form. The accurate treatment of the problem requires including this effect self-consistently in our calculations, but we can avoid this by noting that influence of a change of the shell’s shape on the process of sliding can be attributed to the change in gravitational acceleration $`g`$, rather than the projection angle. We carried out calculations identical to those with spherical top but have taken $`g`$ in equation (17) to be $`4`$ times larger than it is in reality. The result was that, as the luminosity of the SB was increased, they expanded faster, due to the more effective sliding of the matter from the top. But this in turn lead to a more rapid onset of Rayleigh-Taylor instability in the shell: it started to develop for SBs with $`N>69`$. The bubble with $`N=69`$ stops at a time $`t_{RT}=45.7`$ Myr and size $`R_{max}=2398`$ pc. The conclusion is obvious: the change of the shell shape might influence the sliding of the matter, but it leads to the development of the Rayleigh-Taylor instability for even smaller luminosities than in the case of the bubbles with the spherical top, making the impossibility of the expulsion of the flux from Galactic disc even more certain. That is why we think a more self-consistent approach to the problem of the shell shape will not change the general result. ## 7 Different pressure law The pressure law (7) which we used in all our calculations was derived actually for the case of the SB expanding in a uniform medium and thus may be not a very good approximation for our purposes especially when the sliding of the matter influences the expansion of the shell and its size cannot be described by equation (5). For that reason we decided to use a different, more realistic pressure dependence to check if it makes a significant difference in our results. Maciejewski $`\&`$ Cox (1999) proposed a simple, explicit, analytical approximation for the kinematics of the blast wave propagating in an exponentially stratified medium: $$\rho _0=\rho _{}e^{z/h},$$ (28) where $`\rho _{}`$ is the density at the explosion site and $`h`$ is the stratification length scale. They considered an explosion in the framework of the Kompaneets approximation (Kompaneets 1960) with inner pressure constant throughout the volume engulfed by the shock. They showed that the form of the shell is very close to an ellipsoid with minor and major axes $`b`$ and $`a`$ related by $$\mathrm{tan}\frac{b}{2h}=\mathrm{sinh}\frac{a}{2h}$$ (29) and the distance from the explosion site to the ellipsoid center, $`s`$, $$\mathrm{tan}\frac{s}{2h}=\mathrm{cosh}\frac{a}{2h}.$$ (30) At the same time, Weaver et al (1977) showed that internal energy of the SB interior is constant fraction of the total energy and for a spherical SB expanding in a homogeneous medium $$E_{in}=\frac{5}{11}L_{SN}t.$$ (31) We assume that this is also valid for the case of SB in a nonhomogeneous density distribution, so that that inner pressure $$P_{in}=\frac{E_{in}}{V}=\frac{5}{11}\frac{L_{SN}t}{\pi ab^2}.$$ (32) We relate the semi-major axis $`a`$ to the distance from the center of explosion to the top of the shell: $`a+s=R`$, by $$\frac{a}{2h}+\mathrm{log}\left(\mathrm{cosh}\frac{a}{2h}\right)=\frac{R}{2h}.$$ (33) Formulae (32),(33), and (29) give us $`P_{in}`$ for a given $`R`$. The approach of Maciejewski $`\&`$ Cox (1999) includes neither galactic gravity nor the slippage from the top but it takes the inhomogeneity of the surrounding medium into account and enables us to test the stability of our results against different model assumptions. The real ISM contains many components distributed with various length scales so we take rather arbitrarily $`h=200`$ pc in our case. The particular choice of $`h`$ turns out not to play a significant role. The results for superbubbles of various luminosities located at $`z_0=200`$ pc near the Sun are shown in Figures 5a and 5b. We see that differences are only quantitative compared to the case of pressure law (7). Strong Rayleigh-Taylor instability starts to dominate the kinematics of the shell when the number of SNs in SB is larger than $`N=18`$. This SB reaches the maximum size $`R_{max}=447`$ pc at a time $`t_{max}=41`$ Myr. Due to the specifics of the the chosen pressure law, SB with the higher luminosity, developing the Rayleigh-Taylor instability, reach somewhat larger size, maximum is $`R_{RT}=630`$ pc. If we take, for example, $`h=100`$ pc, then the shells start to be not stopped by the gravity but disrupted by Rayleigh-Taylor instability even for smaller $`N`$. Thus, again, for powerful SBs, expansion is limited by the shell fragmentation so that all major results of the consideration with the simplified pressure law remain valid. ## 8 Possible importance of the CR pressure in the shell Kulsrud (1999) proposed that sliding of the matter from the top of the SB shell may be inhibited to some extent by cosmic ray (CR) pressure gradient, thus further supporting our conclusion about the impossibility of flux expulsion from the Galaxy. Now, on the basis of the better knowledge of the processes going on in SB shell, we can check this idea in more detail. Let us consider the magnetic flux tube with the cross section constant along the tube in the shell which reached the size $`R`$ and expands with velocity $`u=\dot{R}`$. This assumption is good for the cylindrical shell, with its axis lying in the Galactic plane, but it will be clear further that this assumption is not important. The continuity equation for the CR along the flux tube reads: $$\frac{}{t}\left(n_{CR}R\right)+\frac{}{\theta }\left[n_{CR}\left(v_{}+v_d\right)\right]=0,$$ (34) where $`n_{CR}`$ is the number density of the CR, $`v_{}`$ and $`v_d`$ are the velocities of matter along the shell with respect to the rest system and of CR with respect to the matter correspondingly. The drift of the CR along the magnetic field is determined by the gradient of their number density and by the scattering of the CR by the Alfv$`\stackrel{´}{\mathrm{e}}`$nic turbulence. The scattering of CR by the self-generated Alfv$`\stackrel{´}{\mathrm{e}}`$n waves was first considered by Kulsrud $`\&`$ Pearce (1969) and Wentzel (1969) and they showed that CR cause the Alfv$`\stackrel{´}{\mathrm{e}}`$n waves to grow with a rate $$\mathrm{\Gamma }=C\frac{\pi }{4}\mathrm{\Omega }_0\frac{v_dv_A}{v_A}\frac{n_{CR}(ϵ)}{n},$$ (35) where $`\mathrm{\Omega }_0`$ is the nonrelativistic cyclotron frequency, $`n_{CR}(ϵ)`$ is the number density of CR with energy greater than $`ϵ`$, $`n`$ is the density of the ISM, $`v_A=B/\sqrt{4\pi nm_H}`$ is the Alfv$`\stackrel{´}{\mathrm{e}}`$n speed ($`m_H`$ is the mass of hydrogen atom), and $`C`$ is a constant of the order of unity whose value depends upon the energy spectrum of CR. These waves are also nonlinearly damped by the scattering of the beat waves on the ions: $$\gamma _d=\sqrt{\frac{\pi }{8}}\frac{v_T}{c}\mathrm{\Omega }\left(\frac{\delta B}{B}\right)^2,$$ (36) (Lee $`\&`$ Völk 1973; Kulsrud 1978) with $`\delta B`$ being the magnetic field perturbations due to the Alfv$`\stackrel{´}{\mathrm{e}}`$n waves, $`\mathrm{\Omega }`$ – relativistic Larmor frequency, and $`v_T`$ – thermal velocity of ions in ISM. Scattering of CR occurs at a rate $$\nu =\mathrm{\Omega }\left(\frac{\delta B}{B}\right)^2,$$ (37) (Kulsrud 1995), so that the typical mean free path of the CR is $$\lambda =\frac{c}{\nu }=\frac{c}{\mathrm{\Omega }}\left(\frac{B}{\delta B}\right)^2.$$ (38) On the other hand, the drift velocity is given by $$n_{CR}v_d=c\frac{\lambda }{R}\frac{n_{CR}}{\theta }.$$ (39) Assuming that $`v_dv_A`$, that Alfv$`\stackrel{´}{\mathrm{e}}`$n waves are in equilibrium, $`\mathrm{\Gamma }=\gamma _d`$, and combining equations (35),(36), (38), and (39) we get $`\|v_d\|=\left(\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{cv_Tnv_A}{C\mathrm{\Omega }_0n_{CR}^2}}\left\|{\displaystyle \frac{1}{R}}{\displaystyle \frac{n_{CR}}{\theta }}\right\|\right)^{1/2}`$ $`=\left(A\left\|{\displaystyle \frac{1}{Rn_{CR}^2}}{\displaystyle \frac{n_{CR}}{\theta }}\right\|\right)^{1/2}.`$ (40) If we introduce the characteristic drift velocity $$v_D=\sqrt{\frac{A}{n_{CR}R}},$$ (41) and suppose that $$v_Dv_{},$$ (42) then it is easy to see from equation (34) that when $`v_dv_{}`$ the variation of the CR density can be expressed as $$n_{CR}R=n_{CR_0}R_0\left(1+\left(\frac{v_{}}{v_D}\right)^2\chi _1(\theta )\right),$$ (43) where $`\chi _1(\theta )`$ is some function of order of unity. To see if the condition (42) is fulfilled, we calculated $`v_D`$ given by (41) and (40) at various time moments for the SB with $`N=75`$ SNs, taking into account the compression of the matter and CR in the shell, which we calculate in the manner similar to the calculation of the thickness of the shell, and supposing that the temperature in the shell after cooling is $`T_s=10^4`$ K. Compression of matter in the shell (and, consequently, of the CR) depends upon the temperature of ambient ISM $`T_0`$, so that $`v_D`$ scales as $`v_DT_0^{3/4}`$. In Figure 6 we plot $`v_D`$, defined by (41) for $`T_0=10^4`$ K. One can see from comparison of Figures 6 and 2, that condition (42) is always fulfilled for $`T_0=10^4`$ K. It might be violated if $`T_010^610^7`$ K, that is if the SB expands into the predominantly hot ISM component, but near the top the condition (42) is always fulfilled even for these high temperatures, because $`v_{}\mathrm{sin}^2\theta `$. In the approximation given by (42) the spatial CR density perturbations are small and the first term in equation (34) is negligible. This means that the drift velocity is almost equal to the velocity of matter $`v_{}`$ but opposite to it in the direction, so that CR slide through the matter to maintain constant spatial density and the time variations of their density are only due to the shell expansion. Then the equation of continuity reduces to $$\frac{}{\theta }\left(n_{CR}v_{}\sqrt{\frac{A}{n_{CR}}\frac{n_{CR}}{\theta }}\right)=0.$$ (44) This equation can be integrated with initial condition $`n_{CR}/\theta =0`$ at $`\theta =0`$, where $`v_{}=0`$, to give $$\frac{n_{CR}}{\theta }=n_{CR}v_{}^2\left(\frac{A}{n_{CR}R}\right)^1.$$ (45) Then the ratio of CR pressure term in the equation (10) to the gravitational acceleration along the shell is $`{\displaystyle \frac{p_{CR}}{\rho _sg_{}}}={\displaystyle \frac{p_{CR}}{n_{CR}}}{\displaystyle \frac{n_{CR}}{\theta }}{\displaystyle \frac{1}{Rg\mathrm{sin}\theta \rho _s}}`$ $`{\displaystyle \frac{p_{CR}}{\rho _sgR}}\left({\displaystyle \frac{v_{}}{\mathrm{sin}\theta v_D}}\right)^2\mathrm{sin}\theta 10^4\mathrm{sin}\theta ,`$ (46) for $`T_0=10^4`$ K, or maybe $`0.1\mathrm{sin}\theta `$ at $`T_010^7`$ K, so that the effect the CR pressure gradient has on the sliding of the matter from the shell’s top is negligible for all interesting cases, and, actually, not only at the shell’s top, but also for $`\theta 1`$. Obviously, the assumption that the flux tube has constant cross section and the cylindrical geometry implied here are not important, because the uniformity of the CR pressure is achieved due to the very high characteristic drift velocity $`v_D`$ of the CR compared to the sliding velocity, which perturbs the uniformity, and not due to any peculiarities of the flux tube geometry. ## 9 Summary In this paper we consider the effect of the downsliding of the matter which takes place in the expanding superbubbles for application to the expulsion of the magnetic flux from the Galaxy. This expulsion is an important ingredient of the $`\alpha \text{-}\mathrm{\Omega }`$ dynamo theory. However, it is shown that even the inclusion of the sliding into the calculations of the kinematics of the superbubbles, does not enable matter and frozen in flux to leave the Galaxy in SBs. One must note, that the impossibility of the flux escape from the Galaxy weakens significantly the Parker’s (1971) argument against the primordial magnetic field, because there is actually no mechanism to expel it. Some authors (Korpi et al 1999a, Korpi et al 1999b) have considered dynamics of the superbubbles in the gravitational field of the Galaxy in greater detail and they also find that the matter does not reach terminal velocities larger than the escape velocity from the Galaxy. In their simulations they do observe the development of the SB and its blowout from the disk, but at the height of several kiloparsecs the velocity of the matter is too small for the matter and the field to leave the Galaxy, which agrees with our conclusions. However, they do not comment on its relation to the dynamo. Other authors (e.g. Hanasz & Lesch 1998, Moss et al 1999), more directly concerned with dynamos, do not investigate the dynamics of escape but merely assume that once the magnetic field lines reach the boundary of the disk, they are advected away by some mechanism leading to a vacuum boundary conditions. They do mention magnetic buoyancy as a possible escape mechanism but buoyancy is unlikely to be important outside the disk and clearly plays no role in the SB escape mechanism, especially when the magnetic fields are first amplified from weak seed fields. In this paper an analytical formalism for the consideration of sliding was built on reasonable assumptions, which enabled to include the back-reaction of lowering the density of the matter on the top of the SB, on the expansion of the shell in the radial direction. All the SB, depending on their luminosity and conditions in the surrounding ISM were demonstrated to fall into two classes: low luminosity SBs which are stopped by the Galactic gravitational field and fall back, and powerful SBs, which are possibly able to reach low density regions at high altitude. Even without sliding these powerful SBs might be able to expel matter from the Galaxy. However, before this happens the shells of these high luminosity SBs fragment into separate blobs due to the Rayleigh-Taylor instability, which develop when the shells start accelerating at the high altitudes, where the density of the matter and the outer pressure are very small. These separate fragments of the shells continue to move ballistically with velocities too small to leave the Galaxy. It was shown that our conclusion about the impossibility of flux to escape in SB explosions does not depend essentially upon the details of the inner pressure behavior in the SB or the shape of the SB top. The inclusion of the CR pressure gradient in the shell also does not influence the downsliding because of the high uniformity of CR density in the shell caused by the very large diffusion velocity of CR in the shell. It should be remarked that throughout the paper we have assumed that magnetic field lines are tied to the gas. If the ionization is the shell is low, then ambipolar diffusion (under the influence of the CR pressure gradient perpendicular to the shell) could allow some small slippage of the field lines through the neutral component and, thus, some small amount of escape. We do not discuss this possibility in this paper. All this lead us to the conclusion that if there does exist some mechanism responsible for the flux expulsions from the Galaxy needed by the $`\alpha \text{-}\mathrm{\Omega }`$ dynamo this mechanism is not superbubbles. The question of whether such a mechanism exist remain an open one. ## 10 Acknowledgements One of the authors (RRR) would like to acknowledge the financial support of this work by the Princeton University Science Fellowship.
warning/0004/quant-ph0004005.html	ar5iv	text	# References Quantum mechanics with time-dependent parameters G. Sardanashvily<sup>1</sup><sup>1</sup>1Electronic mail: sard@grav.phys.msu.su; sard@campus.unicam.it Department of Theoretical Physics, Moscow State University, 117234 Moscow, Russia Abstract. Composite bundles $`Q\mathrm{\Sigma }𝐑`$, where $`\mathrm{\Sigma }𝐑`$ is the parameter bundle, provide the adequate mathematical description of classical mechanics with time-dependent parameters. We show that the Berry’s phase phenomenon is described in terms of connections on composite Hilbert space bundles. I. Smooth fiber bundles $`Q𝐑`$ over a time axis $`𝐑`$ provide the adequate formulation of classical time-dependent mechanics treated as a particular field theory . Let us consider a mechanical system depending on time-dependent parameters. These parameters can be seen as sections of some smooth fiber bundle $`\mathrm{\Sigma }𝐑`$. Then the configuration space of a mechanical system with time-dependent parameters can be seen as the composite fiber bundle $$Q\mathrm{\Sigma }𝐑.$$ (1) In classical mechanics $`Q\mathrm{\Sigma }`$ is a smooth finite-dimensional fiber bundle. In quantum mechanics $`Q\mathrm{\Sigma }`$ is a $`C^{}`$-algebra fiber bundle or a Hilbert space fiber bundle . The following two facts make the composite fiber bundle (1) useful for our purpose. (i) Given a section $`h`$ of a parameter bundle $`\mathrm{\Sigma }𝐑`$, the pull-back bundle $`h^{}Q`$ over $`𝐑`$ describes a mechanical system under the fixed parameter functions $`h(t)`$. (ii) Given a connection $`A_\mathrm{\Sigma }`$ on the fiber bundle $`Q\mathrm{\Sigma }`$, the pull-back connection $`h^{}A_\mathrm{\Sigma }`$ on the pull-back bundle $`h^{}Q𝐑`$ depends in a certain way on the parameter functions $`h(t)`$, and characterizes the dynamics of a mechanical system with time-dependent parameters. This work is devoted to quantum mechanics with classical parameters where connections on composite Hilbert space bundles play the role of Berry connections. II. Recall that by a smooth composite bundle is meant the composition of fiber bundles $$Y\mathrm{\Sigma }X,$$ (2) where $`\pi _{Y\mathrm{\Sigma }}:Y\mathrm{\Sigma }`$ and $`\pi _{\mathrm{\Sigma }X}:\mathrm{\Sigma }X`$ are smooth fiber bundles . It is provided with an atlas of fibered coordinates $`(x^\lambda ,\sigma ^m,y^i)`$, where $`(x^\mu ,\sigma ^m)`$ are fibered coordinates on the fiber bundle $`\mathrm{\Sigma }X`$ and the transition functions $`\sigma ^m\sigma ^m(x^\lambda ,\sigma ^k)`$ are independent of the fiber coordinates $`y^i`$. Proposition 1: Given a composite fiber bundle (2), let $`h`$ be a global section of the fiber bundle $`\mathrm{\Sigma }X`$. Then the restriction $$Y_h=h^{}Y$$ (3) of the fiber bundle $`Y\mathrm{\Sigma }`$ to $`h(X)\mathrm{\Sigma }`$ is a subbundle $`i_h:Y_hY`$ of the fiber bundle $`YX`$. Let us consider a connection $$A_\mathrm{\Sigma }=dx^\lambda (_\lambda +A_\lambda ^i_i)+d\sigma ^m(_m+A_m^i_i):YJ_\mathrm{\Sigma }^1Y$$ (4) on the fiber bundle $`Y\mathrm{\Sigma }`$. Given a section $`h`$ the fiber bundle $`\mathrm{\Sigma }X`$, the connection $`A_\mathrm{\Sigma }`$ (4) induces the pull-back connection $$A_h=i_h^{}A_\mathrm{\Sigma }=dx^\lambda [_\lambda +((A_m^ih)_\lambda h^m+(Ah)_\lambda ^i)_i]$$ (5) on the pull-back bundle $`Y_h`$ (3). Note that, in quantum theory, one follows the notion of a connection phrased in algebraic terms as a connection on modules in comparison with the pure geometric one in classical theory. Here, we restrict our consideration to connecions on modules over the ring $`C^{\mathrm{}}(X)`$ of smooth real functions on a manifold $`X`$ . Definition 2: A connection on a $`C^{\mathrm{}}(X)`$-module $`𝒮`$ assigns to each vector field $`\tau `$ on a manifold $`X`$ an $`𝒮`$-valued first order differential operator $`_\tau \mathrm{Diff}_1(𝒮,𝒮)`$ on $`𝒮`$ which obeys the Leibniz rule $$_\tau (fs)=(\tau df)s+f_\tau s,fC^{\mathrm{}}(X),s𝒮.$$ (6) If $`𝒮`$ is a module of global sections of a smooth vector bundle $`YX`$ over a manifold $`X`$, Definition 1 is equivalent to the familiar geometric definition of a connection on $`YX`$. III. Let us consider a quantum mechanical systems depending on a finite number of real classical parameters given by sections of a smooth parameter bundle $`\mathrm{\Sigma }𝐑`$. For the sake of simplicity, we fix a trivialization $`\mathrm{\Sigma }=𝐑\times Z`$, coordinated by $`(t,\sigma ^m)`$. Although it may happen that the parameter bundle $`\mathrm{\Sigma }𝐑`$ has no preferable trivialization, e.g., if one of parameters is a velocity of a reference frame. Recall that, in the framework of algebraic quantum theory, a quantum system is characterized by a $`C^{}`$-algebra $`A`$ and a positive (hence, continuous) form $`\varphi `$ on $`A`$ which defines the representation $`\pi _\varphi `$ of $`A`$ in a Hilbert space $`E_\varphi `$ with a cyclic vector $`\xi _\varphi `$ such that $`\varphi (a)=\pi _\varphi (a)\xi _\varphi \|\xi _\varphi ,aA.`$ One says that $`\varphi (a)`$ is a mean value of the operator $`a`$ in the state $`\xi _\varphi `$. It should be emphasized that, in quantum mechanics, a time also plays the role of a classical parameter. Indeed, all relations between operators in quantum mechanics are simultaneous, while a computation of a mean value of an operator in a quantum state does not imply an integration over a time. It follows that, at each moment, we have a quantum system, but these quantum systems are different at different instants. Though they may be isomorphic to each other. This characteristic is extended to other classical parameters. Namely, we assign a $`C^{}`$-algebra $`A_\sigma `$ to each point $`\sigma \mathrm{\Sigma }`$ of the parameter bundle $`\mathrm{\Sigma }`$, and treat $`A_\sigma `$ as a quantum system under fixed values $`(t,\sigma ^m)`$ of the parameters. Remark 1: Let us emphasize that one should distinguish classical parameters from coordinates which a wave function can depend on. Let $`\{A_q\}`$ be a set of $`C^{}`$-algebras parameterized by points of a locally compact topological space $`Q`$. Let all $`C^{}`$-algebras $`A_q`$ are isomorphic to each other and to some $`C^{}`$-algebra $`A`$. We consider a locally trivial topological fiber bundle $`PQ`$ whose typical fiber is the $`C^{}`$-algebra $`A`$, i.e., transition functions of this fiber bundle provide automorphisms of $`A`$. The set $`P(Q)`$ of continuous sections of this fiber bundle is a \*-algebra with respect to fiberwise operations. Let us consider a subalgebra $`A(Q)P(Q)`$ which consists of sections $`\alpha `$ of $`PQ`$ such that the real function $`\alpha (q)`$ vanishes at infinity of $`Q`$. For $`\alpha A(Q)`$, put $`\alpha =\underset{qQ}{sup}\alpha (q)<\mathrm{}.`$ With this norm, $`A(Q)`$ is a $`C^{}`$-algebra . One can consider a quantum system characterized by this $`C^{}`$-algebra. In this case, elements of the set $`Q`$ are not classical parameters as follows. Given an element $`qQ`$, the assignment $$A(Q)\alpha \alpha (q)A$$ (7) is a $`C^{}`$-algebra epimorphism. Let $`\pi `$ be a representation of $`A`$. Then the assignment (7) yields a representation $`\rho (\pi ,q)`$ of the $`C^{}`$-algebra $`A(Q)`$. If $`\pi `$ is an irreducible representation of the $`C^{}`$-algebra $`A`$, then $`\rho (\pi ,q)`$ is an irreducible representation of $`A(Q)`$. Moreover, the irreducible representations $`\rho (\pi ,q)`$ and $`\rho (\pi ,q^{})`$ of $`A(Q)`$ are not equivalent . Therefore there is one-to-one correspondence (but not a homeomorphism) between the spectrum $`\widehat{A(Q)}`$ of the $`C^{}`$-algebra $`A(Q)`$ and the product $`Q\times \widehat{A}`$ of $`Q`$ and the spectrum $`\widehat{A}`$ of the $`C^{}`$-algebra $`A`$. It follows that one can find representations of the $`C^{}`$-algebra $`A(Q)`$ among direct integrals of representations of $`A`$ with respect to some measure on $`Q`$. Let $`\mu `$ be a positive measure of total mass 1 on the locally compact space $`Q`$ , and let $`\varphi `$ be a positive form on $`A`$. Then the function $`q\varphi (\alpha (q))`$, $`\alpha A(Q)`$, is a $`\mu `$-measurable, while the integral $`\varphi (\alpha )={\displaystyle \varphi (\alpha (q))\mu (q)}`$ provides a positive form on the $`C^{}`$-algebra $`A(Q)`$. Roughly speaking, a computation of a mean value of an operator $`\alpha A(Q)`$ implies an integration with respect to some measure on $`Q`$ in general. This is not the case of quantum systems depending on classical parameters $`qQ`$. We simplify our consideration in order to single out the manifested Berry’s phase phenomenon. Let us assume that all algebras $`C^{}`$-algebras $`A_\sigma `$, $`\sigma \mathrm{\Sigma }`$, are isomorphic to the von Neumann algebra $`B(E)`$ of bounded operators in some Hilbert space $`E`$, and consider a locally trivial Hilbert space bundle $`\mathrm{\Pi }\mathrm{\Sigma }`$ with the typical fiber $`E`$ and smooth transition functions . Smooth sections of $`\mathrm{\Pi }\mathrm{\Sigma }`$ constitute a module $`\mathrm{\Pi }(\mathrm{\Sigma })`$ over the ring $`C^{\mathrm{}}(\mathrm{\Sigma })`$ of real functions on $`\mathrm{\Sigma }`$. In accordance with Definition 1, a connection $`\stackrel{~}{}`$ on $`\mathrm{\Pi }(\mathrm{\Sigma })`$ assigns to each vector field $`\tau `$ on $`\mathrm{\Sigma }`$ a first order differential operator $$\stackrel{~}{}_\tau \mathrm{Diff}_1(\mathrm{\Pi }(\mathrm{\Sigma }),\mathrm{\Pi }(\mathrm{\Sigma }))$$ (8) which obeys the Leibniz rule $`\stackrel{~}{}_\tau (fs)=(\tau df)s+f\stackrel{~}{}_\tau s,s\mathrm{\Pi }(\mathrm{\Sigma }),fC^{\mathrm{}}(\mathrm{\Sigma }).`$ Let $`\tau `$ be a vector field on $`\mathrm{\Sigma }`$ such that $`dt\tau =1`$. Given a trivialization chart of the Hilbert space bundle $`\mathrm{\Pi }\mathrm{\Sigma }`$, the operator $`\stackrel{~}{}_\tau `$ (8) reads $$\stackrel{~}{}_\tau (s)=(_ti(t,\sigma ^i))s+\tau ^m(_mi\widehat{A}_m(t,\sigma ^i))s,$$ (9) where $`(t,\sigma ^i)`$, $`\widehat{A}_m(t,\sigma ^i)`$ for each $`\sigma \mathrm{\Sigma }`$ are bounded self-adjoint operators in the Hilbert space $`E`$. Let us consider the composite fiber bundle $`\mathrm{\Pi }\mathrm{\Sigma }𝐑`$. Similarly to the case of smooth composite fiber bundles (see Proposition 1), every section $`h(t)`$ of the parameter bundle $`\mathrm{\Sigma }𝐑`$ defines the subbundle $`\mathrm{\Pi }_h=h^{}\mathrm{\Pi }𝐑`$ of the composite fiber bundle $`\mathrm{\Pi }𝐑`$ whose typical fiber is the Hilbert space $`E`$. Accordingly, the connection $`\stackrel{~}{}`$ (9) on the $`C^{\mathrm{}}(\mathrm{\Sigma })`$-module $`\mathrm{\Pi }(\mathrm{\Sigma })`$ defines the pull-back connection $$_h(\psi )=[_ti(\widehat{A}_m(t,h^i(t))_th^m+(t,h^i(t))]\psi $$ (10) on the $`C^{\mathrm{}}(𝐑)`$-module $`\mathrm{\Pi }_h(𝐑)`$ of sections $`\psi `$ of the Hilbert space bundle $`\mathrm{\Pi }_h𝐑`$. As in the case of smooth fiber bundles, we say that a section $`\psi `$ of the fiber bundle $`\mathrm{\Pi }_h𝐑`$ is an integral section of the connection (10) if $$_h(\psi )=[_ti(\widehat{A}_m(t,h^i(t))_th^m+(t,h^i(t))]\psi =0.$$ (11) One can think of the equation (11) as being the Shrödinger equation for a quantum system depending on the parameter function $`h(t)`$. Its solutions take the form $$G_t=T\mathrm{exp}\left[i\underset{0}{\overset{t}{}}(\widehat{A}_m_t^{}h^m+)dt^{}\right],$$ (12) where $`G_t`$ is the time-ordered exponent. The term $`i\widehat{A}_m(t,h^i(t))_th^m`$ in the Shrödinger equation (11) is responsible for the Berry’s phase phenomenon, while $``$ is treated as an ordinary Hamiltonian of a quantum system. To show the Berry’s phase phenomenon clearly, we simplify again the system under consideration. Given a trivialization of the fiber bundle $`\mathrm{\Pi }𝐑`$ and the above mentioned trivialization $`\mathrm{\Sigma }=𝐑\times Z`$ of the parameter bundle $`\mathrm{\Sigma }`$, let us suppose that the components $`\widehat{A}_m`$ of the connection $`\stackrel{~}{}`$ (9) are independent of $`t`$ and that the operators $`(\sigma )`$ commute with the operators $`\widehat{A}_m(\sigma )`$ at all points of the curve $`h(t)\mathrm{\Sigma }`$. Then the operator $`G_t`$ (12) takes the form $$G_t=T\mathrm{exp}\left[i\underset{h([0,t])}{}\widehat{A}_m(\sigma ^i)d\sigma ^m\right]T\mathrm{exp}\left[i\underset{0}{\overset{t}{}}(t^{})dt^{}\right].$$ (13) One can think of the first factor in the right-hand side of the expression (13) as being the operator of a parallel transport along the curve $`h([0,t])Z`$ with respect to the pull-back connection $$=i^{}\stackrel{~}{}=_mi\widehat{A}_m(t,\sigma ^i)$$ (14) on the fiber bundle $`\mathrm{\Pi }Z`$, defined by the imbedding $`i:Z\{0\}\times Z\mathrm{\Sigma }.`$ Note that, since operators $`\widehat{A}_m`$ are independent of time, one can utilize any imbedding of $`Z`$ to $`\{t\}\times Z`$. Moreover, the connection $``$ (14), called the Berry connection, can be seen as a connection on some principal fiber bundle $`PZ`$ for the group $`U(E)`$ of unitary operators in the Hilbert space $`E`$. Let the curve $`h([0,t])`$ be closed, while the holonomy group of the connection $``$ at the point $`h(t)=h(0)`$ is not trivial. Then the unitary operator $$T\mathrm{exp}\left[i\underset{h([0,t])}{}\widehat{A}_m(\sigma ^i)d\sigma ^m\right]$$ (15) is not the identity. For example, if $$i\widehat{A}_m(\sigma ^i)=iA_m(\sigma ^i)\mathrm{Id}_E$$ (16) is a $`U(1)`$-principal connection on $`Z`$, then the operator (15) is the well-known Berry phase factor $`\mathrm{exp}\left[i{\displaystyle \underset{h([0,t])}{}}A_m(\sigma ^i)d\sigma ^m\right].`$ If (16) is a curvature-free connection, Berry’s phase is exactly the Aharonov–Bohm effect on the parameter space $`Z`$. The following variant of the Berry’s phase phenomenon leads us to a principal bundle for familiar finite-dimensional Lie groups. Let $`E`$ be a separable Hilbert space which is the Hilbert sum of $`n`$-dimensional eigenspaces of the Hamiltonian $`(\sigma )`$, i.e., $`E={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}E_k,E_k=P_k(E),`$ where $`P_k`$ are the projection operators, i.e., $`H(\sigma )P_k=\lambda _k(\sigma )P_k`$ (in the spirit of the adiabatic hypothesis). Let the operators $`\widehat{A}_m(z)`$ be time-independent and preserve the eigenspaces $`E_k`$ of the Hamiltonian $``$, i.e., $$\widehat{A}_m(z)=\underset{k}{}\widehat{A}_m^k(z)P_k,$$ (17) where $`\widehat{A}_m^k(z)`$, $`zZ`$, are self-adjoint operators in $`E_k`$. It follows that $`\widehat{A}_m(\sigma )`$ commute with $`(\sigma )`$ at all points of the parameter bundle $`\mathrm{\Sigma }𝐑`$. Then, restricted to each subspace $`E_k`$, the parallel transport operator (15) is a unitary operator in $`E_k`$. In this case, the Berry connection (14) on the $`U(E)`$-principal bundle $`PZ`$ can be seen as a composite connection on the composite bundle $`PP/U(n)Z,`$ which is defined by some principal connection on the $`U(n)`$-principal bundle $`PP/U(n)`$ and the trivial connection on the fiber bundle $`P/U(n)Z`$. The typical fiber of $`P/U(n)Z`$ is exactly the classifying space $`B(U(n))`$ for $`U(n)`$-principal bundles. Moreover, one can consider the parallel transport along a curve in the bundle $`P/U(n)`$. In this case, a state vector $`\psi (t)`$ acquires a geometric phase factor in addition to the dynamical phase factor. In particular, if $`\mathrm{\Sigma }=𝐑`$ (i.e., classical parameters are absent and Berry’s phase has only the geometric origin) we come to the case of a Berry connection on the $`U(n)`$-principal bundle on the classifying space $`B(U(n))`$ . If $`n=1`$, this is the variant of Berry’s geometric phase of Ref. .
warning/0004/quant-ph0004040.html	ar5iv	text	# 1 Introduction ## 1 Introduction A quantum system, prepared in a state that does not belong to an eigenvalue of its total Hamiltonian, starts to evolve quadratically in time . This behavior leads to the so-called quantum Zeno phenomenon: by performing frequent measurements on the system, in order to check whether it is still in its initial state, one can “slow down” its temporal evolution (hindering transitions to states different from the initial one) . This curious feature of the quantal evolution has recently attracted much attention in the physics community. This is mainly due to a nice idea put forward by Cook , who proposed to check this effect on a two-level system, and to a related experimental test , that motivated an interesting discussion . In turn, this has led to new proposals and experiments . However, it should be emphasized that these studies do not deal with bona fide unstable systems, following (approximately) exponential laws, as in the original proposals . The presence of a non-exponential decay at short times has been detected only recently . The aim of the present paper is to investigate an interesting (and often overlooked) feature of what we might call a quantum Zeno dynamics. We shall see that a series of “measurements” (von Neumann’s projections ) does not necessarily hinder the evolution of the quantum system. On the contrary, the system can evolve away from its initial state, provided it remains in the subspace defined by the “measurement” itself. This interesting feature is readily understandable in terms of a theorem proved by Misra and Sudarshan (MS) , but it seems to us that it is worth clarifying it further by analyzing some interesting examples. ## 2 Misra and Sudarshan’s theorem Consider a quantum system Q, whose states belong to the Hilbert space $``$ and whose evolution is described by the unitary operator $`U(t)=\mathrm{exp}(iHt)`$, where $`H`$ is a time-independent semi-bounded Hamiltonian. Let $`E`$ be a projection operator that does not commute with the Hamiltonian, $`[E,H]0`$, and $`EE=_E`$ the subspace spanned by its eigenstates. The initial density matrix $`\rho _0`$ of system Q is taken to belong to $`_E`$. If Q is let to follow its “undisturbed” evolution, under the action of the Hamiltonian $`H`$ (i.e., no measurements are performed in order to get information about its quantum state), the final state at time $`T`$ reads $$\rho (T)=U(T)\rho _0U^{}(T)$$ ($`2.1`$) and the probability that the system is still in $`_E`$ at time $`T`$ is $$P(T)=\text{Tr}\left[U(T)\rho _0U^{}(T)E\right].$$ ($`2.2`$) We call this a “survival probability:” it is in general smaller than 1, since the Hamiltonian $`H`$ induces transitions out of $`_E`$. We shall say that the quantum system has “survived” if it is found to be in $`_E`$ by means of a suitable measurement process . We stress that we do not distinguish between one- and many-dimensional projections: in the examples to be considered in this note, $`E`$ will be infinite-dimensional. Assume that we perform a measurement at time $`t`$, in order to check whether Q has survived. Such a measurement is formally represented by the projection operator $`E`$. By definition, $$\rho _0=E\rho _0E,\text{Tr}[\rho _0E]=1.$$ ($`2.3`$) After the measurement, the state of Q changes into $$\rho _0\rho (t)=EU(t)\rho _0U^{}(t)E,$$ ($`2.4`$) with probability $`P(t)`$ $`=`$ $`\text{Tr}\left[U(t)\rho _0U^{}(t)E\right]=\text{Tr}\left[EU(t)E\rho _0EU^{}(t)E\right]`$ ($`2.5`$) $`=`$ $`\text{Tr}\left[V(t)\rho _0V^{}(t)\right].(V(t)EU(t)E)`$ This is the probability that the system has “survived” in $`_E`$. There is, of course, a probability $`1P`$ that the system has not survived (i.e., it has made a transition outside $`_E`$) and its state has changed into $`\rho ^{}(t)=(1E)U(t)\rho _0U^{}(t)(1E)`$. The states $`\rho `$ and $`\rho ^{}`$ together make up a block diagonal matrix: The initial density matrix is reduced to a mixture and any possibility of interference between “survived” and “not survived” states is destroyed (complete decoherence). We shall concentrate henceforth our attention on the measurement outcome ($`2.4`$)-($`2.5`$). We observe that the evolution just described is time-translation invariant and the dynamics is not reversible (not only not time-reversal invariant). The above is the Copenhagen interpretation: the measurement is considered to be instantaneous. The “quantum Zeno paradox” is the following. We prepare Q in the initial state $`\rho _0`$ at time 0 and perform a series of $`E`$-observations at times $`t_j=jT/N(j=1,\mathrm{},N)`$. The state of Q after the above-mentioned $`N`$ measurements reads $$\rho ^{(N)}(T)=V_N(T)\rho _0V_N^{}(T),V_N(T)[EU(T/N)E]^N$$ ($`2.6`$) and the probability to find the system in $`_E`$ (“survival probability”) is given by $$P^{(N)}(T)=\text{Tr}\left[V_N(T)\rho _0V_N^{}(T)\right].$$ ($`2.7`$) Equations ($`2.6`$)-($`2.7`$) display the “quantum Zeno effect:” repeated observations in succession modify the dynamics of the quantum system; under general conditions, if $`N`$ is sufficiently large, all transitions outside $`_E`$ are inhibited. Notice again that the dynamics ($`2.6`$)-($`2.7`$) is not reversible. In order to consider the $`N\mathrm{}`$ limit (“continuous observation”), one needs some mathematical requirements: assume that the limit $$𝒱(T)\underset{N\mathrm{}}{lim}V_N(T)$$ ($`2.8`$) exists in the strong sense. The final state of Q is then $$\rho (T)=𝒱(T)\rho _0𝒱^{}(T)$$ ($`2.9`$) and the probability to find the system in $`_E`$ is $$𝒫(T)\underset{N\mathrm{}}{lim}P^{(N)}(T)=\text{Tr}\left[𝒱(T)\rho _0𝒱^{}(T)\right].$$ ($`2.10`$) One should carefully notice that nothing is said about the final state $`\rho (T)`$, which depends on the characteristics of the model investigated and on the very measurement performed (i.e. on the projection operator $`E`$, by means of which $`V_N`$ is defined). By assuming the strong continuity of $`𝒱(t)`$ $$\underset{t0^+}{lim}𝒱(t)=E,$$ ($`2.11`$) one can prove that under general conditions the operators $$𝒱(T)\text{exist for all real }T\text{ and form a semigroup.}$$ ($`2.12`$) Moreover, by time-reversal invariance $$𝒱^{}(T)=𝒱(T),$$ ($`2.13`$) so that $`𝒱^{}(T)𝒱(T)=E`$. This implies, by ($`2.3`$), that $$𝒫(T)=\text{Tr}\left[\rho _0𝒱^{}(T)𝒱(T)\right]=\text{Tr}\left[\rho _0E\right]=1.$$ ($`2.14`$) If the particle is “continuously” observed, in order to check whether it has survived inside $`_E`$, it will never make a transition to $`_E^{}`$. This was named “quantum Zeno paradox” . The expression “quantum Zeno effect” seems more appropriate, nowadays. Two important remarks are now in order: first, it is not clear whether the dynamics in the $`N\mathrm{}`$ limit is time reversible. Although one ends up, in general, with a semigroup, there are concrete elements of reversibility in the above equations. Second, the theorem just summarized does not state that the system remains in its initial state, after the series of very frequent measurements. Rather, the system is left in the subspace $`_E`$, instead of evolving “naturally” in the total Hilbert space $``$. This subtle point, implied by Eqs. ($`2.9`$)-($`2.14`$), is not duely stressed in the literature . Incidentally, we stress that there is a conceptual gap between Eqs. ($`2.7`$) and ($`2.10`$): to perform an experiment with $`N`$ finite is only a practical problem, from the physical point of view. On the other hand, the $`N\mathrm{}`$ case is physically unattainable, and is rather to be regarded as a mathematical limit (although a very interesting one). In this paper, we shall not be concerned with this problem (investigated in ; see also , where an interesting perspective is advocated) and shall consider the $`N\mathrm{}`$ limit for simplicity. This will make the analysis more transparent. ## 3 Evolution in the “Zeno” subspace We start off by looking at some explicit examples. Consider a free particle of mass $`m`$ on the real line. The Hamiltonian and the corresponding evolution operator are $$H=\frac{p^2}{2m},U(t)=\mathrm{exp}(itH).$$ ($`3.1`$) Observe that $`H`$ is a positive-definite self-adjoint operator in $`L^2()`$ and $`U(t)`$ is unitary. We shall study the quantum Zeno effect when the system undergoes a measurement defined by the projector $$E_A=𝑑x\chi _A(x)\|xx\|,$$ ($`3.2`$) where $`\chi _A`$ is the characteristic function $$\chi _A(x)=\{\begin{array}{c}1\text{for }xA\hfill \\ 0\text{otherwise}\hfill \end{array}$$ ($`3.3`$) and $`A`$ an interval of $``$. In a few words, we check whether a particle, initially prepared in a state with support in $`A`$ and free to move on the real line, is still found in $`A`$ at a later time $`T`$. Our objective is to understand how the system evolves in the “Zeno” subspace $`_{E_A}=E_AE_A`$. We call this a “quantum Zeno dynamics with a nonholonomic constraint.” We shall work with the Euclidean Feynman integral. Let the particle be initially ($`t=0`$) at position $`yE_A`$. The propagator at time $`t=T/N`$, when a measurement is carried out, reads $$G(x,t;y)x\|E_AU(t)\|y=\chi _A(x)x\|U(t)\|y.$$ ($`3.4`$) For imaginary time $`t=i\tau `$, we get the Green function of the heat equation $`x\|U(i\tau )\|y`$ $`=`$ $`x\|\mathrm{exp}(\tau H)\|y={\displaystyle 𝑑px\|pe^{\tau p^2/2m}p\|y}`$ ($`3.5`$) $`=`$ $`{\displaystyle \frac{dp}{2\pi }e^{\tau p^2/2m+ip(xy)}}=\sqrt{{\displaystyle \frac{m}{2\pi \tau }}}\mathrm{exp}\left[{\displaystyle \frac{m(xy)^2}{2\tau }}\right],`$ so that the Euclidean propagator for a single “step” reads $$W(x,\tau ;y)G(x,i\tau ;y)=\chi _A(x)\sqrt{\frac{m}{2\pi \tau }}\mathrm{exp}\left[\frac{m(xy)^2}{2\tau }\right].$$ ($`3.6`$) The evolution operator after $`N`$ measurements, see ($`2.6`$), can be written as $$V_N(T)[E_AU(T/N)]^NE_A$$ ($`3.7`$) and the resulting propagator is $$G_N(x_\mathrm{f},T;x_\mathrm{i})=x_\mathrm{f}\|V_N(T)\|x_\mathrm{i}.$$ ($`3.8`$) For imaginary $`𝒯=iT`$ this becomes $`W_N(x_\mathrm{f},𝒯;x_\mathrm{i})`$ $``$ $`G_N(x_\mathrm{f},i𝒯;x_\mathrm{i})`$ ($`3.9`$) $`=`$ $`{\displaystyle 𝑑x_1\mathrm{}𝑑x_{N1}W(x_\mathrm{f},\tau ;x_{N1})\mathrm{}W(x_1,\tau ;x_\mathrm{i})\chi _A(x_i)},`$ whose relation with Wiener integration is manifest. Notice that if we could drop the characteristic function $`\chi _A`$ in the propagator ($`3.6`$), then ($`3.9`$) would be a sequence of nested Gaussian integrals, that could be evaluated exactly for every $`N`$ by applying Feynman’s recipe . In ($`3.9`$) the characteristic functions restrict at every step the set of possible paths, modifying the structure of the functional integral . Let us therefore try to reduce the integral ($`3.9`$) to a Gaussian form. To this end we apply a trick that is often used when one endeavours to relate probability and potential theory . We first rewrite the characteristic function in terms of a potential, which is infinite outside $`A`$ , so that the Brownian paths of the Wiener process ($`3.9`$) can never leak out of $`A`$: $$\chi _A(x)=\mathrm{exp}\left(\tau V_A(x)\right),\text{with}V_A(x)=\{\begin{array}{c}0\text{for }xA\hfill \\ +\mathrm{}\text{otherwise}\hfill \end{array},$$ ($`3.10`$) Hence, by using ($`3.10`$), the Euclidean one-step propagator ($`3.6`$) becomes $`W(x,\tau ;y)=\sqrt{{\displaystyle \frac{m}{2\pi \tau }}}\mathrm{exp}\left[{\displaystyle \frac{m(xy)^2}{2\tau }}\tau V_A(x)\right]=x\|e^{\tau V_A}e^{\tau H}\|y`$ ($`3.11`$) and returning to real time $$G(x,t;y)=W(x,it;y)=x\|e^{itV_A}e^{itH}\|y.$$ ($`3.12`$) Consider now the limit of continuous observation $`N\mathrm{}`$. The limiting propagator reads $$𝒢(x_\mathrm{f},T;x_\mathrm{i})=\underset{N\mathrm{}}{lim}G_N(x_\mathrm{f},T;x_\mathrm{i})=\underset{N\mathrm{}}{lim}x_\mathrm{f}\|\left(e^{iTV_A/N}e^{iTH/N}\right)^NE_A\|x_\mathrm{i},$$ ($`3.13`$) which, by using the Trotter product formula, yields $$𝒢(x_\mathrm{f},T;x_\mathrm{i})=x_\mathrm{f}\|e^{iT(H+V_A)}E_A\|x_\mathrm{i}=x_\mathrm{f}\|𝒱(T)\|x_\mathrm{i},$$ ($`3.14`$) where the evolution operator is $`𝒱(T)=\mathrm{exp}(iTH_\mathrm{Z})E_A,`$ ($`3.15`$) $`\text{with}H_\mathrm{Z}{\displaystyle \frac{p^2}{2m}}+V_A(x).`$ ($`3.16`$) The above formula is of general validity: the dynamics within the Zeno subspace $`_{E_A}`$ is governed by the operators ($`3.15`$)-($`3.16`$). It is worth stressing that the previous calculation only makes use of the properties of the kinetic energy operator $`p^2`$: we have not considered the momentum operator $`p`$. It goes without saying that $`p`$ can be symmetric, maximally symmetric or self-adjoint, according to the structure of $`A`$ and the boundary conditions. This will be thoroughly discussed in the following. However, we emphasize that any requirement on $`p`$ would be a physical requirement: the mathematical properties of the “Zeno” evolution only involve the Hamiltonian (which is defined in terms of the kinetic energy). Before we proceed further, let us look at two particular cases: $`A_1`$ $`=`$ $`[0,1],`$ ($`3.17`$) $`A_2`$ $`=`$ $`[0,+\mathrm{}).`$ ($`3.18`$) In the first case, the free Hamiltonian $$H_\mathrm{Z}^0=\frac{^2}{2m}\left(\frac{d}{dx}\right)$$ ($`3.19`$) is a self-adjoint operator on the space $$D_{[0,1]}(H_\mathrm{Z}^0)=\left\{\varphi \text{AC}^2[0,1]\right\|\varphi (0)=\varphi (1)=0\},$$ ($`3.20`$) where AC$`{}_{}{}^{2}[S]`$ is the set of functions in $`L^2[S]`$ whose weak derivatives are in AC$`[S]`$. (AC$`[S]`$ is the set of absolutely continuous functions whose weak derivatives are in $`L^2[S]`$.) Notice that these are the “correct” boundary condition for the potential ($`3.10`$). For this reason, the evolution operators $`𝒱(T)`$ in ($`3.15`$) form a one-parameter group. We notice, incidentally, that MS’s mathematical hypotheses ($`2.8`$) and ($`2.11`$) are satisfied and acquire in this example an appealing physical meaning. We also stress that the theorem ($`2.12`$) appears in this case too restrictive: indeed the operators $`𝒱(T)`$ form a group and not simply a semigroup. One might say that in the example considered, the quantum Zeno effect (engendered by the projection operators) automatically yields the “natural” dynamics in the Zeno subspace, with the correct boundary conditions for the “new” Hamiltonian $`H_\mathrm{Z}`$. This is an interesting observation in itself. We also notice that in this example the momentum operator $`i`$ is symmetric, but not self-adjoint: its deficiency indices in ($`3.20`$) are (1,1). Therefore, a self-adjoint extension of $`i`$ is possible. It is important to stress that the Hamiltonian $`H_\mathrm{Z}`$ is self-adjoint because it involves only $`^2`$ \[which is self-adjoint in ($`3.20`$)\]. There is here an interesting classical analogy: when a classical particle elastically bounces between two rigid walls, any trajectory is characterized by a definite value of energy ($`p^2/2m`$), although momentum changes periodically between $`\pm p`$. This is reflected in the symmetry (rather than self-adjointness) of the quantum mechanical $`p`$ operator. Let us now look at the example $`A_2`$ in ($`3.18`$). The free Hamiltonian ($`3.19`$) is self-adjoint on the space $$D_{[0,\mathrm{})}(H_\mathrm{Z}^0)=\left\{\varphi \text{AC}^2[0,\mathrm{})\right\|\varphi (0)=0\},$$ ($`3.21`$) Once again, this is just the “correct” boundary condition for the potential ($`3.10`$), so that the evolution operators $`𝒱(T)`$ form a one-parameter group. One can draw the same conclusions as in the previous example. There is only one difference: the momentum operator $`i`$ is again symmetric, but its deficiency indices are (0,1). This is irrelevant as far as one’s attention is restricted to the Hamiltonian and the Zeno dynamics; however, if one is motivated (on physical grounds) to consider the properties of momentum, the best one can do in this case is to obtain the most appropriate maximally symmetric momentum operator. (We wonder whether this has spin-offs at a fundamental quantum mechanical level.) ## 4 The problem of the lower-boundedness of the Hamiltonian Let us consider now the model Hamiltonian $`H_\mathrm{Z}^0=p=i`$ in $`A_1=[0,1]`$, describing an ultrarelativistic particle in an interval. The mathematical features of this example are very interesting and deserve careful investigation. A similar example was considered in , although in a different perspective. The Zeno dynamics yields $`𝒱(T)=\mathrm{exp}(iTH_\mathrm{Z})E_{A_1},`$ ($`4.1`$) $`\text{with}H_\mathrm{Z}p+V_{A_1}(x),`$ ($`4.2`$) where $`V_A`$ is defined in ($`3.10`$). The “natural” boundary conditions imposed by the Zeno dynamics are $$D_\mathrm{Z}(p)=\left\{\varphi \text{AC}[0,1]\right\|\varphi (0)=0=\varphi (1)\}.$$ ($`4.3`$) In this domain the Hamiltonian $`p`$ is symmetric but not self-adjoint: its deficiency indices are $`(1,1)`$. Therefore, by Stone’s theorem, the Zeno dynamics is not governed by a group and is certainly not time-reversal invariant. More to this, this Hamiltonian is not lower bounded and therefore violates one of the premises of the MS theorem. In order to understand what happens during a Zeno dynamics, look at the first row in Figure 1, where an arbitrary wave packet evolves under the action of the free Hamiltonian $`p`$ (incidentally, notice that the wave packet does not disperse, due to the form of the Hamiltonian). The probability of “surviving” inside $`A_1`$ decreases with time: in other words, even though a “continuous” measurement is performed, in order to check whether the particle is outside $`A_1`$, the particle does leak out of $`A_1`$ and no quantum Zeno effect takes place. Let us now assume, on physical grounds, the validity of periodic boundary conditions: $$D^\alpha (p)=\left\{\varphi \text{AC}[0,1]\right\|\varphi (0)=\varphi (1)e^{i\alpha }\},$$ ($`4.4`$) where the phase $`\alpha `$ determines the specific self-adjoint extension. Notice that this is a physical requirement: it is not a consequence of the Zeno dynamics. The Hamiltonian is now self-adjoint and the dynamics is governed by a unitary group (Stone’s theorem). Obviously, the physical picture given by this self-adjoint extension is completely different from the previous case. See the second row in Figure 1: a quantum Zeno effect takes place. We also stress that the dependence of the Hamiltonian on the $`p`$ operator is not a sufficient condition to yield the behavior described above. In order to clarify this point, let us consider an additional example. Let (we set $`m=1/2`$) $$H=p^2+pH_\mathrm{Z}=p^2+p+V_A(x).$$ ($`4.5`$) We first observe that $`H`$ is lower bounded \[$`p^2+p=(p+1/2)^21/4`$; notice also that this Hamiltonian can be tranformed into the usual form by adding a phase $`x/2`$ to the wave function.\] Consider again the quantum Zeno dynamics on the sets $`A_1`$ and $`A_2`$. Since this is not a classical textbook example, we explicitly derive the deficiencies. In the first case ($`A_1`$) one gets $`(H_\mathrm{Z}\varphi ,\psi )(\varphi ,H_\mathrm{Z}^{}\psi )`$ $`=`$ $`i\overline{\varphi (0)}\psi (0)+\overline{\varphi ^{}(0)}\psi (0)\overline{\varphi (0)}\psi ^{}(0)`$ ($`4.6`$) $`+i\overline{\varphi (1)}\psi (1)\overline{\varphi ^{}(1)}\psi (1)+\overline{\varphi (1)}\psi ^{}(1).`$ It is easy to check that $`H_\mathrm{Z}`$ is lower bounded and self adjoint on the space ($`3.20`$). The Zeno evolution is therefore unitary. In the second case ($`A_2`$) one gets $$(H_\mathrm{Z}\varphi ,\psi )(\varphi ,H_\mathrm{Z}^{}\psi )=i\overline{\varphi (0)}\psi (0)+\overline{\varphi ^{}(0)}\psi (0)\overline{\varphi (0)}\psi ^{}(0).$$ ($`4.7`$) It is straightforward to check that the Hamiltonian is lower bounded and self-adjoint on the space ($`3.21`$). Once again, the Zeno evolution is unitary. ## 5 Discussion One is led to the following question: is it possible to find an example in which the Zeno dynamics is governed by a dynamical semigroup? The answer to this question would be positive if one could find a quantum Zeno dynamics yielding a symmetric, but not self-adjoint, Hamiltonian operator. Indeed, in such a case, by Stone’s theorem one cannot have a group, and by MS’s theorem one must have a semigroup. It would be incorrect to think that the model Hamiltonian $`H=p=i`$ in $`A_1=[0,1]`$ (or even more in $`A_2=[0,\mathrm{})`$) provides us with the counterexample we seek. Indeed, such a Hamiltonian is not a satisfactory example, because it violates one of the premises of the MS theorem, that requires a lower-bounded Hamiltonian from the outset (see beginning of Section 2). We are unable, at the present stage, to give a clear-cut answer to this problem. However, some comments are in order. If, for some reason, the quantum Zeno dynamics yields a symmetric Hamiltonian operator, the search for its self-adjoint extensions seems to us a very important one, on physical grounds. Suppose then that one is willing to consider a self-adjoint extension of the Zeno Hamiltonian $`H_\mathrm{Z}`$. If this is the case, close inspection shows that a quantum Zeno dynamics always yields a group, at least in the class of systems considered in this note. Indeed, a theorem due to von Neumann, Stone and Friedrichs states that “every semi-bounded symmetric transformation $`S`$ can be extended to a semi-bounded self-adjoint transformation $`S^{}`$ in such a way that $`S^{}`$ has the same (greatest lower or least upper) bound as $`S`$.” Therefore, if the Hamiltonian is lower-bounded on the real line, as for instance in ($`3.1`$) (one could even add a non-pathological potential to the kinetic energy), the Zeno dynamics in an interval of $``$ will also be engendered by a lower-bounded Hamiltonian, like in ($`3.16`$); this would always admit a self-adjoint extension (due to the above-mentioned theorem), which in turn would yield a group of evolution operators. Therefore, in order to avoid the consequences of von Neumann’s theorem, the operators arising from the Zeno dynamics must not be lower bounded. Only in such a case the Zeno Hamiltonian might not admit self-adjoint extensions. In conclusion, we have seen that in the situations considered in this paper the quantum Zeno effect yields a unitary dynamics, governed by groups, not by semigroups. We are therefore left with two possible options: i) The MS theorem can be made stronger and the Zeno dynamics is always governed by a group; ii) Different projections, more general than ($`3.2`$)-($`3.3`$), and/or different Hamiltonian operators may yield symmetric Zeno Hamiltonian operators that are not self-adjoint (or, even more, maximally symmetric operators with no self-adjoint extensions) and therefore (due to the MS theorem) a semigroup of evolution operators. The answer to the above alternative would clarify whether a quantum Zeno dynamics introduces some elements of irreversibility in the evolution of a quantum system. This is an interesting open problem. ## Acknowledgments We thank I. Antoniou for interesting remarks.
warning/0004/cond-mat0004051.html	ar5iv	text	# Island nucleation in thin-film epitaxy: A first-principles investigation \[ ## Abstract We describe a theoretical study of the role of adsorbate interactions in island nucleation and growth, using Ag/Pt(111) heteroepitaxy as an example. From density-functional theory, we obtain the substrate-mediated Ag adatom pair interaction and we find that, past the short range, a repulsive ring is formed about the adatoms. The magnitude of the repulsion is comparable to the diffusion barrier. In kinetic Monte Carlo simulations, we find that the repulsive interactions lead to island densities over an order of magnitude larger than those predicted by nucleation theory and thus identify a severe limitation of its applicability. Copyright 2000 by The American Physical Society. \] Island nucleation is often the first step in thin-film epitaxy and is, thus, relevant to the synthesis of a wide variety of interfacial materials. Achieving a quantitative understanding of the island morphologies (i.e., sizes, shapes, density, spatial distribution, etc.) that develop in the initial stages of thin-film growth is also important for fundamental reasons. Since thin-film epitaxy frequently occurs away from equilibrium, the kinetics of deposition and surface diffusion play a key role in governing island morphology and there is great variety in the resulting structures. Considering shapes , for example, islands can be fractal-like or compact and triangular, hexagonal, square, rectangular, etc.. Each of these structures is a signature of an intricate kinetic balance and reflects a complex set of interatomic interactions that is unique for each material. Despite the complexity and potentially enormous variety in growth morphologies, certain aspects of island nucleation and growth appear to be common to many different systems. In a general description, gas-phase species are deposited onto an initially bare solid substrate with a rate $`F`$. These species hop on the surface with a rate $`D=\nu _0e^{E_b^0/k_BT}`$, where $`\nu _0`$ is the preexponential factor, $`E_b^0`$ is the diffusion-energy barrier for an isolated species, $`k_B`$ is Boltzmann’s constant, and $`T`$ is temperature. Hopping mediates the aggregation of adspecies into nuclei, which either dissociate with an energy barrier $`E_{d,i}`$, if they are below a critical size $`i`$, or grow subsequently to become stable islands. Initially, the formation of island nuclei is the main process taking place. As the surface coverage increases, it becomes increasingly likely that deposited species will add to stable islands and promote their growth instead of forming new nuclei. These general features can be captured in a mean-field theory for the stable island density $`N_x`$ . In the island growth regime, this expression has the form $$N_x(F/D)^{i/(i+2)}\mathrm{exp}(E_{d,i}/k_BT)^{1/(i+2)}.$$ (1) Although the utility of a general expression cannot be overstated, Eq. 1 cannot describe all aspects of thin-film epitaxy and it is important to understand its limitations. In the interest of achieving a complete and predictive model for thin-film morphology, it is clearly desirable to have an approach that is as free as possible from arbitrary parameters or assumptions. In this Letter, with an aim toward this ideal approach, we present the results of a combined kinetic Monte Carlo (kMC) and first-principles, density-functional theory (DFT) study of island nucleation in a model for the growth of Ag on a monolayer (ML) of Ag on Pt(111). Our choice of this model system was motivated by intriguing results from recent, low-temperature, scanning-tunneling microscopy (STM) studies , in which Eq. 1 was used to obtain the energy barrier and preexponential factor for adatom hopping. Shown in Table I are the parameters obtained in these studies. A striking feature of the experimental results is that the preexponential factors are significantly smaller than would be anticipated for systems such as these. For example, from ab initio calculations, Ratsch and Scheffler find a preexponential factor of $`\nu _0=1.3\times 10^{12}`$ s<sup>-1</sup> for a Ag adatom on 1-ML-Ag/Pt(111), with a diffusion barrier of $`E_b^0`$ = 63 meV. Inserting the experimental and theoretical values for the diffusion parameters into Eq. 1 in the low-temperature limit where $`i=1`$ and $`E_{d,i}=0`$, we see that the experimental island densities are about an order of magnitude higher than predictions based on the theoretical diffusion parameters. Here, we investigate the origins of this discrepancy. Our DFT-kMC model includes many features of the complex potential-energy surface experienced by Ag adatoms during thin-film growth and is free from several of the assumptions in Eq. 1. We find that one of these assumptions – that interactions between adsorbed species do not extend beyond a short range – is violated. For systems with low diffusion-energy barriers \[such as Ag on 1-ML-Ag/Pt(111)\], we show that these long-range, adatom-adatom interactions play an important and previously underestimated role in island nucleation and growth. The DFT calculations are performed using the plane-wave, pseudopotential method within the generalized gradient approximation . Previously, Ratsch et al. showed in DFT calculations that the diffusion-energy barrier of an Ag atom on the 1-ML-Ag/Pt(111) substrate is essentially the same as that on a strained Ag(111) substrate, in which Ag is given the lattice constant of Pt. Thus, to model the heteroepitaxial system we use strained Ag(111), in which the lattice constant is set to a value of 4.01 Å. This value is 4.61 percent smaller than our calculated lattice constant for bulk Ag. We use the supercell approach to describe the surface, which is modeled as a $`(4\times 4\times 4)`$ slab with a vacuum spacing of five interlayer distances. The cut-off energy is 50 Ry and we use 4 $`𝐤`$ points to sample the full surface Brillouin zone. The top layer of a bare slab is fully relaxed. Subsequently, an adatom is placed in a binding site (fcc and hcp three-fold hollow sites), and its height is optimized with respect to the fixed substrate. To calculate adatom interaction energies, two (or more) adatoms are placed on the relaxed (and fixed) substrate with heights fixed to values from the single-adatom calculations. In this way, we seek to isolate the electronic interaction between adatoms in binding sites. With simultaneous relaxation of both the adatoms and surface atoms, we can resolve the role of substrate-mediated, elastic interactions in the total interaction energy. Full relaxation of a few trial structures and inspection of the forces in our partially relaxed slabs indicates that elastic interactions are not highly dependent on adsorbate configuration and that our results will change by 10 meV or less with full relaxation. The total interaction energy $`\mathrm{\Delta }E`$ for a periodic slab containing $`N`$ adatoms, of which $`M`$ are at binding site $`a`$ and $`(NM)`$ are at binding site $`b`$, is given by $`\mathrm{\Delta }E=E_{S+N}^{a,b}ME_{S+1}^a(NM)E_{S+1}^b+(N1)E_S`$. Here, $`E_{S+N}^{a,b}`$ is the total energy of a slab with $`N`$ adatoms, $`E_{S+1}^a`$ and $`E_{S+1}^b`$ are the total energies of slabs containing one adatom, and $`E_S`$ is the total energy of a bare slab. In the DFT supercell approach, the total interaction energy is comprised of interactions between different adatoms in the slab and interactions between adatoms in the slab and the periodic-image adatoms. We can express $`\mathrm{\Delta }E`$ as a function of these interactions using the lattice-gas Hamiltonian approach (see, e.g.), which yields $`\mathrm{\Delta }E`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j}{}}V^{(2)}(𝐑_{i,j})n_in_j`$ (3) $`+{\displaystyle \frac{1}{3}}{\displaystyle \underset{i,j,k}{}}V^{(3)}(𝐑_{i,j},𝐑_{i,k})n_in_jn_k+\mathrm{}`$ Here, the summations run over all sites $`i`$ in the slab and all sites $`j`$ and $`k`$ in the supercell (which includes both the slab and its periodic images), $`n_m`$ is unity if site $`m`$ ($`m=i,j,k`$) is occupied and zero, otherwise, $`V^{(2)}(𝐑_{i,j})`$ is the pair interaction between two adatoms on sites $`i`$ and $`j`$, and $`V^{(3)}(𝐑_{i,j},𝐑_{i,k})`$ is the trio interaction between three adatoms on sites $`i`$, $`j`$, and $`k`$. We neglect higher-order interactions. Another assumption implicit in Eq. 3 is that the interaction between adatoms at a fixed distance is independent of whether these atoms occupy fcc or hcp sites. We confirmed this assumption in one trial calculation. Finally, the adatom binding energies on fcc and hcp sites are virtually equal: The fcc site is favored by less than 3 meV. Thus, for a given adatom configuration, we express $`\mathrm{\Delta }E`$ as a sum of pair and trio interactions with unknown coefficients. From 18 different configurations, we obtain a system of linear equations and solve these for pair-interaction coefficients up to the $`13^{\mathrm{th}}`$-neighbor, as well as for 5 different trio interactions. We assume that all other interaction coefficients are zero. To verify our parameterization of Eq. (2), we used our interaction parameters to predict the total interaction energy in several, additional test structures. All of the predicted values agreed well with values from DFT calculations. The pair interaction is shown in Fig. 1, where we also show results for Ag on unstrained Ag(111). For both surfaces, this interaction is strongly attractive at the nearest-neighbor distance and repulsive at longer distances. It is interesting to consider the origins of the long-range repulsion. At these distances, the interaction could be due to substrate-mediated elastic interactions or of electronic origin . Since we find that elastic interactions play a small role here, the repulsion is primarily an electronic effect. Each adatom induces a small perturbation in the electron density, which decays with distance from the adatom in an oscillatory manner. The asymptotic tail, which is expected to decay with distance $`d`$ as $`d^5`$ (or as $`d^2`$, if a partially filled surface state is involved), is a Friedel-type oscillation. Friedel oscillations have been imaged as concentric, ring-like, features around defects in low-temperature STM studies of several noble-metal surfaces \[including Ag(111)\] . We expect the Friedel tail to extend to much longer distances than can be probed in DFT calculations. However, interactions associated with the Friedel tail should be weaker than those probed here. Thus, the central-ring interaction resolved here will have the most significant ramifications for thin-film morphology. From Fig. 1(a), we see that the magnitude of the repulsive ring for Ag on strained Ag(111) is comparable to the diffusion-energy barrier for an isolated adatom. For Ag on Ag(111), it appears that the repulsive interaction is weaker and the diffusion-energy barrier is larger. The diffusion barriers reported here are obtained with full relaxation of both the first-layer slab atoms and the adatom. Our barriers are in good agreement with experimental values for Ag on 1-ML-Ag/Pt(111) (60 meV) and on Ag(111) (97 meV) and with those of Ratsch and Scheffler . If the interaction energy and the diffusion barrier are of comparable size, we expect interatomic interactions to significantly influence adatom diffusion and island formation. Since Eq. 1 neglects the influence of long-range interactions, it is unclear if this expression is accurate under these circumstances. To resolve the effect of long-range interactions on thin-film growth, we developed a kMC model employing the general method of Fichthorn and Weinberg and incorporating the pair potential for Ag on strained Ag(111) shown in Fig. 1(a). In the initial stages of thin-film epitaxy, the surface coverage is low and pair interactions are likely to be the only significant interactions governing island nucleation and growth. In our kMC model, atoms are deposited onto a fcc(111) substrate with a rate of $`F=0.1`$ ML/s. An adatom hops from site $`i`$ to site $`j`$ with a rate given by $`D_{ij}=\nu _0e^{E_{ij}/k_BT}`$, where $`E_{ij}`$ is the energy barrier to hop from site $`i`$ to $`j`$. For the hopping-rate parameters, we use $`\nu _0=10^{12}`$ s<sup>-1</sup> . The energy barrier is given by $`E_{ij}=E_{i,j}^{}E_i`$, where $`E_i`$ is the energy with an atom at site $`i`$ and $`E_{i,j}^{}`$ is the energy of the transition state between sites $`i`$ and $`j`$. In general, $`E_{i,j}^{}`$ should depend on both $`E_i`$ and $`E_j`$. Considering possible permutations of adatom configurations with $`13^{\mathrm{th}}`$-neighbor interactions, $`10^{14}`$ different, diffusion-energy barriers could occur. To make the problem tractable, we adopt a simple model, in which $`E_{ij}=E_b^0+\frac{1}{2}(E_jE_i)`$. All of the quantities in this equation are obtained from DFT calculations. We have tested this equation for a trial geometry in which an adatom with four fcc $`9^{\mathrm{th}}`$ neighbors hops to a nearest-neighbor hcp site where it has two $`7^{\mathrm{th}}`$ and two $`12^{\mathrm{th}}`$ neighbors. From our simple model, we find $`E_{ij}`$ = 53 meV, which is in remarkable agreement with the value of 46 meV from DFT calculations. We simulated thin-film epitaxy over temperatures ranging from 40-70 K and determined island densities in the beginning of the island growth regime. These low temperatures are in the range of the experimental studies (cf., Table I). At such low temperatures, Eq. 1 reduces to the form $`N_x(F/D)^{1/3}`$. Fig. 2 shows an Arrhenius plot of the island density from our DFT-kMC model as a function of temperature. Also shown in Fig. 2 is the island density predicted by nucleation theory for the values of $`F`$, $`\nu _0`$, and $`E_b^0`$ used here. To quantitatively compare nucleation theory with the simulations, a proportionality coefficient $`\eta `$ is needed in Eq. 1 (i.e., $`N_x=\eta (F/D)^{1/3}`$). This coefficient is related to the efficiency of the islands in capturing adatoms. Using a self-consistent approach, $`\eta `$ = 0.25 and values of $`\eta `$ ranging from 0.2 to 0.23 have been found in kMC simulations of Ag island nucleation on Pt(111) . Here, we use $`\eta `$ = 0.25. In Fig. 2, we see that the DFT-kMC island densities are an order of magnitude (or more) above the theoretical values. To understand this, we construct a caricature model, in which we replace the set of pair interactions shown in Fig. 1(a) with a nearest-neighbor attractive interaction and a uniform, repulsive ring of strength $`\epsilon _R`$ at distances 10-13. By varying the magnitude of $`\epsilon _R`$, we span the entire range of possible behaviors in this system. As $`\epsilon _R\mathrm{}`$, the island density assumes a constant, maximum value that is independent of temperature (cf., Fig. 2). This is because island nucleation can only occur when one atom is deposited within the repulsive ring of another and it is governed by the temperature-independent deposition rate. In this regime, many adatoms are isolated by repulsion in the initial stages of deposition. Each isolated adatom becomes a stable island when another atom is deposited into its ring and the resulting island density is significantly higher than in the absence of such a ring. As $`\epsilon _R`$ is decreased, diffusing adatoms are increasingly able to surmount the ring barrier and a second channel for island nucleation and growth (via long-range, adatom diffusion) opens up. The extent to which long-range diffusion contributes to island nucleation and growth depends on the temperature. In Fig. 2, we see that at 40 K, the DFT-kMC island density is the same as that for an infinitely repulsive ring (i.e., diffusing adatoms are unable to penetrate the ring on the time scale for nucleation). As the temperature increases, adatoms are increasingly able to penetrate the ring barrier to aggregate and add to existing islands via long-range diffusion. Consequently, the island density decreases with increasing temperature. It is interesting to note that for the conditions studied, even a relatively weak repulsive ring with $`\epsilon _R`$ = 25 meV can lead to significantly higher island densities than those predicted by nucleation theory. Returning to our discussion of the experimental results shown in Table I, we point out that the order-of-magnitude difference between the island densities predicted from ab initio calculations and those found experimentally for Ag on 1-ML-Ag/Pt(111) is also seen in our study, comparing the island densities predicted by nucleation theory to those found in our DFT-kMC “computer experiments” (cf., Fig. 2). Thus, we conclude that our results can explain the theoretical-experimental gap in the island density for Ag on 1-ML-Ag/Pt(111). Further, our results indicate that for Ag on Ag(111), the theoretical-experimental gap should be weaker or non-existent. This result is also consistent with a comparison of theoretical diffusion parameters for Ag on Ag(111) ($`\nu _0=8.2\times 10^{11}`$s<sup>-1</sup>, $`E_b^0=82`$ meV) to experimental values obtained using Eq. 1 ($`\nu _0=2\times 10^{11}`$s<sup>-1</sup>, $`E_b^0=`$ 97 meV). Finally, Bogicevic and co-workers recently found similar DFT and kMC results for both Al(111) and Cu(111) homoepitaxy. Thus, we conclude that long-range, electronic, substrate-mediated adatom interactions exist and, if their strength is comparable to the diffusion barrier, they can significantly influence surface diffusion and the growth morphology in thin-film epitaxy. For Ag on strained Ag(111), the adatom pair interaction becomes repulsive past the short range and the repulsion forms a ring around isolated adatoms. The magnitude of the repulsion is comparable to the diffusion barrier. By inhibiting island nucleation and growth via long-range adatom diffusion, these interactions lead to island densities that are substantially larger than those predicted by nucleation theory. We acknowledge helpful conversations with A. Bogicevic, H. Brune, P. Kratzer, C. Ratsch, A. Seitsonen, and C. Stampfl. Support for this work is from the Alexander von Humboldt Foundation and the NSF (DMR-9617122).
warning/0004/astro-ph0004108.html	ar5iv	text	# Circumnuclear regions in barred spiral galaxies I. Near-infrared imagingPartly based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract NAS 5-26555. ## 1 Introduction Observations and modelling of the circumnuclear regions (CNRs) of barred galaxies give important clues about the nature of active galactic nuclei (AGN), (circum)nuclear starbursts and the gas dynamics in the central kpc (e.g. Knapen et al. 1995a,b; Buta & Combes 1996; Elmegreen et al. 1997; Quillen et al. 1999; Regan & Mulchaey 1999; see also reviews by Knapen 1999 and Shlosman 1999). The presence of non-axisymmetries, such as oval distortions or bars, in the gravitational potential of a disc galaxy can lead to the infall of gaseous material, from galactocentric radii of a few kpc, into the central kpc. The nonaxisymmetry facilitates the removal of angular momentum from the gas, and directs it into the CNR or possibly the nucleus (Shlosman et al. 1990; Athanassoula 1992; Phinney 1994). This gas fuelling process is expected to induce morphological signatures in the bars and CNRs, which are useful tools in the understanding of the dynamics and kinematics of galaxies. Dust lanes and nuclear rings appear to be the dominant circumnuclear features. Dust lanes have been interpreted as the location of shocks in the gas flow. Thus, the morphology of the dust distribution can reveal characteristics of the dynamics of these galaxies (Athanassoula 1992) and of the kinematics of the gas. According to Athanassoula’s models, the degree of curvature of the dust lanes is a direct indicator of the strength of the bar for bars which are not very strong. Nuclear rings, which are usually sites of active star formation (SF), occur at the location of strong density enhancements in the gas, where the bar-driven inflow of gas slows down in the vicinity of inner Lindblad resonances (ILRs; Athanassoula 1992; Heller & Shlosman 1994; Buta & Combes 1996; Shlosman 1999). SF may result from the gas becoming gravitationally unstable (Elmegreen 1994, 1997), or from triggering in miniature spiral arms (Knapen et al. 1995b, 2000). Whereas imaging in the blue, ultraviolet, or in spectral lines like H$`\alpha `$ traces the massive SF (e.g. Sersic & Pastoriza 1967; Pogge 1989a,b; Knapen et al. 1995a; Maoz et al. 1996; Colina et al. 1997), near-infrared (NIR) imaging offers very important advantages. Firstly, the presence of substantial and mostly unquantified amounts of extinguishing (absorbing and scattering) dust in the CNRs hampers the interpretation of imaging at especially the UV and optical wavelengths. NIR, particularly $`K`$-band, emission is much less susceptible to extinction by dust (a factor of 10 when comparing $`K`$ to $`V`$), whereas colour index images such as $`IK`$ (Knapen et al. 1995a) or $`JH`$ are clear morphological dust indicators due to the relatively modest changes in those colours caused by different stellar populations. Secondly, NIR observations are of particular importance to the study of the CNRs because they allow a better determination of the mass distribution (excluding dark matter). NIR imaging is more sensitive to light primarily from cool giants and dwarfs which dominate the mass or are at least directly proportional to it. However, as shown in the literature and by us in this paper, there can be significant, or even dominant, contributions from young stars to the NIR emission even in the $`K`$-band, so for objects with strong SF any estimate of the mass distribution made from NIR maps must take the varying $`M/L`$ ratio into account (Knapen et al. 1995a,b; Elmegreen et al. 1997; Wada, Sakamoto & Minezaki 1998; Ryder & Knapen 1999). The atlas of images presented in this paper (Paper I) represents the first results of a programme which studies the circumnuclear structures that appear at small scales (a few hundred parsecs to a kpc) and their connection with the global disc structure in a sample of 12 barred spiral galaxies. We present $`J,H`$ and $`K`$-images at subarcsec resolution for all 12 sample galaxies, and Hubble Space Telescope (HST) archive $`H`$-band images for 10 of them. In Paper II (Pérez–Ramírez & Knapen 2000a), we present an accompanying optical data set of broad-band and H$`\alpha `$ images of the complete discs of all our sample galaxies, while in Paper III (Pérez–Ramírez & Knapen 2000b) the morphological information is interpreted in terms of the structure and evolution of CNRs of barred galaxies. In Section 2, we describe the observations and data reduction techniques. The imaging data are presented as a series of multi-panel figures, which show broad-band images, colour index images highlighting dust and SF features against the dominant stellar luminosity, and radial profiles of ellipticity, position angle, surface brightness and colour. The results are summarized in Section 3 for individual galaxies, and briefly discussed in Section 4, while concluding remarks are given in Section 5. ## 2 Sample Selection, Observations and Data Reduction ### 2.1 The sample The main selection criterion for our sample is the presence of a bar and evidence of some circumnuclear structure associated with it, such as rings, nuclear bars or regions of SF. Furthermore, the sample galaxies should be nearby, bright and observable from the northern hemisphere. Some of the sample galaxies are taken from the lists published by Sersic & Pastoriza (1967) and Pogge (1989a,b). The sample can be considered more anecdotal than complete in any sense. However, it does significantly increase the sample size when compared to published NIR studies of CNRs in barred galaxies, which rarely include more than three objects and usually discuss only one (e.g. Knapen et al. 1995a,b; Elmegreen et al. 1997; Laine et al. 1999; Ryder & Knapen 1999). With our new imaging, we also improve the spatial resolution, and show different colour index maps, $`JK`$ (primarily outlining dust extinction) and $`HK`$ (additionally indicating the possible emission due to hot dust in the case of the more active sample galaxies). Regan & Mulchaey (1999) show HST 1.6$`\mu `$m and optical-infrared (basically $`RH`$) colour index images of a sample of Seyfert (Sy) galaxies, including a number of CNR galaxies. Two of the galaxies in their sample are also included in our sample (NGC 3516 and NGC 3982). The $`RH`$ colour, however, is much more sensitive to changes in stellar populations than the NIR colours used in our paper. For 10 of our sample galaxies, we have retrieved, in addition to our own data, $`H`$-band images from the HST archive (Fig. 1), and we comment on the high-resolution NIR morphology of those objects. In Table 1 we give some information about the classification of the galaxies of our sample, the predominant circumnuclear feature, as determined from the data, and the type of nuclear activity, if known. ### 2.2 Ground-based observations We have obtained NIR broad-band images in the $`J`$ (1.25$`\mu `$m), $`H`$ (1.65 $`\mu `$m) and $`K`$ (2.2 $`\mu `$m) bands of the 12 barred galaxies that make up our sample, during three observing runs: 1994 September 26, 1995 November 5–6, and 1996 February 3–5. The observations were made with the 3.6 m Canada–France–Hawaii Telescope (CFHT). We used the Montreal NIR camera (MONICA; Nadeau et al. 1994), equipped with a 256 x 256 pixel HgCdTe array detector with a projected pixel size of 0$`.^{\prime \prime }`$248. Since the sky background in the NIR is bright and changes rapidly, it is essential to obtain sky frames frequently. We observed the sky usually just before and after each galaxy observation. Frames typically consisted of four co-added individual exposures. Separate frames with the nucleus of the galaxy at slightly offset positions were co-added in the reduction process to produce the final mosaic images. The weather was clear and generally photometric thoughout these nights, and the seeing values as measured from our final images were 0$`.^{\prime \prime }`$7 – 1$`.^{\prime \prime }`$0. ### 2.3 Data reduction The main steps in the reduction of the NIR data include subtracting the sky background, interpolating across known bad pixels, and registering and combining the images. The data reduction was done partly with private programmes and partly with standard iraf routines. We first combined four sky exposures taken at different positions on the sky immediately before and after a series of galaxy exposures by median-averaging them, while iteratively rejecting those values in the averaged sky exposure that deviate by more than 3$`\sigma `$ from the average value in that frame. This procedure allows the automated rejection of any star images that might have been present in the sky exposures, which were observed by offsetting the telescope blindly to a position a few arcmin from the galaxy centre. We then subtracted the averaged sky image from the relevant galaxy images. A close inspection of the sky-subtracted galaxy images of the 1996 run showed the presence of periodic, horizontal lines which could be traced back to electronic crosstalk in the system. This raises the noise level in these images. In order to correct the problem, we used a Fourier transform technique to locate the frequencies of the maximum intensity in the images, and filtered them out. This resulted in a considerable improvement in most images, although in a few images interference stripes are still notable at low levels. After flat-fielding with dome flats, and masking the unreliable pixels in the array, we combined the sky-subtracted and, where appropriate, de-striped images. Ideally, one would like to locate one or more common field stars in the images to achieve the most accurate alignment possible. Unfortunately, the field of view of our images is small and in most cases we could only use the nucleus of the galaxy itself to determine the reference position. The individual images were thus registered and shifted to a common position (in all cases accurate to a fraction of the pixel size) and averaged. In the process, the pixel size was halved to increase sampling (reduced from 0$`.^{\prime \prime }248`$ to 0$`.^{\prime \prime }124`$). Finally, we rotated the images by an angle of 88.8 in order to obtain the correct north (up), east (left) orientation. ### 2.4 Photometric Calibration Although we subtracted the background sky as described in Section 2.3, we cannot check on how accurate our sky subtraction is since the sky is variable, and the frame is small compared with the galaxy. This makes it very hard to estimate to what level the flux registered in the photometric aperture is influenced by any residual background emission. For that reason, we have used aperture photometry from the literature. By calibrating with multiple aperture measurements we can better estimate the value of the sky and the efficiency of the system. If only one useful aperture is available, we can use only the relatively bright part of the image. We were able to find the aperture photometry in $`J`$, $`H`$ and $`K`$-bands for all but one of the galaxies, NGC 1530. For this galaxy, we used the average magnitude offset as determined for the rest of the galaxies. Literature photometry sources include Glass (1976), Aaronson (1977), McAlary, McLaren & Crabtree (1979), Balzano & Weedman (1981), Willner et al. (1985) and Spinoglio et al. (1995). To check the quality of our photometry, we also made comparisons with data from the literature. We selected values from Glass (1984), Cidziel, Wynn-Williams & Becklin (1985), Devereux, Becklin & Scoville (1987) and Hunt et al. (1994). The average absolute difference between our results and the results from different authors was 0.12 mag. Direct comparison of the surface brightness profile in the $`H`$-band and the $`HK`$ colour profile plots was possible for two of our sample galaxies, NGC 3516 and NGC 3982, which are in common with the Peletier et al. (1999) sample. The photometric calibration procedures used were slightly different, but the agreement in both cases was good. ### 2.5 Colour index maps and radial profiles Colour index maps were created after assuring that the relevant pairs of images were at the same pixel scale and orientation and at comparable resolution before combining them. After calibrating the images, we used galphot (Jørgensen, Franx & Kjaergaard 1992) to fit ellipses to the images. The centre position, ellipticity and the position angle of the fitted ellipses were allowed to change freely as a function of radius for the $`H`$ passband. The values for these parameters were then used as input for the fits to the images in other passbands, in order to ensure that we produce reliable colour profiles. In the case of NGC 2903, the position angle and ellipticity were kept constant because the large amount of structure present in the core did not allow a meaningful fit to these parameters. We thus produced plots (shown in Fig. 2) of the surface brightness of the $`H`$-band as a function of radius, and colour profiles in $`JK`$ and $`HK`$ for all the galaxies which were imaged in these bands. Errorbars in these plots indicate the uncertainty in the determination of the sky background. We have estimated these uncertainties using the procedure described by Peletier et al. (1999), where 1$`\sigma `$ errorbars correspond to $`\mu _H`$= 21.5 mag arcsec<sup>-2</sup>, $`\mu _J`$= 21.2 mag arcsec<sup>-2</sup> and $`\mu _K`$= 20.5 mag arcsec<sup>-2</sup> (about 0.1% of the sky background in $`H`$ and $`K`$ and 0.3 % in $`J`$). In Fig. 2, we show for each galaxy (except NGC 3351 for which we only obtained a $`K`$-band image) greyscale representations of the $`K`$ broad-band and $`JK`$ and $`HK`$ colour index images, radial profiles of position angle and ellipticity as determined from the fit to the $`H`$-band image ($`K`$ for NGC 3351), and radial surface brightness ($`H`$) and colour ($`JK`$ and $`HK`$) profiles. We tabulated the values plotted as a function of radius, but publish the tables in electronic form only<sup>1</sup><sup>1</sup>1The data tables are available electronically from the Centre de Données astronomiques de Strasbourg (CDS), on: ftp://cdsarc.u-strasbg.fr/pub/cats/J/MNRAS/volume/first\_page .. They are also available from the authors. ### 2.6 HST NIR imaging For 10 of our 12 sample galaxies, NICMOS imaging with the HST is available from the HST archive. We retrieved the re-reduced F160W (comparable to $`H`$-band) images from the archive. However, we improved the quality of some of these images by doing additional data reduction to remove artifacts. We relied on header information to place the images on an astrometrically correct grid. The images are all taken with the NIC2 camera, with a pixel size of 0$`.^{\prime \prime }`$075. In most cases, the image retrieved from the archive is a combination of several individual exposures. For NGC 3516 and NGC 3982, however, one single exposure was available, and these images in fact improved most due to our additional data reduction. Regan & Mulchaey (1999) published two of these images (NGC 3516 and NGC 3982). We show the central areas of all images in Fig. 1, with the same scale and orientation. In most of the objects, a wealth of detail can be seen in the CNR. The emission usually coincides with the location of the SF ring, and most of the individual bright knots are due to regions of current SF. Dust lanes and/or SF regions often outline spiral-like patterns, which will be discussed in more detail below, and in Paper III. ## 3 Results on individual galaxies In this section, we describe some of the results of the NIR imaging of our sample galaxies, as shown in detail in Fig. 1 (HST NIR images) and Fig. 2 (ground based multi-band images and profile fits). A more systematic study of parameters derived from these data in combination with optical imaging of the complete galaxy discs is forthcoming (Papers II and III). ### 3.1 NGC 1300 Our broad-band NIR images are remarkably smooth, and do not show any structure in the CNR (Fig. 2a). The colour index images, however, show a red ring-like structure, possibly outlining a single spiral arm that departs from the nucleus towards the west side of the image and continues to wrap around the nucleus until it closes in a ring. This ring appears more continuous and broader in the northern part than in the south where it looks patchy. The feature seen in our NIR colour index map cannot be an artifact of combining two images with slightly different spatial resolution, because such an artifact would show up as a complete ring, whereas the observed red feature is not at constant galactocentric radius. Pogge (1989a) saw an incomplete nuclear ring in H$`\alpha `$ emission with a number of distinct “hot spots”. The HST $`H`$-band image resolves the nuclear ring into a series of tightly wound spiral armlets, outlined in dust and stellar populations. The ring region is of a relatively low amplitude compared to the central bulge component, in contrast to nuclear rings in other objects, e.g. NGC 3351, as can be observed in Fig. 1. The location of the ring corresponds to a bump in all radial profiles (i.e., surface brightness, colour, ellipticity and position angle). The $`JK`$ colour of the ring is redder by 0.1 magnitudes than the background. No isophotal twists are seen in the ellipticity and position angle profiles or in the contour map of the $`K`$-band image. Thus we confirm that there is no evidence for a nuclear bar (Regan & Elmegreen 1997). ### 3.2 NGC 1530 We find very well-defined mini-spiral structure in the central region of this galaxy (Fig. 2b). Reynaud & Downes (1997; see also Reynaud & Downes 1998, 1999) suggest that the molecular gas distribution in the central 5 kpc of this galaxy is concentrated along arc-like shock features. Our $`JK`$ image shows that these arcs are part of a well-defined nuclear spiral structure. Furthermore, Reynaud & Downes suggest that a patchy molecular ring may lie inside these shock fronts, possibly connected to them. Fig. 2b reveals that the spiral structure forms a pseudoring around the nucleus. Even the broad-band image shows considerable structure in the CNR. No isophote twists can be seen in the position angle plot. Distinct regions, presumably of enhanced SF, can be noticed along the mini-spiral in the broad-band NIR images. These regions coincide in position with the dark (red) lanes seen to outline the mini-spiral in the $`JK`$ image. This implies that the mini-spiral in $`JK`$ is not necessarily a dust spiral (as would have been tempting to conclude from the colour index image alone) but may well be delineated by either emission from young stars (e.g. red supergiants), or from hot dust. We conclude that we see a mini-spiral in the NIR, outlined by emission from SF and dust. This mini-spiral is well traced also on the HST $`H`$-band image (Fig. 1) as a collection of distinct luminous regions outlining the armlets, which are accompanied by dust lanes. Like in the case of NGC 1300, a bump in the radial profiles is visible at the location of the ring. The $`JK`$ colour of the ring is 0.05 magnitudes redder than the background colour. No pronounced change of slope is detected in the $`H`$-band surface brightness profile, which overall decreases more smoothly than in the other galaxies discussed here. ### 3.3 NGC 2903 Fig. 2c shows several peaks of SF in the CNR of this starburst galaxy, which, as expected, are resolved in much more detail in the HST NICMOS image we obtained from the HST archive (Fig. 1). These peculiar ‘hot spots’ in the nuclear region have been identified and described in different ways by various authors. Marcelin, Boulesteix & Georgelin (1983) found six ‘hot spots’ forming a linear structure that crosses the central region. A later study in infrared and radio by Wynn–Williams & Becklin (1985) showed that bursts of SF were not confined to the visible hot spots. Simons et al. (1988) noticed the presence of organized dust structure in the CNR using NIR and optical imaging. We see patches of dust in our images, but no coherent structure. Regan & Elmegreen (1997) found that this galaxy has an isophotal twist, using a $`K`$-band image of the inner bar region, although the results from an earlier study by Elmegreen et al. (1996) were ambiguous. Similar twists have been interpreted as nuclear bars or triaxial structures (e.g., Shaw et al. 1995; Friedli & Martinet 1993; Wozniak et al. 1995). However, our images show such an abundance of structure that it is hard to imagine that the measured changes in the radial behaviour of ellipticity or PA would have any significance in this respect. For that reason, we set the ellipticity and the PA to fixed values corresponding to the parameters of the isophotes on the outskirts of our images when we fit our radial colour profiles. As in, e.g., NGC 1300, the circumnuclear ring-like region of enhanced SF shows up in the surface brightness profile as a change of slope, and in the colour profiles as a significant red bump at the same radius. Our colour index maps suggest that the nuclear hot spots may be part of a pseudoring. We can identify more than six clumps in the broad-band images, and even more in the colour index map. This is in good agreement with the idea that obscuration by dust is the reason for a previous underestimate of the extent of this SF region (cf. Jackson et al. 1991). As in NGC 1530, we see significant emission from SF regions in the NIR, including the $`K`$-band, which must at least in part be due to red supergiant emission. The red $`JK`$ and $`HK`$ colours of these regions makes them stand out clearly in the colour index maps. ### 3.4 NGC 3351 We only obtained a $`K`$-band image (Fig. 2d), which shows an incomplete ring of presumably SF clumps. The HST $`H`$-band image (Fig. 1) shows the same ring but at significantly higher resolution, and resolved into a large number of individual luminous regions, presumably star-forming. Some of the brightest knots can be recognised in the UV ($`2200`$Å) HST image published by Colina et al. (1997), but others, e.g. the one north of the nucleus, are absent from the UV image. Dust extinction is the most likely candidate mechanism for this absence, although stellar population differences can be envisaged as culprits as well. Spiral arm structure is not obvious in the CNR, and if it is present at all it is very tightly wound. There is only little evidence of dust organised into lanes, and the best example of such a dust lane is seen towards the south-east in the HST image. In contrast, and as a direct result of their much lower spatial resolution, Shaw et al. (1995), found evidence for a circumnuclear ring only in their $`JK`$ colour map. As in the other galaxies discussed before in this section, there is a dichotomy in the slope of the $`K`$-band surface brightness profile. The change occurs at the radius of the SF “ring”, and can be explained as the transition between an inner active SF region and the quiescent disc around it. PA and ellipticity are difficult to fit in that area, but their radial profiles show consistent behaviour in- and outside the nuclear ring. ### 3.5 NGC 3504 Our NIR imaging of this barred galaxy reveals a lot of structure, including a double peak in the central part of the NIR broad-band images (Fig. 2e). Recent adaptive optics images obtained with the CFHT (by F. Combes & J.H. Knapen, private communication) confirm this double peak, but spectroscopic follow-up observations are needed to confirm whether these peaks are in fact two nuclei. No HST NIR imaging is available. The colour index images show a pair of long and straight dust lanes that come into the CNR through the main bar, and intricate dust lane structure in the central region. The nuclear double peak is very obvious in the colour index maps. Elmegreen et al. (1997) showed a ring in $`JK`$ with five discrete clumps of SF. The signal to noise ratio in our images is too low to reveal these SF clumps. The previously reported isophote twist (e.g., Pompea & Rieke 1990) is obvious from our imaging, as well as from the PA and ellipticity profiles. The double bar, outlined by the two separate peaks in the ellipticity profile, is most probably not exclusively due to the nuclear double peak, because the accompanying isophote twist in fact starts somewhat outside the region of influence of the double peak. The surface brightness and colour profiles show the hump at the CNR radius that is seen in many similar galaxies presented in this paper. ### 3.6 NGC 3516 The only obvious feature in our images is the very red nucleus. Neither the broad-band HST $`H`$-image, nor our ground-based NIR colour maps reveal any other structure (Fig. 1, 2f). Regan & Mulchaey (1999) used an HST WFPC2–NICMOS colour index map to show that a single spiral dust pattern dominates the circumnuclear morphology of this galaxy. They described a strong red dust lane that emerges from a blue feature north of the nucleus at a radius of 3<sup>′′</sup>. The spatial resolution of our images is not high enough to reveal such fine details. The differences between our images and HST images are probably due to the longer spectral baseline that Regan & Mulchaey used (0.55 $`\mu `$m to 1.6 $`\mu `$m) and the different spatial resolution. Quillen et al. (1999) noticed that the inner $`J`$-band isophotes are slightly elongated in a direction roughly perpendicular to the outer bar, so the galaxy may be doubly barred. Our $`H`$-band profile indeed shows an isophote twist of about 50 (Fig. 2f). Quillen et al. (1999) saw a curved dust feature at about 4 arcsec south of the nucleus in their HST WFPC2–NICMOS colour index map. ### 3.7 NGC 3982 We see a multi-armed spiral pattern in our images (Fig. 2g), in agreement with Regan & Mulchaey (1999) who describe the global morphology of this galaxy as multi-armed. There is a lot of structure in our ground-based images, but a detailed comparison with the colour index maps obtained by Regan & Mulchaey (1999) is difficult due to the lower spatial resolution of our images. The HST NIR image (Fig. 1) is the one used by Regan & Mulchaey, and shows the rather faint spiral structure in the CNR. Our colour index map shows that this galaxy has a small red nucleus, classified as Sy2. The ellipticity of the isophotes reaches a maximum at a radius of 9 arcsec. This could correspond to a ring, a small bar or a triaxial bulge. The radial colour profiles ($`JK`$ and $`HK`$) follow a characteristic shape, becoming very red close to the nucleus, with the colour most likely due to emission from dust heated by the AGN radiation field. The colours become bluer until a certain radius (the location of the ring), after which they remain constant. Such a profile shape is only seen in the AGNs of our sample. The difference in colour between the nucleus and the ring radius is about 0.2 magnitudes in both $`JK`$ and $`HK`$. ### 3.8 NGC 4303 Buta & Crocker (1993) classify this galaxy as having a nuclear ring, based on their H$`\alpha `$ data. Elmegreen et al. (1997), using NIR observations, did not detect any ring. According to them, the ring consists of very young stars which do not show up in the NIR. However, we can see a well-defined ring in our $`JK`$ and $`HK`$ images (Fig. 2h). There is a pair of dust lanes which connect the bar to the nuclear ring in the south and north. The reddest colours are seen where the dust lanes merge with the nuclear ring. Colina et al. (1997) present a UV ($`2200`$Å) HST image of the CNR of NGC 4303, which shows spiral structure outlining massive SF, continuing all the way into the unresolved core on the NE side of the nucleus. The UV SF spiral is strongest on the side opposite to where we see the largest concentrations of dust (darkest patches in Fig. 2h) in our NIR colour index maps. Our HST $`H`$-band image shows some spiral structure in the NIR but emission is dominated by the central bulge component. There are strong isophotal twists within the central 7 arcsec, and a possible nuclear bar with a radius of 2 arcsec where the ellipticity reaches a peak of 0.2 and the position angle is constant at about 220. The signature of the ring can be recognised as a red peak in the $`JK`$ and $`HK`$ colour profiles, as well as in all other profiles. ### 3.9 NGC 4314 Our ellipticity and PA profiles provide some evidence for a nuclear bar with a radius of 1–2 arcsec. The corresponding PA twist can be recognized in the contours of the broad-band image. Benedict et al. (1993), using the WFPC camera on the HST, found an oval distortion with a length of 8 arcsec in the nuclear region. Wozniak et al. (1995) suggested that this galaxy could be another example of a double-barred galaxy, although they found it difficult to confirm this with their ground-based image at a considerably lower spatial resolution than our new data. The colour maps published by Shaw et al. (1995) revealed the presence of the circumnuclear ring, which is so prominent in the optical (e.g., Morgan 1958), and especially in H$`\alpha `$ (e.g., Pogge 1989a). This ring is redder than all the other regions in the galaxy, having colours consistent with those of typical old stellar populations in ellipticals and spiral bulges. We see a smooth and continuous ring in our $`JK`$ colour index map (Fig. 2i). The HST $`H`$-image (Fig. 1) shows a wealth of emitting structure in narrow tightly wound spiral armlets in the nuclear ring. Fine dust lanes are seen to accompany the spirals. The radial profiles also show the signature of the ring, most clearly with colours that are about 0.1 magnitude redder than the background. ### 3.10 NGC 4321 The prominent star-forming CNR of NGC 4321 (M100) has been studied in great detail using, e.g., optical and NIR imaging, CO interferometry, and modelling (e.g., Pogge 1989; Knapen et al. 1995a,b; Knapen 1998; Wada et al. 1998; García–Burillo et al. 1998; Ryder & Knapen 1999; Knapen et al. 2000). Knapen et al. (1995b) concluded that M100 has a circumnuclear starburst maintained by a global bar-driven density wave. As already shown in detail by Knapen et al. (1995a), the $`K`$-image of this region is generally smooth, in contrast to the appearance in optical and H$`\alpha `$. Two symmetrically placed ‘hot spots’, named K1 and K2, are obvious in Fig. 2j, both in emission in the broad-band image, and as red features in the colour index images. Ryder & Knapen (1999) recently used NIR imaging and spectroscopy to confirm the suspicion that K1 and K2 are in fact regions of enhanced SF, and the $`K`$ emission from those regions is partly due to young stars (Knapen et al. 1995a,b). Our NIR imaging confirms the location of dust lanes and suspected SF regions, shown by Knapen et al. (1995a) in their $`IK`$ colour index map. The locus of the circumnuclear ring-like structure shows up prominently in all radial profiles. Unfortunately, no HST NIR images are available. ### 3.11 NGC 5248 NGC 5248 is a galaxy with an Hii nucleus and a lot of SF activity in the CNR. Elmegreen et al. (1997) found very conspicuous central spiral arms, and several hotspots that form a ring-like spiral pattern. Buta and Crocker (1993) detected a nuclear ring with a diameter of 10–17 arcsec. This activity shows up clearly in our broad-band NIR images, as well as in our colour index maps (Fig. 2k). Spiral structure with star-forming arms, accompanied by dust lanes, is the dominant feature. The western spiral arm has colours which are redder by about 0.1 mag in $`JK`$ than its counterpart in the east. Our $`JK`$ image has been published earlier by Laine et al. (1999). They compared it with the images obtained using adaptive optics, which show a nuclear grand-design spiral structure. This nuclear spiral, at scales of tens of pc, is not expected to show up in a single broad-band NIR image, even at HST resolution, and in fact does not show up (Fig. 1). The HST $`H`$-band image does show a wealth of structure in the CNR, again in the form of emitting regions distributed along spiral arm fragments, and accompanied by less luminous regions which may well be dusty. As in other galaxies, we can see the signature of the ring as a peak at a radius of $`7`$ arcsec in all radial profiles. There is no evidence for nested bars. ### 3.12 NGC 6951 The dust structure in the circumnuclear ring is the most conspicuous feature in our $`JK`$ colour index map, but the star-forming regions in the ring can also be seen in the broad-band images (Fig. 2l). The bar dust lanes connect to this nuclear ring in the northeast and southwest. Several sites of SF are located along the ring, and its presence is also seen in the ellipticity and PA profiles, as well as in all other profiles (Fig. 2l). We can also see a blue ring at about 2 arcsec, within the red nuclear ring surrounding it. This blue ring shows up as a dip in the $`JK`$ profile. Its nature is not clear, and needs further study. The HST $`H`$-band image, used also by Pérez et al. (2000), shows a picture also seen in, e.g., NGC 3351 and NGC 4314, namely of bright emitting knots distributed along tightly wound spiral armlets in the nuclear ring, accompanied by dust lanes. ## 4 Discussion We have presented examples of circumnuclear rings and dust lane patterns in the CNRs of a dozen barred galaxies. In most galaxies, rings are well defined, in others, nuclear spirals with different curvature radii are a predominant circumnuclear feature. There are also some galaxies where both features coexist. Our sample is composed of 12 galaxies, most of which have nuclear rings. Six of them also have clear nuclear mini-spirals (NGC 1530, NGC 3504, NGC 4314, NGC 4321, NGC 5248 and NGC 6951), while in a few others there is less clear or circumstantial evidence for mini-spirals. ### 4.1 Circumnuclear Rings Simulations (e.g. Athanassoula 1992; Byrd et al. 1994; Heller & Shlosman 1994; Knapen et al. 1995b; Piner, Stone & Teuben 1995; see also review by Shlosman 1999) have shown that circumnuclear rings can arise as a consequence of bar-driven inflow and the existence of dynamical resonances in the bar. Shocks form in the gas in regions of orbit crowding along the leading edges of the bar, and dense gas accumulates in these regions. In the shock region the gas loses angular momentum through torques exerted by the bar, and flows inward. The pattern of the gas flow can be quite complex, but if two ILRs are present, inflowing gas accumulates in a ring between them (reviewed by Shlosman 1999). If there is no ILR, the gas may continue to flow inward, resulting in a nuclear starburst rather than a ring (Telesco, Dressel & Wolstencroft 1993). Other mechanisms for the origin of circumnuclear rings that have been put forward as possible origins include shear in the differentially rotating disc (Buta & Combes 1996), supermassive black hole binaries (Taniguchi & Wada 1996), and minor mergers (Taniguchi 1999). Our sample galaxies were selected to have some kind of circumnuclear structure, and except in NGC 3516 and NGC 3982, all have clear circumnuclear star-forming ring-like regions. SF is visible in most galaxies in the high-resolution HST NICMOS images (Fig. 1), and also in a number of galaxies in the ground-based images. However, in all galaxies with circumnuclear SF, dust lane structure outlining small-scale spiral arms is present, as can be seen in the colour index maps (Fig. 2). This is true even for NGC 1300 where the SF regions are conspicuously absent from our ground-based NIR data and very weak in HST NIR images. We conclude that in this class of barred galaxies with circumnuclear SF, dust lane structure in the CNR is always direct evidence of the accompanying SF, which may also show up in broad-band NIR imaging. NIR data are not optimal for determining the sizes of the circumnuclear features, and we will come back to this issue in future papers, where H$`\alpha `$ imaging will be used to localize the SF activity in the CNR. However, the radial profiles often allow the determination of the diameter of the “rings”: typical major axis diameters are 1-2 kpc. ### 4.2 Dust Lanes The first suggestion about the link between dust lanes and shocks in the gas flow was made by Prendergast (1962). Gas can follow simple periodic orbits that do not intersect when the bar or oval distortion is not very strong. When the bar is strong, the families of periodic orbits intersect and shocks are formed (e.g. Sanders, Teuben and Van Albada 1983). Dust lanes in bars, as described by Athanassoula (1992), can be: 1. straight and parallel to the bar major axis, or 2. curved and have their concave sides towards the bar major axis. The detailed shape of the loci of offset shocks changes with the parameters of the potential (Athanassoula 1992). The simulations presented by Athanassoula show that there is a very clear sequence of shapes for the shock loci. For axial ratios corresponding to fat bars or ovals, and presumably to what could theoretically be called a weak bar, the shock lanes are curved with their concave parts towards the bar major axis. For slim (“strong”)bars, the dust lanes straighten out, their curvature decreases and their shape becomes gradually more straight. For even higher axial ratios the shock loci tend towards the bar major axis. Thus, strong bars should have straight dust lanes, while ovals or weak bars should have more curved dust lanes. Athanassoula (1992) also proposed that the strength of SF near the dust lanes depends on the strength of the bar. The dust lanes observed on the leading edges of most bars with high ellipticity (strong) are very smooth and show no Hii regions or other signs of recent SF. In the case of curved dust lanes strings of Hii regions are accompanying the dust lanes. Due to the enhanced gas mass fraction, dust is expected to be prominent in the CNRs themselves, and this, in fact, is seen clearly in most cases in the colour index images, and also often in the broad-band images, especially the HST NICMOS images. Such dust lanes are expected to be continuations of the dust lanes in the bar, and this is what is seen in our images (see Shlosman 1999 for a theoretical review). NIR colour index images are not very sensitive to changes in stellar populations, and can outline dust structure clearly (see e.g. the $`IK`$ map of M100 in Knapen et al. 1995a). We see clear and abundant observational evidence for dust lanes on several scales, but most clearly in the CNRs. Dust lanes in the bars are not well visible in general in our NIR imaging due to the lower signal to noise ratios achieved in the bar regions. In Paper II, we present optical colour index maps which outline the dust lane structure in the bars and discs of our sample galaxies more clearly, and we study the relationship between the shape of the dust lanes, the axial ratio, and the SF in the bar in more detail. ### 4.3 Colour Distributions Our colour maps have shown the presence of circumnuclear rings in most of the galaxies in our sample. These rings are generally redder than other regions in the galaxy by about 0.05– 0.1 magnitude in $`JK`$ and $`HK`$. From the colour index maps or colour profiles alone it is not possible to make meaningful statements about quantities of extinguishing dust implied by the redder colours. In fact, there are indications from both optical and NIR imaging and spectroscopy that the red colours in CNRs in galaxies like the ones studied here may be influenced by young stars, e.g. red supergiants (Knapen et al. 1995a,b; Knapen 1996; Elmegreen et al. 1997; Ryder & Knapen 1999). In Paper III, we will compare the precise location of the $`JK`$ and $`HK`$ features with those of SF regions as seen in H$`\alpha `$ emission, and try to place quantitative limits on the origins of the red light in the CNRs. Two of the host galaxies of Sy nuclei (NGC 3516 and NGC 3982), which have colours consistent with those measured by Peletier et al. (1999), also show a peculiar shape in the colour profiles, starting from a very red value and steeply decreasing until the radius of the ring, remaining constant afterwards. Peletier et al. (1999) suggested that the red $`HK`$ (or $`JK`$) colours in the cores of many Sy galaxies could be due to a significant fraction of thermal radiation from hot dust heated by the Sy nucleus. However, the two other AGN hosts in our sample (NGC 4303 and NGC 6951) do not show significantly red nuclei. ### 4.4 Radial profiles One of the main data products presented in this paper is a set of detailed radial profile fits to the NIR images. We present (Fig. 2a-l) radial profiles of ellipticity, major axis position angle, NIR colour, and $`H`$-band surface brightness. We will discuss the existence and properties of (primary and secondary) bars in Paper II, where we combine the NIR profiles presented here with $`I`$-band profiles of the complete discs, covering also the primary bar, if present. Certain aspects of the radial profiles have already been discussed in Section 3, but there is one characteristic which is worth mentioning here. Apart from NGC 3516 and NGC 3982, the radial profiles show characteristic behaviour near the radius of the circumnuclear “ring”. In seven (out of 10) galaxies, all profiles (surface brightness, colour, ellipticity and position angle) show either a change in slope, where the surface brightness profile becomes steeper, or a bump, when the colour profile shows a limited redder region caused by either dust or SF in or near the ring – see Section 4.3. These characteristics are accompanied by some change in ellipticity and position angle. In a further two galaxies (NGC 1530 and NGC 4314) only the surface brightness profile does not follow this trend, whereas in NGC 5248 the colour profile is the deviant one. Of course in NGC 2903 and NGC 3351 we cannot judge the behaviour of ellipticity/position angle and colour profiles, respectively. We postulate that the features in the surface brightness and colour profiles are due directly to a relatively young population. Firstly, an old, bulge-like, population would not be expected to give rise to strong high-resolution features such as the ones discussed here. Secondly, if the red bumps in the colour profiles were due to extinction by dust only, they should be accompanied by relative dips in the corresponding surface brightness profiles (such dust would remove light in each of the $`J,H`$ and $`K`$-bands). Instead, red bumps in the colour profiles are in all cases accompanied by bumps in the surface brightness profiles: relative excesses in emission, red in colour, that are naturally explained by excess emission produced locally by relatively young stars, probably accompanied by dust. We will come back to this issue in Paper III, where we combine the NIR data with H$`\alpha `$ imaging. ## 5 Conclusions This is the first paper in a series exploring the relationships between ring-like star-forming regions in the central few kpc regions of barred galaxies and the morphological properties of their hosts. In this paper, we present new sub-arcsec resolution NIR images obtained with the CFHT of the central regions of a sample of 12 barred galaxies with circumnuclear SF activity, and HST archival NIR images of most of the sample objects. In Paper III, we will combine these data with optical imaging of the complete hosts, as presented in Paper II. Our NIR images reveal a wealth of structure, caused by young stars and dust. In colour index maps, structures such as SF sites and dust lanes are generally apparent, even in those galaxies where the broad-band images appear featureless. The colour maps of most galaxies show evidence of an often tightly wound spiral structure. In some cases, this spiral structure is visible also in the broad-band images, and as expected due to the higher spatial resolution, more often so in the HST than in the ground-based images. Where circumnuclear structure is visible in the HST images (in about three quarters of all cases), it shows large numbers of small emitting regions, presumably SF regions, often outlined in spiral arm fragments, and accompanied by dust lanes. Circumnuclear spiral structure thus appears to be common in barred spiral galaxies with circumnuclear SF. In most of our sample galaxies, and in all those where the star-forming nuclear ring is well defined, radial profiles of surface brightness, colour, ellipticity and position angle show a characteristic bump or change in slope at the radius where the circumnuclear ring is located. Acknowledgements This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration.
warning/0004/hep-ph0004038.html	ar5iv	text	# Renormalizability and model independent description of 𝑍' signals at low energies ## I Introduction The existence of the heavy $`Z^{}`$ gauge boson is predicted by a number of grand unified theories (GUT’s) and superstring theories . The mass of this particle is expected to be of order $`m_Z^{}500`$ GeV, and therefore it cannot be produced at present day accelerators. Various strategies of searching for signals of $`Z^{}`$ as a virtual heavy state were developed and different observables convenient for its experimental detection have been introduced (see the survey and references therein). The model-dependent and model-independent $`Z^{}`$ searches at $`e^+e^{}`$ colliders are discussed (see, for instance, the report ). A popular model assumes that at low energies the $`Z^{}`$ interactions with ordinary particles of the Standard Model (SM) can be described by the effective gauge group $`\mathrm{SU}(2)_L\times \mathrm{U}(1)_Y\times \stackrel{~}{\mathrm{U}}(1)`$. An alternative choice is the gauge group $`\mathrm{SU}(2)_L\times \mathrm{SU}(2)_R\times \mathrm{U}(1)_{BL}`$ . These models are considered as the remnants of underlying theories which are not specified. The low-energy effective Lagrangians (EL) take into consideration the most general property of renormalizable theories, ensured by the decoupling theorem – the dominance of renormalizable interactions at low energies. The interactions of non-renormalizable types, being generated at high energies due to radiation corrections, are suppressed by the inverse heavy mass $`1/m_Z^{}`$. Therefore, it is possible not to consider them in leading order at lower energies. Another popular description is the introduction of the EL, considered as the sum of all effective operators with dimensions $`n4`$, constructed from the fields of light particles. The coefficients at these operators are treated as independent unknown numbers to be determined in experiments. For more details see Ref. . In general, the number of possible $`Z^{}`$ couplings is large. So, it is difficult to introduce observables allowing a unique detection of $`Z^{}`$ signals. In this regard, it is desirable either to decrease the number of the independent $`Z^{}`$ parameters on some reasonable grounds and to introduce observables most sensitive to the $`Z^{}`$ virtual states. In any case, the main idea is to find correlations between the $`Z^{}`$ couplings at low energies. A straightforward way to find the correlations is to specify the underlying theory describing interactions at energies $`\mathrm{\Lambda }_{\mathrm{GUT}}`$ and to consider running of the couplings from high to low energies $`m_W`$ by using the renormalization group (RG) equations. In this approach, each underlying theory leads to the unique values of the parameters and, hence, the corresponding correlations are model dependent ones. Another way is to specify a basis low-energy theory (for instance, the SM can be chosen) and to determine the relations between the $`Z^{}`$ parameters, following from some model independent arguments. These correlations are to be model independent. Naturally, they remain dependent on the chosen basis low-energy theory. In Refs. the method for derivation model independent correlations between the parameters of physics beyond the SM has been developed, and new observables convenient in searching for the $`Z^{}`$ boson in four-fermion processes were introduced. This approach is based on principles of the RG and the decoupling theorem . As it was argued, any virtual heavy particle can be treated as an “external field” scattering the SM particles. The vertex describing interaction with the field contains a numeric factor, which is considered as an arbitrary parameter. Actually, it is generated by the decoupling and therefore depends on the underlying model. Due to renormalizability, the scattering amplitude in the “external field” satisfies some simple relation (named RG relation), which includes the $`\beta `$ and $`\gamma `$ functions entering the RG equation. These functions have to be calculated with the light particles only, and the vertex factor. Hence, relations between different vertex factors follow. Then, they can be implemented in a number of model independent observables corresponding to the specific heavy virtual state, in particular, to the $`Z^{}`$ gauge boson . In Ref. as the low-energy basis model the minimal SM (with one scalar doublet) has been chosen. However, at present there is a few information about the scalar fields. In this regard, the theory with two scalar doublets is intensively studied . The two-Higgs-doublet model (THDM) is also known as the low-energy limit of some $`\mathrm{E}_6`$ based GUT’s, which predict the $`Z^{}`$ gauge boson. In the present paper, the results of Ref. are generalized to the THDM case. We analyse in detail both the Abelian and the so called “chiral” types of the $`Z^{}`$ couplings to light particles. As the latter type is concerned, it was derived as follows. We first assumed the most general parametrization of $`Z^{}`$ interactions with the SM fields and then derived the generator structures, compatible with the renormalizability. As it will be shown in what follows, there is an important difference between these two types of interactions. Thus, in order to derive the model independent constraints we choose the THDM as the low-energy basis theory (notice, the minimal SM is a particular case of the THDM). Then, we introduce a general parametrization of linear in $`Z^{}`$ couplings, which is independent of the specific underlying theory. As a result, the derived RG correlations are model independent ones. They hold for the THDM as well as for the minimal SM. Moreover, the existence of other heavy particles with masses $`m_im_Z^{}`$ does not affect these correlations. As it will be shown, there are two completely different sets of the $`Z^{}`$ couplings to the SM fields compatible with renormalizability. The first one describes the Abelian $`Z^{}`$, which respects the additional $`\stackrel{~}{\mathrm{U}}(1)`$ symmetry of the low energy EL. In this case the $`Z^{}`$ couplings to the axial-vector fermion currents have a universal absolute value. The second set corresponds to the chiral $`Z^{}`$, which interacts with the SM doublets, only. One has to distinguish these neutral $`Z^{}`$ gauge bosons because they are described by different operators. The content is as follows. In Sec. II the general parametrization of interactions involving the $`Z^{}`$ and the SM fields is introduced. The RG correlations between the $`Z^{}`$ couplings are derived in Sec. III. In Sec. IV they are compared with the specific values of the $`Z^{}`$ couplings in the GUT’s based on the $`\mathrm{E}_6`$ group. In Sec. V the observables convenient in detection of the $`Z^{}`$ signals are proposed. The results of our investigation are discussed in Sec. VI. ## II Parametrization of the $`Z^{}`$ couplings In the present paper we analyze the four-fermion scattering amplitudes of order $`m_Z^{}^2`$ generated by the virtual $`Z^{}`$ states. Vertices of interactions with more than one $`Z^{}`$ field contribute to the amplitudes involving several virtual $`Z^{}`$ states. The latter processes have order $`m_Z^{}^4`$ and higher. Therefore, in what follows we consider the linear in $`Z^{}`$ vertices, only. To introduce a general parametrization of the vertices involving the SM fields and being linear in the $`Z^{}`$ field, let us impose a number of natural conditions. First of all, the renormalizable type interactions are dominant at low energies $`m_W`$. The non-renormalizable interactions generated at high energies due to radiation corrections are suppressed by the inverse heavy mass $`1/m_Z^{}`$ (or by other heavier scales $`1/\mathrm{\Lambda }_i1/m_Z^{}`$) and therefore at low energies can be neglected in leading order. We assume that the $`Z^{}`$ is the only neutral vector boson with the mass $`m_Z^{}`$, and the $`Z^{}`$ gauge field enters the theory through covariant derivatives with a corresponding charge. We also assume that the $`\mathrm{SU}(2)_L\times \mathrm{U}(1)_Y`$ gauge group of the SM is a subgroup of the GUT group. In this case, a product of generators associated with the SM subgroup is a linear combination of these generators. As a consequence, the all structure constants connecting two SM gauge bosons with $`Z^{}`$ have to be zero. Hence, the interactions of gauge fields of the types $`Z^{}W^+W^{}`$, $`Z^{}ZZ`$, and other are absent at tree level. Let $`\varphi _i`$ ($`i=1,2`$) be two complex scalar doublets: $$\varphi _i=\{a_i^+,\frac{v_i+b_i+ic_i}{\sqrt{2}}\},$$ (1) where $`v_i`$ marks corresponding vacuum expectation values, $`a_i^+`$ are complex fields, and $`b_i`$, $`c_i`$ are real fields. By diagonalizing the quadratic terms of the scalar potential $`V(\varphi _1,\varphi _2)`$ one obtains the mass eigenstates: two neutral $`CP`$-even scalar particles, $`H`$ and $`h`$, the neutral $`CP`$-odd scalar particle, $`A_0`$, the Goldstone boson partner of the $`Z`$ boson, $`\chi _3`$, the charged Higgs field, $`H^\pm `$, and the Goldstone field associated with the $`W^\pm `$ boson, $`\chi ^\pm `$: $`a_1^+=`$ $`\chi ^+\mathrm{cos}\beta H^+\mathrm{sin}\beta ,`$ $`a_2^+=H^+\mathrm{cos}\beta +\chi ^+\mathrm{sin}\beta ,`$ (2) $`c_1=`$ $`\chi _3\mathrm{cos}\beta A_0\mathrm{sin}\beta ,`$ $`c_2=A_0\mathrm{cos}\beta +\chi _3\mathrm{sin}\beta ,`$ (3) $`b_1=`$ $`H\mathrm{cos}\alpha h\mathrm{sin}\alpha ,`$ $`b_2=h\mathrm{cos}\alpha +H\mathrm{sin}\alpha ,`$ (4) where $$\mathrm{tan}\beta =\frac{v_2}{v_1},$$ (5) and the angle $`\alpha `$ is determined by the explicit form of the potential $`V(\varphi _1,\varphi _2)`$. For instance, the $`CP`$-conserving potential, which has only $`CP`$-invariant minima, can be used : $`V`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}\left[\mu _i^2\varphi _i^{}\varphi _i+\lambda _i(\varphi _i^{}\varphi _i)^2\right]+\lambda _3(\text{Re}[\varphi _1^{}\varphi _2])^2`$ (7) $`+\lambda _4(\text{Im}[\varphi _1^{}\varphi _2])^2+\lambda _5(\varphi _1^{}\varphi _1)(\varphi _2^{}\varphi _2).`$ It is consistent with the absence of the tree-level flavor-changing neutral currents (FCNC’s) in the fermion sector. The corresponding value of $`\alpha `$ is $$\mathrm{tan}2\alpha =\frac{v_1v_2\left(\lambda _3+\lambda _5\right)}{\lambda _2v_2^2\lambda _1v_1^2}.$$ (8) At low energies, when all heavy states are decoupled, the $`Z^{}`$ interactions with the scalar doublets can be parametrized in a model independent way as follows : $`_\varphi `$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}\|(_\mu {\displaystyle \frac{ig}{2}}\sigma _aW_\mu ^a{\displaystyle \frac{ig^{}}{2}}Y_{\varphi _i}B_\mu `$ (10) $`{\displaystyle \frac{i\stackrel{~}{g}}{2}}\stackrel{~}{Y}_{\varphi _i}\stackrel{~}{B}_\mu \left)\varphi _i\right\|^2,`$ where $`g`$, $`g^{}`$, $`\stackrel{~}{g}`$ are the charges associated with the $`\mathrm{SU}(2)_L`$, $`\mathrm{U}(1)_Y`$, and the $`Z^{}`$ gauge groups, respectively, $`\sigma _a`$ are the Pauli matrices, $`\stackrel{~}{Y}_{\varphi _i}`$ $`=`$ $`\left(\begin{array}{cc}\stackrel{~}{Y}_{\varphi _i,1}& 0\\ 0& \stackrel{~}{Y}_{\varphi _i,2}\end{array}\right)`$ (13) is the generator corresponding to the gauge group of the $`Z^{}`$ boson, and $`Y_{\varphi _i}`$ is the $`\mathrm{U}(1)_Y`$ hypercharge. The condition $`Y_{\varphi _i}=1`$ guarantees that the vacuum is invariant with respect to the gauge group of photon. The vector bosons, $`A`$, $`Z`$, and $`Z^{}`$, are related with the symmetry eigenstates as follows: $`B`$ $``$ $`A\mathrm{cos}\theta _W(Z\mathrm{cos}\theta _0Z^{}\mathrm{sin}\theta _0)\mathrm{sin}\theta _W,`$ (14) $`W_3`$ $``$ $`A\mathrm{sin}\theta _W+(Z\mathrm{cos}\theta _0Z^{}\mathrm{sin}\theta _0)\mathrm{cos}\theta _W,`$ (15) $`\stackrel{~}{B}`$ $``$ $`Z\mathrm{sin}\theta _0+Z^{}\mathrm{cos}\theta _0,`$ (16) where $`\mathrm{tan}\theta _W=g^{}/g`$ is the adopted in the SM value of the Weinberg angle, and $`\mathrm{tan}\theta _0`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{g}m_W^2\left(\stackrel{~}{Y}_{\varphi _1,2}\mathrm{cos}^2\beta +\stackrel{~}{Y}_{\varphi _2,2}\mathrm{sin}^2\beta \right)}{g\mathrm{cos}\theta _W\left(m_Z^{}^2m_W^2/\mathrm{cos}^2\theta _W\right)}}.`$ (17) As is seen, the mixing angle $`\theta _0`$ is of order $`m_W^2/m_Z^{}^2`$. That results in the corrections of order $`m_W^2/m_Z^{}^2`$ to the interactions between the SM particles. To avoid the tree-level mixing of the $`Z`$ boson and the physical scalar field $`A_0`$ one has to impose the condition $`\stackrel{~}{Y}_{\varphi _1,2}=\stackrel{~}{Y}_{\varphi _2,2}\stackrel{~}{Y}_{\varphi ,2}`$. Now, let us parametrize the fermion-vector interactions introducing the effective low-energy Lagrangian : $`_f`$ $`=`$ $`i{\displaystyle \underset{f_L}{}}\overline{f}_L\gamma ^\mu (_\mu {\displaystyle \frac{ig}{2}}\sigma _aW_\mu ^a{\displaystyle \frac{ig^{}}{2}}B_\mu Y_{f_L}`$ (20) $`{\displaystyle \frac{i\stackrel{~}{g}}{2}}\stackrel{~}{B}_\mu \stackrel{~}{Y}_{f_L})f_L`$ $`+i{\displaystyle \underset{f_R}{}}\overline{f}_R\gamma ^\mu \left(_\mu ig^{}B_\mu Q_f{\displaystyle \frac{i\stackrel{~}{g}}{2}}\stackrel{~}{B}_\mu \stackrel{~}{Y}_{R,f}\right)f_R,`$ where the renormalizable type interactions are admitted and the summation over the all SM left-handed fermion doublets, $`f_L=\{(f_u)_L,(f_d)_L\}`$, and the right-handed singlets, $`f_R=(f_u)_R,(f_d)_R`$, is understood. $`Q_f`$ denotes the charge of $`f`$ in the positron charge units, $`\stackrel{~}{Y}_{f_L}`$ $`=`$ $`\left(\begin{array}{cc}\stackrel{~}{Y}_{L,f_u}& 0\\ 0& \stackrel{~}{Y}_{L,f_d}\end{array}\right),`$ (23) and $`Y_{f_L}`$ equals to $`1`$ for leptons and $`1/3`$ for quarks. Renormalizable interactions of fermions and scalars are described by the Yukawa Lagrangian. To avoid the existence of the tree-level FCNC’s one has to ensure that at the diagonalization of the fermion mass matrix the diagonalization of the scalar-fermion couplings is automatically fulfilled. In this case the Yukawa Lagrangian, which respects the $`\mathrm{SU}(2)_L\times \mathrm{U}(1)_Y`$ gauge group, can be written in the form: $`_{\mathrm{Yuk}}`$ $`=`$ $`\sqrt{2}{\displaystyle \underset{f_L}{}}{\displaystyle \underset{i=1}{\overset{2}{}}}\{G_{f_d,i}[\overline{f}_L\varphi _i(f_d)_R+(\overline{f}_d)_R\varphi _i^{}f_L]`$ (25) $`+G_{f_u,i}[\overline{f}_L\varphi _i^c(f_u)_R+(\overline{f}_u)_R\varphi _i^cf_L]\},`$ where $`\varphi _i^c=i\sigma _2\varphi _i^{}`$ is the charge conjugated scalar doublet, and the Cabibbo-Kobayashi-Maskawa mixing is neglected. Then, the fermion masses are $$m_f=\frac{2m_W}{g}\left(G_{f,1}\mathrm{cos}\beta +G_{f,2}\mathrm{sin}\beta \right).$$ (26) As was shown by Glashow and Weinberg , the tree-level FCNC’s mediated by Higgs bosons are absent in case when all fermions of a given electric charge couple to no more than one Higgs doublet. This restriction leads to four different models, as discussed in Ref. . In what follows, we will use the most general parametrization (25) including the models mentioned as well as other possible variations of the Yukawa sector without the tree-level FCNC’s. By using Eqs. (10), (20), and (25) it is easy to derive the Feynman rules which are collected in Appendix A. ## III RG relations In this section we consider the correlations between the parameters $`\stackrel{~}{Y}_{L,f}`$, $`\stackrel{~}{Y}_{R,f}`$, $`\stackrel{~}{Y}_{\varphi _i,1}`$, and $`\stackrel{~}{Y}_{\varphi _i,2}`$ appearing due to the renormalizability of an underlying theory. As is known, $`S`$-matrix elements are to be invariant with respect to the RG transformations, which express the independence of the location of a normalization point $`\kappa `$ in the momentum space. In a theory with different mass scales the decoupling of heavy loop contributions at the thresholds of heavy masses, $`\mathrm{\Lambda }`$, results in the important property of low energy amplitudes: the running of all functions is regulated by the loops of light particles. Therefore, the $`\beta `$ and $`\gamma `$ functions at low energies are determined by the SM particles, only. This fact is the consequence of the decoupling theorem . Actually, the decoupling results in the redefinition of the parameters of the theory at the scale $`\mathrm{\Lambda }`$ and removing the all heavy particle loop contributions proportional to $`\mathrm{ln}\kappa `$ from the RG equation : $`\lambda _a`$ $`=`$ $`\widehat{\lambda }_a+a_{\lambda _a}\mathrm{ln}{\displaystyle \frac{\widehat{\mathrm{\Lambda }}^2}{\kappa ^2}}+b_{\lambda _a}\mathrm{ln}^2{\displaystyle \frac{\widehat{\mathrm{\Lambda }}^2}{\kappa ^2}}+\mathrm{},`$ (27) $`X`$ $`=`$ $`\widehat{X}\left(1+a_X\mathrm{ln}{\displaystyle \frac{\widehat{\mathrm{\Lambda }}^2}{\kappa ^2}}+b_X\mathrm{ln}^2{\displaystyle \frac{\widehat{\mathrm{\Lambda }}^2}{\kappa ^2}}+\mathrm{}\right),`$ (28) where we use the notation $`\lambda _a`$ to refer to the charges, and $`X`$ represents the all fields and masses. Hats over quantities mark the parameters of the underlying theory. They include the loops of both the SM and the heavy particles, whereas the quantities without hats are calculated assuming that no heavy particles are excited inside loops. The matching between the both sets of parameters ($`\lambda _a`$, $`X`$ and $`\widehat{\lambda }_a`$, $`\widehat{X}`$) is chosen to be done at the normalization point $`\kappa \mathrm{\Lambda }`$, $$\lambda _a_{\kappa =\mathrm{\Lambda }}=\widehat{\lambda }_a_{\kappa =\mathrm{\Lambda }},X_{\kappa =\mathrm{\Lambda }}=\widehat{X}_{\kappa =\mathrm{\Lambda }}.$$ (29) Since the sets of parameters $`\lambda _a`$, $`X`$ and $`\widehat{\lambda }_a`$, $`\widehat{X}`$ differ at one-loop level, it is possible to substitute one set by another. As is shown in Ref. , the redefinition of fields and charges (27) allows one to eliminate the one-loop mixing between heavy and light virtual states. Therefore, virtual states of heavy particles can be treated as the “external fields” scattering SM particles. The renormalizability of the underlying theory leads to some relations for vertices describing this scattering, called the RG relations. Let us consider the four-fermion amplitudes caused by the $`Z^{}`$ boson exchange. In the lower order in ratio $`m_W^2/m_Z^{}^2`$ the process $`\overline{f}_1f_1Z_{}^{}{}_{}{}^{}\overline{f}_2f_2`$ can be presented as scattering of the initial, $`f_1`$, and the final, $`f_2`$, fermions in the “external field” $`1/m_Z^{}`$ with the corresponding vertex factors $`\mathrm{\Gamma }_{f_1Z^{}}`$, $`\mathrm{\Gamma }_{f_2Z^{}}`$. The quantity $`\mathrm{\Gamma }_{fZ^{}}`$ contains no contributions of heavy particle loops. Thus, it can be computed as a linear combination of the parameters $`\stackrel{~}{Y}_{L,f}`$, $`\stackrel{~}{Y}_{R,f}`$, $`\stackrel{~}{Y}_{\varphi _i,1}`$, and $`\stackrel{~}{Y}_{\varphi _i,2}`$. The RG invariance of the vertex leads to equation $$𝒟\left(\overline{f}\mathrm{\Gamma }_{fZ^{}}f\frac{1}{m_Z^{}}\right)=0,$$ (30) where the effective low-energy RG operator is defined as follows: $`𝒟`$ $``$ $`{\displaystyle \frac{}{\mathrm{ln}\kappa }}+{\displaystyle \underset{a}{}}\beta _a{\displaystyle \frac{}{\lambda _a}}{\displaystyle \underset{X}{}}\gamma _X{\displaystyle \frac{}{\mathrm{ln}X}},`$ (31) $`\beta _a`$ $`=`$ $`{\displaystyle \frac{d\lambda _a}{d\mathrm{ln}\kappa }},\gamma _X={\displaystyle \frac{d\mathrm{ln}X}{d\mathrm{ln}\kappa }},`$ (32) and the coefficient functions $`\beta _a`$ and $`\gamma _X`$ are computed taking into account the loops of light particles. Relation (30) ensures that, as a consequence of renormalizability, the mathematical structure of the leading logarithmic terms of the vertices, calculated in one- and higher-loop approximations, coincides with that of the tree-level structures. The standard usage of Eq. (30) is to improve scattering amplitudes calculated in a fixed order of perturbation theory. In contrast, in what follows we will apply Eq. (30) to obtain an algebraic relation between the parameters $`\stackrel{~}{Y}_{L,f}`$, $`\stackrel{~}{Y}_{R,f}`$, $`\stackrel{~}{Y}_{\varphi _i,1}`$, $`\stackrel{~}{Y}_{\varphi _i,2}`$, which are to be considered as arbitrary numbers, since the underlying theory is not specified. Let us explain the idea in more detail. In case when the underlying theory is specified ($`\stackrel{~}{Y}_{L,f}`$, $`\stackrel{~}{Y}_{R,f}`$, $`\stackrel{~}{Y}_{\varphi _i,1}`$, $`\stackrel{~}{Y}_{\varphi _i,2}`$ have to be computed as discussed before), and the $`\beta `$ and $`\gamma `$ functions as well as the $`S`$-matrix elements are calculated in a fixed order of perturbation theory, Eq. (30) is just the identity. If the underlying theory is not specified, whereas the $`\beta `$, $`\gamma `$ functions and $`S`$-matrix elements are computed in a fixed order of perturbation theory, equality (30) may serve to correlate the unknown parameters $`\stackrel{~}{Y}`$. In case of the four-fermion processes mediated by the gauge $`Z^{}`$ boson, the number of independent $`\beta `$ functions is less than the number of RG equations. Therefore, the non-trivial system of equations correlating the originally independent parameters occurs. The one-loop RG relation for the fermion-$`Z^{}`$ vertex is $$\overline{f}\frac{\mathrm{\Gamma }_{fZ^{}}^{(1)}}{\mathrm{ln}\kappa }f\frac{1}{m_Z^{}}+𝒟^{(1)}\left(\overline{f}\mathrm{\Gamma }_{fZ^{}}^{(0)}f\frac{1}{m_Z^{}}\right)=0,$$ (33) where $`\mathrm{\Gamma }_{fZ^{}}^{(0)}`$ and $`\mathrm{\Gamma }_{fZ^{}}^{(1)}`$ denote the tree-level and the one-loop level contributions to the fermion-$`Z^{}`$ vertex, and $`𝒟^{(1)}`$ is the one-loop level part of the RG operator, $$𝒟^{(1)}\underset{a}{}\beta _a^{(1)}\frac{}{\lambda _a}\underset{X}{}\gamma _X^{(1)}\frac{}{\mathrm{ln}X}.$$ (34) As it follows from Eq. (33), only the divergent parts of the one-loop vertices $`\mathrm{\Gamma }_{fZ^{}}^{(1)}`$ are to be calculated. The corresponding diagrams are shown in Fig. 1. The following expressions for the right-handed and the left-handed fermions, respectively, have been obtained, $`{\displaystyle \frac{\mathrm{\Gamma }_{f_RZ^{}}^\mu }{\mathrm{ln}\kappa }}`$ $`=`$ $`{\displaystyle \frac{\gamma ^\mu }{8\pi ^2}}\{g^2Q_f^2\stackrel{~}{Y}_{R,f}\mathrm{tan}^2\theta _W+{\displaystyle \frac{4}{3}}g_{s,f}^2\stackrel{~}{Y}_{R,f}`$ (38) $`+G_{f,1}^2\left[2T_f^3\left(\stackrel{~}{Y}_{\varphi ,2}+\stackrel{~}{Y}_{\varphi _1,1}\right)+\stackrel{~}{Y}_{L,f}+\stackrel{~}{Y}_{L,f^{}}\right]`$ $`+G_{f,2}^2\left[2T_f^3\left(\stackrel{~}{Y}_{\varphi ,2}+\stackrel{~}{Y}_{\varphi _2,1}\right)+\stackrel{~}{Y}_{L,f}+\stackrel{~}{Y}_{L,f^{}}\right]`$ $`+O\left({\displaystyle \frac{m_W^2}{m_Z^{}^2}}\right)\},`$ $`{\displaystyle \frac{\mathrm{\Gamma }_{f_LZ^{}}^\mu }{\mathrm{ln}\kappa }}`$ $`=`$ $`{\displaystyle \frac{\gamma ^\mu }{8\pi ^2}}\{{\displaystyle \frac{g^2}{2}}\stackrel{~}{Y}_{L,f^{}}+{\displaystyle \frac{4}{3}}g_{s,f}^2\stackrel{~}{Y}_{L,f}`$ (44) $`+g^2\stackrel{~}{Y}_{L,f}\left[{\displaystyle \frac{1}{4\mathrm{cos}^2\theta _W}}+\left(Q_f^2\left\|Q_f\right\|\right)\mathrm{tan}^2\theta _W\right]`$ $`+\left(G_{f,1}^2+G_{f,2}^2\right)\left(\stackrel{~}{Y}_{R,f}2T_f^3\stackrel{~}{Y}_{\varphi ,2}\right)`$ $`+G_{f^{},1}^2\left(2T_f^3\stackrel{~}{Y}_{\varphi _1,1}+\stackrel{~}{Y}_{R,f^{}}\right)`$ $`+G_{f^{},2}^2\left(2T_f^3\stackrel{~}{Y}_{\varphi _2,1}+\stackrel{~}{Y}_{R,f^{}}\right)`$ $`+O\left({\displaystyle \frac{m_W^2}{m_Z^{}^2}}\right)\},`$ where $`f`$ and $`f^{}`$ are the partners of a $`\mathrm{SU}(2)_L`$ fermion doublet (namely, $`l^{}=\nu _l`$, $`\nu _l^{}=l`$, $`q_u^{}=q_d`$, and $`q_d^{}=q_u`$), $`T_f^3`$ is the third component of the weak isospin, and $`g_{s,f}`$ is the QCD charge for quarks, and zero for leptons. The fermion anomalous dimensions can be calculated by using the diagrams of Fig. 2: $`\gamma _{f_R}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}[g^2Q_f^2\mathrm{tan}^2\theta _W+{\displaystyle \frac{4}{3}}g_{s,f}^2+2(G_{f,1}^2+G_{f,2}^2)`$ (46) $`+O\left({\displaystyle \frac{m_W^2}{m_Z^{}^2}}\right)],`$ $`\gamma _{f_L}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}[g^2(Q_f^2\|Q_f\left\|\right)\mathrm{tan}^2\theta _W+{\displaystyle \frac{4}{3}}g_{s,f}^2+{\displaystyle \frac{g^2}{2}}`$ (49) $`+{\displaystyle \frac{g^2}{4\mathrm{cos}^2\theta _W}}+G_{f,1}^2+G_{f,2}^2+G_{f^{},1}^2+G_{f^{},2}^2`$ $`+O\left({\displaystyle \frac{m_W^2}{m_Z^{}^2}}\right)].`$ RG relations (33) considered in a lower order in $`m_W^2/m_Z^{}^2`$ lead to the equations for the parameters $`\stackrel{~}{Y}_{L,f}`$, $`\stackrel{~}{Y}_{R,f}`$, $`\stackrel{~}{Y}_{\varphi _i,1}`$, and $`\stackrel{~}{Y}_{\varphi _i,2}`$: $`4\pi ^2\stackrel{~}{Y}_{R,f}\left({\displaystyle \frac{\beta _{\stackrel{~}{g}}^{(1)}}{\stackrel{~}{g}^2}}+\gamma _{m_Z^{}^2}^{(1)}\right)=`$ (56) $`G_{f,1}^2\left[2T_f^3\left(\stackrel{~}{Y}_{\varphi ,2}+\stackrel{~}{Y}_{\varphi _1,1}\right)+\stackrel{~}{Y}_{L,f}+\stackrel{~}{Y}_{L,f^{}}2\stackrel{~}{Y}_{R,f}\right]`$ $`G_{f,2}^2\left[2T_f^3\left(\stackrel{~}{Y}_{\varphi ,2}+\stackrel{~}{Y}_{\varphi _2,1}\right)+\stackrel{~}{Y}_{L,f}+\stackrel{~}{Y}_{L,f^{}}2\stackrel{~}{Y}_{R,f}\right],`$ $`4\pi ^2\stackrel{~}{Y}_{L,f}\left({\displaystyle \frac{\beta _{\stackrel{~}{g}}^{(1)}}{\stackrel{~}{g}^2}}+\gamma _{m_Z^{}^2}^{(1)}\right)={\displaystyle \frac{g^2}{2}}\left(\stackrel{~}{Y}_{L,f}\stackrel{~}{Y}_{L,f^{}}\right)`$ $`+\left(G_{f,1}^2+G_{f,2}^2\right)\left(2T_f^3\stackrel{~}{Y}_{\varphi ,2}+\stackrel{~}{Y}_{L,f}\stackrel{~}{Y}_{R,f}\right)`$ $`G_{f^{},1}^2\left(2T_f^3\stackrel{~}{Y}_{\varphi _1,1}\stackrel{~}{Y}_{L,f}+\stackrel{~}{Y}_{R,f^{}}\right)`$ $`G_{f^{},2}^2\left(2T_f^3\stackrel{~}{Y}_{\varphi _2,1}\stackrel{~}{Y}_{L,f}+\stackrel{~}{Y}_{R,f^{}}\right).`$ One has to derive two sets of relations, which ensure the compatibility of Eqs. (56). The first one is $`\stackrel{~}{Y}_{\varphi _2,1}=\stackrel{~}{Y}_{\varphi _1,1}=\stackrel{~}{Y}_{\varphi ,2}\stackrel{~}{Y}_\varphi ,`$ (57) $`\stackrel{~}{Y}_{L,f}+\stackrel{~}{Y}_{L,f^{}}=0,\stackrel{~}{Y}_{R,f}=0.`$ (58) It describes the $`Z^{}`$ boson analogous to the third component of the $`\mathrm{SU}(2)_L`$ gauge field. The characteristic features of these interactions are the zero traces of generators and the absence of couplings to the right-handed singlets. In what follows, we shall call this type of interaction the “chiral” $`Z^{}`$. The second set, $`\stackrel{~}{Y}_{\varphi _1,1}=\stackrel{~}{Y}_{\varphi _2,1}=\stackrel{~}{Y}_{\varphi ,2}\stackrel{~}{Y}_\varphi ,`$ (59) $`\stackrel{~}{Y}_{L,f}=\stackrel{~}{Y}_{L,f^{}},\stackrel{~}{Y}_{R,f}=\stackrel{~}{Y}_{L,f}+2T_f^3\stackrel{~}{Y}_\varphi ,`$ (60) corresponds to the Abelian $`Z^{}`$ boson. In this case the SM Lagrangian appears to be invariant with respect to the $`\stackrel{~}{\mathrm{U}}(1)`$ group associated with the $`Z^{}`$. The first and the second relations in Eqs. (59) mean that appropriate generators are proportional to the unit matrix, whereas the third relation ensures the Yukawa terms to be invariant with respect to the $`\stackrel{~}{\mathrm{U}}(1)`$ transformations. Introducing the $`Z^{}`$ couplings to the vector and the axial-vector fermion currents, $`v_Z^{}^f(\stackrel{~}{Y}_{L,f}+\stackrel{~}{Y}_{R,f})/2`$, $`a_Z^{}^f(\stackrel{~}{Y}_{R,f}\stackrel{~}{Y}_{L,f})/2`$, one can rewrite the second and the third of Eqs. (59) in the following form: $$v_Z^{}^fa_Z^{}^f=v_Z^{}^f^{}a_Z^{}^f^{},a_Z^{}^f=T_f^3\stackrel{~}{Y}_\varphi .$$ (61) As is seen, the couplings of the Abelian $`Z^{}`$ to the axial-vector fermion currents have a universal absolute value proportional to the $`Z^{}`$ coupling to the scalar doublets. The solutions derived are the same as in case of the minimal SM considered in Ref. . Notice that both of correlations (57) and (59) lead to the same $`Z^{}`$ couplings to each of the scalar doublets. Notice, in case of the Abelian $`Z^{}`$ boson the correlations (59),(61) can be derived on related but formally different grounds. The point is that the renormalizability and gauge invariance of interactions are closely connected. Therefore, the requirement of renormalizability can be substituted by the requirement of gauge invariance of the effective low-energy Lagrangian. In general, the EL respects by construction various \[and, in particular, $`\stackrel{~}{U}(1)`$\] symmetries. But if non-renormalizable interactions are admitted, no relations between the arbitrary parameters can be found. If only the renormalizable interactions are taken into account, as in Eq. (20), some correlations appear. In fact, to obtain formulae (59),(61) it is sufficient to require the $`\stackrel{~}{U}(1)`$ gauge invariance of the Yukawa terms. Note also that the correlations in Eq. (61) are the same as in the SM for the specific values of the hypercharges $`Y_f`$ and $`Y_\varphi `$ corresponding to the $`U(1)_Y`$ gauge transformations of fermion and scalar fields. On the other hand, we did not find any symmetry requirement describing the all possible relations following from Eq. (57). Therefore, the renormalizability requirement looks as more general one. ## IV RG correlations and the $`Z^{}`$ in $`\mathrm{E}_6`$ based models Over last decades the GUT’s based on the $`\mathrm{E}_6`$ gauge group are intensively studied. They predict the Abelian $`Z^{}`$ boson with the mass $`m_Z^{}m_W`$. Since the low-energy limit of the $`\mathrm{E}_6`$ GUT’s is the THDM considered, it is of interest to check whether relations (61) hold for the specific values of the $`Z^{}`$ couplings in these models. There are different schemes of the $`\mathrm{E}_6`$-symmetry breaking. One of them is $`\mathrm{E}_6`$ $``$ $`\mathrm{SO}(10)\times \mathrm{U}(1)_\psi ,`$ (62) $`\mathrm{SO}(10)`$ $``$ $`\mathrm{SU}(3)_c\times \mathrm{SU}(2)_L\times \mathrm{SU}(2)_R\times `$ (64) $`\times \mathrm{U}(1)_{BL}.`$ This leads to the so called left-right (LR) model. Another scheme, $$\mathrm{E}_6\mathrm{SO}(10)\times \mathrm{U}(1)_\psi \mathrm{SU}(5)\times \mathrm{U}(1)_\chi \times \mathrm{U}(1)_\psi ,$$ (65) predicts the Abelian $`Z^{}`$, which is a linear combination of the neutral vector bosons $`\psi `$ and $`\chi `$, $$Z^{}=\chi \mathrm{cos}\stackrel{~}{\beta }+\psi \mathrm{sin}\stackrel{~}{\beta },$$ (66) where $`\stackrel{~}{\beta }`$ is the mixing angle. In Table I (see Ref. ) we show the $`Z^{}`$ couplings to the SM fermions in models mentioned (notice, the sign of axial-vector couplings in Ref. is opposite to the sign of $`a_Z^{}^f`$). At first glance, some of the couplings in Table I are inconsistent with relations (61). This requires to be discussed in more detail. First of all, let us consider the $`Z^{}`$ couplings to neutrinos. It is usually supposed in theories based on the $`\mathrm{E}_6`$ group that the Yukawa terms responsible for generation of the Dirac masses of neutrinos must be set to zero . Therefore, there are no RG relations for the $`Z^{}`$ interactions with the neutrino axial-vector currents, because the terms proportional to $`G_{\nu ,i}`$ vanish in Eq. (56). In this case the couplings $`a_Z^{}^\nu `$ given in Table I are not restricted by relations (61). Now, let us discuss the $`Z^{}`$ couplings to charged leptons and quarks. The values of the couplings satisfy relations (61) in case of the LR model. As for models described by the $`\mathrm{E}_6`$ breaking scheme (65), two possibilities of choosing $`\stackrel{~}{\beta }`$ are of interest. First is if the $`\psi `$ boson is much heavier than the $`\chi `$ field. In general, this is a natural condition, since the fields $`\psi `$ and $`\chi `$ arise at different energy scales. As a consequence, the field $`\psi `$ is decoupled, and the mixing angle $`\stackrel{~}{\beta }`$ is small ($`\stackrel{~}{\beta }1`$). In this case RG relations (61) hold in lower order in $`\stackrel{~}{\beta }`$ for the $`Z^{}`$ couplings to quarks and charged leptons. The second possibility is if the masses of $`\chi `$ and $`\psi `$ are of the same order. This means the tuning of the vacuum expectation values generating the vector boson masses. This case cannot be treated straightforwardly on the basis of relations (61) since the mixed states of the $`Z^{}`$ bosons have to be included into consideration explicitly. Although our approach is applicable in this case, it requires additional investigation. Moreover, the $`Z^{}`$ mixed states cause some different exchange amplitudes, which have to be incorporated into low-energy observables. In what follows, we will not discuss the case of two $`Z^{}`$ bosons having masses of the same order. ## V Observables Now, let us introduce the observables convenient for detection of the $`Z^{}`$ in electron-positron annihilation into fermion pairs $`e^+e^{}V^{}\overline{f}f`$ ($`fe,t`$). The center-of-mass energy is taken in the range $`\sqrt{s}500`$ GeV. Consider the case of non-polarized initial and final fermions. Since the $`t`$ quark is not considered, other fermions at these energies can be treated as massless particles, $`m_f0`$. In this approximation the left-handed and the right-handed fermions can be substituted by the helicity states, which will be marked as $`\lambda `$ and $`\xi `$ for the incoming electron and the outgoing fermion, respectively ($`\lambda ,\xi =L,R`$). Let $`𝒜_V`$ be the Born amplitude of the process $`e^+e^{}V^{}\overline{f}f`$ ($`fe,t`$) with the virtual $`V`$-boson state in the $`s`$ channel ($`V=A,Z,Z^{}`$). The $`Z^{}`$ boson existence leads to the deviation of order $`m_Z^{}^2`$ of the cross section from its SM value. In general, the tree-level deviations originate from two types of contributions. The first is caused by the $`Z`$-$`Z^{}`$ mixing. Using the results of Sec. III the mixing angle $`\theta _0`$ \[see Eq. (17)\] can be calculated as follows, $`\theta _0`$ $``$ $`{\displaystyle \frac{\stackrel{~}{g}m_W^2\stackrel{~}{Y}_\varphi }{g\mathrm{cos}\theta _Wm_Z^{}^2}}.`$ (67) Because of the mixing there are corrections of order $`\theta _0m_Z^{}^2`$ to the vertex describing interaction of $`Z`$ boson and fermions. Hence, the amplitude $`𝒜_Z(\theta _0)`$ deviates from its SM value $`𝒜_Z(\theta _0=0)`$. The second type describes the interference between the SM amplitude, $`𝒜_{\mathrm{SM}}`$, and the $`Z^{}`$ exchange amplitude, $`𝒜_Z^{}`$. Thus, for the process $`e^+e^{}\overline{f}f`$ the deviation of the cross section is $$\mathrm{\Delta }\frac{d\sigma _f}{d\mathrm{\Omega }}=\frac{d\sigma _f}{d\mathrm{\Omega }}\frac{d\sigma _{f,\mathrm{SM}}}{d\mathrm{\Omega }}=\frac{\text{Re}\left[𝒜_{\mathrm{SM}}^{}\mathrm{\Delta }𝒜\right]}{32\pi s}+O\left(\frac{s^2}{m_Z^{}^4}\right),$$ (68) where $$𝒜_{\mathrm{SM}}=𝒜_A+𝒜_Z\|_{\theta _0=0},\mathrm{\Delta }𝒜=𝒜_Z^{}+\left(\frac{d𝒜_Z}{d\theta _0}\right)_{\theta _0=0}\theta _0.$$ (69) The quantity $`\mathrm{\Delta }d\sigma /d\mathrm{\Omega }`$ can be calculated in the form $$\mathrm{\Delta }\frac{d\sigma _f}{d\mathrm{\Omega }}=\underset{\lambda ,\xi =L,R}{}\left[_{\lambda \xi }^{ef}(s)+_{\lambda \xi }^{ef}(s)\right]\left(z+P_\lambda P_\xi \right)^2,$$ (70) where $`P_L=1`$, $`P_R=1`$, $`z\mathrm{cos}\theta `$ ($`\theta `$ is the angle between the incoming electron and the outgoing fermion), $`_{\lambda \xi }^{ef}`$ denotes the $`Z`$-$`Z^{}`$ interference term, and $`_{\lambda \xi }^{ef}`$ accounts of the contributions from the $`Z`$-$`Z^{}`$ mixing: $`_{\lambda \xi }^{ef}`$ $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}\stackrel{~}{g}^2T_f^3N_f}{4\pi m_Z^{}^2}}\stackrel{~}{Y}_{\lambda ,e}\stackrel{~}{Y}_{\xi ,f}[\|Q_f\|`$ (72) $`+\chi (s)(P_\lambda \epsilon )(P_\xi 1+\|Q_f\|\|Q_f\|\epsilon )],`$ $`_{\lambda \xi }^{ef}`$ $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}g\stackrel{~}{g}T_f^3N_f\theta _0}{4\pi \mathrm{cos}\theta _W(sm_Z^2)}}[\stackrel{~}{Y}_{\xi ,f}(\delta _{\lambda ,L}2\mathrm{sin}^2\theta _W)`$ (75) $`+2T_f^3\stackrel{~}{Y}_{\lambda ,e}(2\|Q_f\|\mathrm{sin}^2\theta _W\delta _{\xi ,L})\left]\right[\|Q_f\|`$ $`+\chi (s)(P_\lambda \epsilon )(P_\xi 1+\|Q_f\|\|Q_f\|\epsilon )],`$ where $`\alpha _{\mathrm{em}}`$ is the fine structure constant, $`N_f=3`$ for quarks and $`N_f=1`$ for leptons, $`\epsilon 14\mathrm{sin}^2\theta _W0.08`$, $`\chi ^1(s)=16\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W(1m_Z^2/s)`$, and $`\delta _{\lambda ,\xi }`$ is the Kronecker delta. The leading contribution comes from the $`Z`$-$`Z^{}`$ interference term $`_{\lambda \xi }^{ef}`$, whereas the $`Z`$-$`Z^{}`$ mixing terms are suppressed by the additional factor $`m_Z^2/s`$. At energies $`\sqrt{s}500`$ GeV $`_{\lambda \xi }^{ef}_{\lambda \xi }^{ef}`$. To take into consideration the correlations (57) or (59) let us introduce the function $`\sigma _f(z)`$ defined as the difference of cross sections integrated in a suitable range of $`\mathrm{cos}\theta `$ : $$\sigma _f(z)_z^1\frac{d\sigma _f}{dz}𝑑z_1^z\frac{d\sigma _f}{dz}𝑑z.$$ (76) The conventionally used observables – the total cross section $`\sigma _{f,T}`$ and the forward-backward asymmetry $`A_{f,FB}`$ – can be obtained by a special choice of $`z`$ \[$`\sigma _{f,T}=\sigma _f(1)`$, $`A_{f,FB}=\sigma _f(0)/\sigma _{f,T}`$\]. One can express $`\sigma _f(z)`$ in terms of $`\sigma _{f,T}`$ and $`A_{f,FB}`$: $$\sigma _f(z)=\sigma _{f,T}\left[A_{f,FB}\left(1z^2\right)\frac{1}{4}z\left(3+z^2\right)\right].$$ (77) Then, the deviation $`\mathrm{\Delta }\sigma _f(z)\sigma _f(z)\sigma _{f,\mathrm{SM}}(z)`$ can be written in the form: $`\mathrm{\Delta }\sigma _f(z)`$ $`=`$ $`4\pi {\displaystyle \underset{\lambda ,\xi }{}}\left[_{\lambda \xi }^{ef}(s)+_{\lambda \xi }^{ef}(s)\right]`$ (79) $`\times \left(P_\lambda P_\xi zz^2P_\lambda P_\xi {\displaystyle \frac{z^3}{3}}\right).`$ Let us compare the observable $`\mathrm{\Delta }\sigma _f(z)`$ with the differential cross section (70). As is seen, the polynomial in the polar angle $`z`$ in Eq. (70) is replaced by the function of the boundary angle $`z`$ in Eq. (79). The overall factor $`4\pi `$ appears due to the angular integration. In what follows, we consider the observable (79) taking into account correlations (57) and (59). ### A Chiral $`Z^{}`$ The case of the chiral $`Z^{}`$ corresponds to correlations (57). Because of absence of the $`Z^{}`$ couplings to right-handed fermions the leading contribution to $`\mathrm{\Delta }\sigma _f(z)`$ is proportional to the same polynomial in $`z`$ for any outgoing fermion $`f`$: $`\mathrm{\Delta }\sigma _f(z)`$ $``$ $`4\pi _{LL}^{ef}(s)\left(1zz^2{\displaystyle \frac{z^3}{3}}\right)`$ (80) $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}\stackrel{~}{g}^2T_f^3N_f}{m_Z^{}^2}}\stackrel{~}{Y}_{L,e}\stackrel{~}{Y}_{L,f}\left(1zz^2{\displaystyle \frac{z^3}{3}}\right)`$ (82) $`\times \left\{\left[\|Q_f\|+2\chi (s)\|Q_f\|\chi (s)\right]+O\left(\epsilon \right)\right\}.`$ Therefore, the differential cross section is completely determined by the total one: $`\mathrm{\Delta }\sigma _f(z)`$ $`=`$ $`\mathrm{\Delta }\sigma _{f,T}[{\displaystyle \frac{3}{4}}(1zz^2{\displaystyle \frac{z^3}{3}})`$ (84) $`+O(\epsilon ,m_Z^2s^1)].`$ Comparing the observables for fermions of the same $`\mathrm{SU}(2)_L`$ isodoublet, $`\{f_u,f_d\}`$, it is possible to derive the correlation: $`\mathrm{\Delta }\sigma _{f_u}(z)`$ $`=`$ $`\mathrm{\Delta }\sigma _{f_d}(z)\left[{\displaystyle \frac{\|Q_{f_u}\|+1}{\|Q_{f_d}\|+1}}+O(\epsilon ,m_Z^2s^1)\right].`$ (85) Hence, the ratio $`\mathrm{\Delta }\sigma _{f_u}(z)/\mathrm{\Delta }\sigma _{f_d}(z)`$ is independent of $`z`$. It equals to 5/4 for quarks and 1/2 for leptons in lower order in $`\epsilon `$, $`m_Z^2s^1`$. So, the values of the observables in the $`\mathrm{\Delta }\sigma _{f_u}(z)`$$`\mathrm{\Delta }\sigma _{f_d}(z)`$ plane are at the same curve (straight line in the approximation used) for any $`z`$ specified. It also follows from Eq. (80) that there is a value $`z=z^{}`$ when $`\mathrm{\Delta }\sigma (z^{})=0`$. As one can check, $`z^{}=2^{2/3}1`$. Notice, the observable $`\mathrm{\Delta }\sigma (z^{})`$ is just the variable $`\mathrm{\Delta }\sigma _{}`$ proposed in Ref. . This quantity is completely insensitive to the chiral $`Z^{}`$ signals. On the other hand, the deviation of the total cross section, $`\mathrm{\Delta }\sigma _T`$, is more sensitive to signals of the chiral $`Z^{}`$, since the maximum of the polynomial $`1zz^2z^3/3`$ is at $`z=1`$. ### B Abelian $`Z^{}`$ The Abelian $`Z^{}`$ beyond the minimal SM was considered recently in Ref. , where sign-definite observables convenient for detection of the Abelian $`Z^{}`$ have been introduced. RG correlations (59) in Sec. III coincide with that of Ref. . Therefore, the observables for Abelian $`Z^{}`$ beyond the THDM are to be the same as in case of the minimal SM. In case of the chiral $`Z^{}`$ the RG correlations (57) suppress amplitudes corresponding to the processes with right-handed fermions. This is not the case for the Abelian $`Z^{}`$. However, one can switch off some contributions to observable (79) by specifying the kinematic parameter $`z`$. In what follows, it will be convenient to use the $`Z^{}`$ couplings to vector and axial-vector fermion currents \[$`v_Z^{}^f(\stackrel{~}{Y}_{L,f}+\stackrel{~}{Y}_{R,f})/2`$, $`a_Z^{}^f(\stackrel{~}{Y}_{R,f}\stackrel{~}{Y}_{L,f})/2`$\]. Because of correlations (61) the absolute value of the axial-vector couplings is universal for the all types of SM fermions, $`a_Z^{}\stackrel{~}{Y}_\varphi `$. So, the observable $`\mathrm{\Delta }\sigma _f(z)`$ has the form $`\mathrm{\Delta }\sigma _f(z)`$ $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}\stackrel{~}{g}^2}{m_Z^{}^2}}[_0^f(z,s)a_Z^{}^2+_1^f(z,s)v_Z^{}^ev_Z^{}^f`$ (87) $`+_2^f(z,s)a_Z^{}v_Z^{}^f+_3^f(z,s)v_Z^{}^ea_Z^{}].`$ As it was argued in Ref. , one is able to choose the value of $`z=z^{}`$, which switches off the leading contributions to the leptonic factors $`_1^l`$, $`_2^l`$, and the factor $`_3^f`$. The appropriate value of $`z^{}`$ can be found from the equation $`\chi (s)\left(1z^2\right)\left(z^{}+{\displaystyle \frac{z^3}{3}}\right)\left[1+\chi (s)\epsilon ^2\right]`$ $`=`$ $`0.`$ (88) The solution $`z^{}(s)`$ is shown in Fig. 3. This switches off the factor at $`v_Z^{}^ev_Z^{}^l`$. As is seen, $`z^{}`$ decreases from 0.317 at $`\sqrt{s}=500`$ GeV to 0.313 at $`\sqrt{s}=700`$ GeV. In what follows the value of $`\sqrt{s}`$ is taken to be 500 GeV, because $`z^{}`$ and $`\mathrm{\Delta }\sigma (z)`$ depend on the center-of-mass energy through the small quantity $`m_Z^2/s`$ (such contributions are of order 3%). With the above discussed choice of $`z^{}`$ made, one can introduce the sign definite observable $`\mathrm{\Delta }\sigma _l(z^{})`$: $`\mathrm{\Delta }\sigma _l(z^{})`$ $`=`$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}\stackrel{~}{g}^2}{m_Z^{}^2}}_0^l(z^{},s)a_Z^{}^2`$ (89) $`=`$ $`0.10{\displaystyle \frac{\alpha _{\mathrm{em}}\stackrel{~}{g}^2\stackrel{~}{Y}_\varphi ^2}{m_Z^{}^2}}\left[1+O\left(0.04\right)\right]<0.`$ (90) Notice, the value of $`\mathrm{\Delta }\sigma _l(z^{})`$ is universal for the all types of SM charged leptons. There are also sign definite observables for the quarks of the same generation: $$\mathrm{\Delta }\sigma _q(z^{})\mathrm{\Delta }\sigma _{q_u}+0.5\mathrm{\Delta }\sigma _{q_d}2.45\mathrm{\Delta }\sigma _l\left(z^{}\right)<0.$$ (91) Hence one can conclude that the values of $`\mathrm{\Delta }\sigma _{q_u}(z^{})`$ and $`\mathrm{\Delta }\sigma _{q_d}(z^{})`$ in the $`\mathrm{\Delta }\sigma _{q_u}(z^{})`$$`\mathrm{\Delta }\sigma _{q_d}(z^{})`$ plane have to be at the line crossing the axes at the points $`\mathrm{\Delta }\sigma _{q_u}(z^{})=2.45\mathrm{\Delta }\sigma _l(z^{})`$ and $`\mathrm{\Delta }\sigma _{q_d}(z^{})=4.9\mathrm{\Delta }\sigma _l(z^{})`$, respectively. Signals of the Abelian and the chiral $`Z^{}`$ are compared in Figs. 4-5. Suppose for a moment that experiments give the non-zero values of leptonic observables $`\mathrm{\Delta }\sigma _l(z^{})`$ ($`l=\mu ,\tau `$). If they correspond to the Abelian $`Z^{}`$, either of the observables has to be the same negative number. Let one also know the values of the neutrino observables $`\mathrm{\Delta }\sigma _\nu (z^{})`$ ($`\nu =\nu _\mu ,\nu _\tau `$). In case of the chiral $`Z^{}`$ the corresponding point in Fig. 4 has to be at the straight line shown (with the accuracy of the approximation). For the Abelian $`Z^{}`$ the shaded region as a whole is available. Now, let us consider observables for the quarks of the same generation (see Fig. 5). If the value of the leptonic observable $`\mathrm{\Delta }\sigma _l(z^{})`$ is measured, one has to expect that the experimental points will be located at one of two possible curves corresponding either to the chiral or to the Abelian $`Z^{}`$. The shaded range represents signals of the Abelian $`Z^{}`$ for the all possible values of the leptonic observable. So, by measuring the observables $`\mathrm{\Delta }\sigma _f(z^{})`$ for fermions of the same $`\mathrm{SU}(2)_L`$ isodoublet, one is able to distinguish the Abelian and the chiral $`Z^{}`$ couplings. ## VI Discussion In the present paper the method of RG relations , developed originally for the minimal SM, is extended to searching for signals of the heavy $`Z^{}`$ gauge boson beyond the THDM. General conditions when our consideration is applicable are the following. 1) The mechanism generating the heavy particle masses is not specified, and the $`Z^{}`$ mass is treated as an arbitrary parameter. 2) The light particle masses are generated in a standard way via the non-zero vacuum values of the scalar fields of the low-energy basis theory. Interactions of light particles with heavy scalar fields, which are responsible for $`m_Z^{}`$, are excluded at tree level. The radiation corrections to the masses due to heavy particle loops are suppressed by factors $`O(m_{\mathrm{light}}/m_Z^{})`$, and therefore not taken into account. This kind of the mass hierarchy corresponds to the case when the basis theory is a subgroup of the underlying high energy model remaining unknown. As our consideration shown, only two types of the $`Z^{}`$ couplings to light particles are consistent with the renormalizability. The first type corresponds to the Abelian couplings respecting the $`\stackrel{~}{\mathrm{U}}(1)`$ symmetry of the effective Lagrangian (20). In this case, the RG correlations fix the gauge symmetry of the Yukawa terms, which relates the fermion and the scalar hypercharges. As a consequence, the $`Z^{}`$ couplings to the axial-vector fermion currents are completely determined by the scalar field hypercharge and the fermion isospin. The second set of solutions – chiral $`Z^{}`$ – describes interactions with the SM particles similar to the third component of the $`\mathrm{SU}(2)_L`$ gauge field. The characteristic feature of the latter couplings is the zero traces of generators associated with the $`Z^{}`$. Notice that the $`Z^{}`$ interactions of the chiral type result in the effective four-fermion couplings $`(\overline{f}_{1L}\gamma ^\mu \sigma ^af_{1L})(\overline{f}_{2L}\gamma ^\mu \sigma ^af_{2L})`$ described by the operators $`𝒪_{ll}^{(3)}`$, $`𝒪_{lq}^{(3)}`$, and $`𝒪_{qq}^{(1,3)}`$ according to the classification in Refs. . Since each type of the $`Z^{}`$ interactions corresponds to one of mentioned operators, there is a possibility to select interactions by constructing the proper observables. As was shown, the observables proposed in Ref. can be chosen in searching for the Abelian $`Z^{}`$ boson. Thus, the bounds on the $`Z^{}`$ couplings calculated therein are also applicable in case of the THDM. The above note is important for the model independent search for $`Z^{}`$ virtual states at LEP2 and future colliders LHC and NLC. In the analysis of experimental data no discriminations between these two cases have been discussed in literature (see, for instance, recent survey or report ). This difference should be important for the model-dependent $`Z^{}`$ search when different scenarios of symmetry breaking are discussed. We believe that the derived RG relations to be useful in improving of experimental bounds on either the parameters of the $`Z^{}`$ interaction with fermions and on the relations between the cross sections of various four-fermion scattering processes. ## Acknoledgments The authors thank S. V. Peletminski and N. F. Shul’ga for discussions. ## A Feynman rules In what follows we use the notation $`\omega _{L,R}=(1\gamma ^5)/2`$, and all the momenta in the vertices are understood to be incoming. The Feynman rules for vertices of Figs. 1, 2 are listed below: 1. Fermion-vector vertices $`\overline{f}fA_\mu :`$ $`g\mathrm{sin}\theta _WQ_f\gamma ^\mu ;`$ (A1) $`\overline{f}fZ_\mu :`$ $`{\displaystyle \frac{g}{\mathrm{cos}\theta _W}}\gamma ^\mu \left(T_f^3\omega _LQ_f\mathrm{sin}^2\theta _W\right)`$ (A3) $`+O(\theta _0);`$ $`\overline{f}fZ_\mu ^{}:`$ $`{\displaystyle \frac{\stackrel{~}{g}}{2}}\gamma ^\mu \left(\omega _L\stackrel{~}{Y}_{L,f}+\omega _R\stackrel{~}{Y}_{R,f}\right)+O(\theta _0);`$ (A4) $`\overline{f}_df_uW_\mu ^{}:`$ $`{\displaystyle \frac{g}{\sqrt{2}}}\gamma ^\mu \omega _L;`$ (A5) $`\overline{f}_uf_dW_\mu ^+:`$ $`{\displaystyle \frac{g}{\sqrt{2}}}\gamma ^\mu \omega _L;`$ (A6) 2. Fermion-scalar vertices $`\overline{f}fH:`$ $`\left(G_{f,1}\mathrm{cos}\alpha +G_{f,2}\mathrm{sin}\alpha \right);`$ (A7) $`\overline{f}fh:`$ $`\left(G_{f,1}\mathrm{sin}\alpha G_{f,2}\mathrm{cos}\alpha \right);`$ (A8) $`\overline{f}fA_0:`$ $`2iT_f^3\left(\omega _L\omega _R\right)`$ (A10) $`\times \left(G_{f,1}\mathrm{sin}\beta G_{f,2}\mathrm{cos}\beta \right);`$ $`\overline{f}f\chi _3:`$ $`2iT_f^3\left(\omega _L\omega _R\right)`$ (A12) $`\times \left(G_{f,1}\mathrm{cos}\beta +G_{f,2}\mathrm{sin}\beta \right);`$ $`\overline{f}_df_uH^{}:`$ $`\sqrt{2}[\omega _L(G_{f_d,1}\mathrm{sin}\beta G_{f_d,2}\mathrm{cos}\beta )`$ (A14) $`+\omega _R(G_{f_u,1}\mathrm{sin}\beta +G_{f_u,2}\mathrm{cos}\beta )];`$ $`\overline{f}_uf_dH^+:`$ $`\sqrt{2}[\omega _R(G_{f_d,1}\mathrm{sin}\beta G_{f_d,2}\mathrm{cos}\beta )`$ (A16) $`+\omega _L(G_{f_u,1}\mathrm{sin}\beta +G_{f_u,2}\mathrm{cos}\beta )];`$ $`\overline{f}_df_u\chi ^{}:`$ $`\sqrt{2}[\omega _L(G_{f_d,1}\mathrm{cos}\beta +G_{f_d,2}\mathrm{sin}\beta )`$ (A18) $`+\omega _R(G_{f_u,1}\mathrm{cos}\beta +G_{f_u,2}\mathrm{sin}\beta )];`$ $`\overline{f}_uf_d\chi ^+:`$ $`\sqrt{2}[\omega _R(G_{f_d,1}\mathrm{cos}\beta +G_{f_d,2}\mathrm{sin}\beta )`$ (A20) $`+\omega _L(G_{f_u,1}\mathrm{cos}\beta +G_{f_u,2}\mathrm{sin}\beta )];`$ 3. $`Z^{}`$ scalar vertices $`Z_\mu ^{}H^+H^{}:`$ $`{\displaystyle \frac{\stackrel{~}{g}}{2}}(p_{H^+}p_H^{})_\mu (\stackrel{~}{Y}_{\varphi _1,1}\mathrm{sin}^2\beta `$ (A22) $`+\stackrel{~}{Y}_{\varphi _2,1}\mathrm{cos}^2\beta )+O(\theta _0);`$ $`Z_\mu ^{}H^+\chi ^{}:`$ $`{\displaystyle \frac{\stackrel{~}{g}\mathrm{sin}2\beta }{4}}\left(p_\chi ^{}p_{H^+}\right)_\mu `$ (A24) $`\times \left(\stackrel{~}{Y}_{\varphi _1,1}\stackrel{~}{Y}_{\varphi _2,1}\right)+O(\theta _0);`$ $`Z_\mu ^{}H^{}\chi ^+:`$ $`{\displaystyle \frac{\stackrel{~}{g}\mathrm{sin}2\beta }{4}}\left(p_H^{}p_{\chi ^+}\right)_\mu `$ (A26) $`\times \left(\stackrel{~}{Y}_{\varphi _1,1}\stackrel{~}{Y}_{\varphi _2,1}\right)+O(\theta _0);`$ $`Z_\mu ^{}\chi ^+\chi ^{}:`$ $`{\displaystyle \frac{\stackrel{~}{g}}{2}}(p_{\chi ^+}p_\chi ^{})_\mu (\stackrel{~}{Y}_{\varphi _1,1}\mathrm{cos}^2\beta `$ (A28) $`+\stackrel{~}{Y}_{\varphi _2,1}\mathrm{sin}^2\beta )+O(\theta _0);`$ $`Z_\mu ^{}HA_0:`$ $`{\displaystyle \frac{i\stackrel{~}{g}}{2}}\left(p_{A_0}p_H\right)_\mu \stackrel{~}{Y}_{\varphi ,2}\mathrm{sin}\left(\alpha \beta \right)`$ (A30) $`+O(\theta _0);`$ $`Z_\mu ^{}H\chi _3:`$ $`{\displaystyle \frac{i\stackrel{~}{g}}{2}}\left(p_{\chi _3}p_H\right)_\mu \stackrel{~}{Y}_{\varphi ,2}\mathrm{cos}\left(\alpha \beta \right)`$ (A32) $`+O(\theta _0);`$ $`Z_\mu ^{}hA_0:`$ $`{\displaystyle \frac{i\stackrel{~}{g}}{2}}\left(p_{A_0}p_h\right)_\mu \stackrel{~}{Y}_{\varphi ,2}\mathrm{cos}\left(\alpha \beta \right)`$ (A34) $`+O(\theta _0);`$ $`Z_\mu ^{}h\chi _3:`$ $`{\displaystyle \frac{i\stackrel{~}{g}}{2}}\left(p_hp_{\chi _3}\right)_\mu \stackrel{~}{Y}_{\varphi ,2}\mathrm{sin}\left(\alpha \beta \right)`$ (A36) $`+O(\theta _0).`$
warning/0004/hep-ph0004011.html	ar5iv	text	# 1 Introduction ## 1 Introduction The axial $`U(1)`$ anomaly in QCD implies an additional term in the Lagrangian, which violates $`P`$, $`T`$ and $`CP`$. This new term is proportional to the so-called vacuum angle $`\theta _0`$, an unknown parameter, and its size may be determined from $`CP`$-violating effects, as e.g. $`\eta \pi \pi `$ or the electric dipole moments of the neutron and the $`\mathrm{\Lambda }`$. The most recent measurements of the electric dipole moment of the neutron $`d_n^\gamma `$ have constrained it to $$\|d_n^\gamma \|<6.3\times 10^{26}e\text{cm}.$$ (1) On the other hand, theoretical estimates for $`d_n^\gamma `$ induced by the $`\theta _0`$-term can be given leading to an upper bound for $`\theta _0`$ . Of particular interest here is the estimate of Pich and de Rafael who used an effective chiral Lagrangian approach and came to the conclusion that it is possible to obtain an estimate for the size of the vacuum angle $`\theta _0`$ with an experimental upper limit of $`\|\theta _0\|5\times 10^{10}`$. However, the authors worked in a relativistic framework which does not have a systematic chiral counting scheme, so that higher loop diagrams contribute to lower chiral orders. This problem is avoided in heavy baryon chiral perturbation theory as proposed in , which allows for a consistent power counting. One should therefore study this effect within the heavy baryon formulation. Furthermore, the baryon Lagrangian in which describes the interactions of the neutron with the pseudoscalar nonet ($`\pi ,K,\eta ,\eta ^{}`$) does not contain explicitly $`CP`$-violating terms. They are rather induced by the vacuum alignment of the purely mesonic Lagrangian in the presence of the $`\theta _0`$-term. As shown in the most general baryonic Lagrangian taking the axial $`U(1)`$ anomaly into account does have explicitly $`CP`$-violating terms even at lower chiral orders. It has to be checked, if such terms lead to sizeable contributions for the present consideration. Finally, the authors of proposed to estimate the contribution from unknown counterterms by varying the scale in the chiral logarithms. This procedure reveals the scale dependence of the involved coupling constants but not their absolute value, and it is desirable to have a somewhat more reliable estimate of the involved couplings. The aim of the present work is to reinvestigate the electric dipole moments of the neutron and the $`\mathrm{\Lambda }`$ by taking the above mentioned points into consideration. One has to check, if it is still possible to give a reliable estimate of the vacuum angle $`\theta _0`$, and if it is different from the one given in . The paper is organized as follows. In the next section we present the purely mesonic effective Lagrangian in the presence of the $`\theta _0`$-term and the vacuum alignment is discussed. Baryon fields are included in the effective theory in Sec. 3 by using the method outlined in . We proceed by calculating in Sec. 4 the electric dipole moment of the neutron and the $`\mathrm{\Lambda }`$ up to one-loop order within the heavy baryon framework. Numerical results and conclusions are given in Sec. 5. ## 2 The mesonic Lagrangian In this section, we will consider the purely mesonic Lagrangian in the presence of the $`\theta _0`$-term. The derivation of this Lagrangian has been given elsewhere, see e.g. , so we will restrict ourselves to the repetition of some of the basic formulae which are needed in the present work. In the topological charge operator coupled to an external field is added to the QCD Lagrangian $$=_{QCD}\frac{g^2}{16\pi ^2}\theta (x)\text{tr}_c(G_{\mu \nu }\stackrel{~}{G}^{\mu \nu })$$ (2) with $`\stackrel{~}{G}_{\mu \nu }=ϵ_{\mu \nu \alpha \beta }G^{\alpha \beta }`$ and $`\text{tr}_c`$ is the trace over the color indices. Under $`U(1)_R\times U(1)_L`$ the axial $`U(1)`$ anomaly adds a term $`(g^2/16\pi ^2)2N_f\alpha \text{tr}_c(G_{\mu \nu }\stackrel{~}{G}^{\mu \nu })`$ to the QCD Lagrangian, with $`N_f`$ being the number of different quark flavors and $`\alpha `$ the angle of the global axial $`U(1)`$ rotation. The vacuum angle $`\theta (x)`$ is in this context treated as an external field that transforms under an axial $`U(1)`$ rotation as $$\theta (x)\theta ^{}(x)=\theta (x)2N_f\alpha .$$ (3) Then the term generated by the anomaly in the fermion determinant is compensated by the shift in the $`\theta `$ source and the Lagrangian from Eq. (2) remains invariant under axial $`U(1)`$ transformations. The symmetry group $`SU(3)_R\times SU(3)_L`$ of the Lagrangian $`_{QCD}`$ is extended<sup>3</sup><sup>3</sup>3To be more precise, the Lagrangian changes by a total derivative which gives rise to the Wess-Zumino term. We will neglect this contribution since the corresponding terms involve five or more meson fields which do not play any role for the discussions here. to $`U(3)_R\times U(3)_L`$ for $``$. The Green functions of QCD are obtained by expanding the generating functional around $`\theta (x)=\theta _0`$ where the phase of the quark mass matrix emerging from the Yukawa couplings of the light quarks in the electroweak sector has been absorbed in $`\theta _0`$. The extended symmetry remains at the level of an effective theory and the additional source $`\theta `$ also shows up in the effective Lagrangian. Let us consider the purely mesonic effective theory first. The lowest lying pseudoscalar meson nonet is summarized in a matrix valued field $`\stackrel{~}{U}(x)`$. The effective Lagrangian is formed with the fields $`\stackrel{~}{U}(x)`$, derivatives thereof and also includes both the quark mass matrix $``$ and the vacuum angle $`\theta `$: $`_{\text{eff}}(\stackrel{~}{U},\stackrel{~}{U},\mathrm{},,\theta )`$. Under $`U(3)_R\times U(3)_L`$ the fields transform as follows: $$\stackrel{~}{U}^{}=R\stackrel{~}{U}L^{},^{}=RL^{},\theta ^{}(x)=\theta (x)2N_f\alpha $$ (4) with $`RU(3)_R`$, $`LU(3)_L`$, but the Lagrangian remains invariant. The phase of the determinant of $`\stackrel{~}{U}(x)`$ transforms under axial $`U(1)`$ as $`\mathrm{ln}\text{det}\stackrel{~}{U}^{}(x)=\mathrm{ln}\text{det}\stackrel{~}{U}(x)+2iN_f\alpha `$ so that the combination $`\theta i\mathrm{ln}\text{det}\stackrel{~}{U}`$ remains invariant. It is more convenient to replace the variable $`\theta `$ by this invariant combination, $`_{\text{eff}}(\stackrel{~}{U},\stackrel{~}{U},\mathrm{},,\theta i\mathrm{ln}\text{det}\stackrel{~}{U})`$. One can now construct the effective Lagrangian in these fields that respects the symmetries of the underlying theory. In particular, the Lagrangian is invariant under $`U(3)_R\times U(3)_L`$ rotations of $`\stackrel{~}{U}`$ and $``$ at a fixed value of the last argument. The Lagrangian up to and including terms with two derivatives and one factor of $``$ reads $`_\varphi `$ $`=`$ $`V_0+V_1_\mu \stackrel{~}{U}^{}^\mu \stackrel{~}{U}+V_2\stackrel{~}{\chi }^{}\stackrel{~}{U}+\stackrel{~}{\chi }\stackrel{~}{U}^{}+iV_3\stackrel{~}{\chi }^{}\stackrel{~}{U}\stackrel{~}{\chi }\stackrel{~}{U}^{}`$ (5) $`+V_4\stackrel{~}{U}_\mu \stackrel{~}{U}^{}\stackrel{~}{U}^{}^\mu \stackrel{~}{U}.`$ The expression $`\mathrm{}`$ denotes the trace in flavor space and $`\stackrel{~}{\chi }`$ is proportional to the quark mass matrix $`\stackrel{~}{\chi }=\stackrel{~}{\chi }^{}=2B_0`$ with $`=\text{diag}(m_u,m_d,m_s)`$ and $`B_0=0\|\overline{q}q\|0/F_\pi ^2`$ the order parameter of the spontaneous symmetry violation. The covariant derivative of $`\stackrel{~}{U}`$ is defined by $$_\mu \stackrel{~}{U}=_\mu \stackrel{~}{U}i(v_\mu +a_\mu )\stackrel{~}{U}+i\stackrel{~}{U}(v_\mu a_\mu ).$$ (6) The external fields $`v_\mu (x),a_\mu (x)`$ represent Hermitian $`3\times 3`$ matrices in flavor space. Note that a term of the type $`iV_5\stackrel{~}{U}^{}_\mu \stackrel{~}{U}^\mu \theta `$ can be transformed away and a term proportional to $`V_6_\mu \theta ^\mu \theta `$ does not enter the calculations performed in the present work and will be neglected. The coefficients $`V_i`$ are functions of the variable $`\theta i\mathrm{ln}\text{det}\stackrel{~}{U}`$, $`V_i(\theta i\mathrm{ln}\text{det}\stackrel{~}{U})`$, and can be expanded in terms of this variable. The terms $`V_{1,\mathrm{},4}`$ are of second chiral order, whereas $`V_0`$ is of zeroth chiral order. Parity conservation implies that the $`V_i`$ are all even functions of $`\theta i\mathrm{ln}\text{det}\stackrel{~}{U}`$ except $`V_3`$, which is odd, and $`V_1(0)=V_2(0)=F_\pi ^2/4`$ gives the correct normalizaton for the quadratic terms of the Goldstone boson octet, where $`F_\pi 92.4`$ MeV is the pion decay constant. In order to use the effective Lagrangian, one must first determine the vacuum expectation value of $`\stackrel{~}{U}`$ by minimizing the potential energy $$V(\stackrel{~}{U})=V_0V_2\stackrel{~}{\chi }^{}\stackrel{~}{U}+\stackrel{~}{\chi }\stackrel{~}{U}^{}iV_3\stackrel{~}{\chi }^{}\stackrel{~}{U}\stackrel{~}{\chi }\stackrel{~}{U}^{}.$$ (7) Since $`\stackrel{~}{\chi }`$ is diagonal, one can assume the minimum $`U_0`$ to be diagonal as well and of the form $$U_0=\text{diag}(\text{e}^{i\phi _u},\text{e}^{i\phi _d},\text{e}^{i\phi _s}).$$ (8) In terms of the angles $`\phi _q`$ the potential becomes $$V(U_0)=V_0(\overline{\theta }_0)4B_0V_2(\overline{\theta }_0)\underset{q}{}m_q\mathrm{cos}\phi _q4B_0V_3(\overline{\theta }_0)\underset{q}{}m_q\mathrm{sin}\phi _q,$$ (9) where we have introduced the notation $`\overline{\theta }_0=\theta _0_q\phi _q`$. The Taylor expansions of the functions $`V_i`$ read $`V_i(\overline{\theta }_0)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}V_i^{(2n)}\overline{\theta }_0^{2n}\text{for}i=0,2`$ $`V_3(\overline{\theta }_0)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}V_i^{(2n+1)}\overline{\theta }_0^{2n+1}`$ (10) with coefficients not fixed by chiral symmetry. Minimizing the potential with respect to the angles $`\phi _q`$ leads to $$2B_0m_q\mathrm{sin}\phi _q=𝒜+2B_0m_q\mathrm{cos}\phi _q$$ (11) with $`𝒜`$ $`=`$ $`2B_0\left({\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}V_2^{(2n)}\overline{\theta }_0^{2n}\right)^1{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}({\displaystyle \frac{1}{2B_0}}nV_0^{(2n)}\overline{\theta }_0^{2n1}`$ (12) $`2nV_2^{(2n)}\overline{\theta }_0^{2n1}{\displaystyle \underset{j=u,d,s}{}}m_j\mathrm{cos}\phi _j`$ $`(2n1)V_3^{(2n1)}\overline{\theta }_0^{2n2}{\displaystyle \underset{j=u,d,s}{}}m_j\mathrm{sin}\phi _j)`$ and $$=\left(\underset{n=0}{\overset{\mathrm{}}{}}V_2^{(2n)}\overline{\theta }_0^{2n}\right)^1\underset{n=0}{\overset{\mathrm{}}{}}V_3^{(2n+1)}\overline{\theta }_0^{2n+1}.$$ (13) To lowest order both in the quark masses $`m_q`$ and 1/$`N_c`$ Eq. (11) reads $$\frac{1}{2}B_0F_\pi ^2m_q\mathrm{sin}\phi _q=V_0^{(2)}\overline{\theta }_0$$ (14) which is the equation for the $`\phi _q`$ considered in . One then writes $$\stackrel{~}{U}=\sqrt{U_0}U\sqrt{U_0}$$ (15) and $`U`$ can be parametrized as $$U(\varphi ,\eta _0)=\mathrm{exp}\{2i\varphi /F_\pi +i\sqrt{\frac{2}{3}}\eta _0/F_0\},$$ (16) where the singlet $`\eta _0`$ couples to the singlet axial current with strength $`F_0`$. The unimodular part of the field $`U(x)`$ contains the degrees of freedom of the Goldstone boson octet $`\varphi `$ $`\varphi ={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\pi ^0+\frac{1}{\sqrt{6}}\eta _8& \pi ^+& K^+\\ \pi ^{}& \frac{1}{\sqrt{2}}\pi ^0+\frac{1}{\sqrt{6}}\eta _8& K^0\\ K^{}& \overline{K^0}& \frac{2}{\sqrt{6}}\eta _8\end{array}\right),`$ while the phase det$`U(x)=e^{i\sqrt{6}\eta _0/F_0}`$ describes the $`\eta _0`$ . The diagonal subgroup $`U(3)_V`$ of $`U(3)_R\times U(3)_L`$ does not have a dimension-nine irreducible representation and consequently does not exhibit a nonet symmetry. We have therefore used the different notation $`F_0`$ for the decay constant of the singlet field. One can now express the effective Lagrangian in terms of the Goldstone boson matrix $`U`$ and the angles $`\phi _q`$ $`_\varphi `$ $`=`$ $`V_0+V_1_\mu U^{}^\mu U+[V_2+V_3]\chi (U+U^{})i𝒜V_2UU^{}`$ (17) $`+`$ $`i[V_3V_2]\chi (UU^{})+𝒜V_3U+U^{}+V_4U_\mu U^{}U^{}^\mu U,`$ where we have absorbed the angles $`\phi _q`$ in hermitian matrices $`\chi `$ and $`H`$ by defining $$\sqrt{U_0^{}}\stackrel{~}{\chi }\sqrt{U_0^{}}=\chi +iH,\sqrt{U_0}\stackrel{~}{\chi }^{}\sqrt{U_0}=\chi iH,$$ (18) so that $`\chi =2B_0\text{diag}(m_q\mathrm{cos}\phi _q)`$ and $`H=2B_0\text{diag}(m_q\mathrm{sin}\phi _q)=𝒜+\chi `$. The $`V_i`$ are functions of $`\sqrt{6}\eta _0/F_0+\overline{\theta }_0`$, $`V_i(\sqrt{6}\eta _0/F_0+\overline{\theta }_0)`$ and we have assumed the external fields $`v_\mu `$ and $`a_\mu `$ to be diagonal which is the case if one considers electromagnetic interactions. Note that the Goldstone boson masses at lowest chiral order are not only functions of the current quark masses $`m_q`$, but also depend on the angles $`\mathrm{cos}\phi _q`$. The kinetic energy of the $`\eta _0`$ singlet field obtains contributions from $`V_1_\mu U^{}^\mu U`$ and $`V_4U_\mu U^{}U^{}^\mu U`$ which read $$\left(\frac{F_\pi ^2}{2F_0^2}+\frac{6}{F_0^2}V_4(0)\right)_\mu \eta _0^\mu \eta _0.$$ (19) We renormalize the $`\eta _0`$ field in such a way that the coefficient in brackets is 1/2 in analogy to the kinetic term of the octet. By redefining $`F_0`$ and keeping for simplicity the same notation both for $`\eta _0`$ and $`F_0`$ one arrives at the same Lagrangian as in Eq. (17) but with $`V_4(0)=(F_0^2F_\pi ^2)/12`$ in order to ensure the usual normalization for the kinetic term of a pseudoscalar particle. ## 3 $`CP`$-violating terms in the baryon Lagrangian As mentioned before, another source for $`CP`$-violation is the baryon Lagrangian. The $`CP`$-violating terms can be divided into two groups. Firstly, the vacuum alignment of the mesonic Lagrangian induces $`CP`$ non-conserving meson-baryon interactions as considered in . But secondly, there are also explicitly $`CP`$-violating terms in the most general Lagrangian in the presence of the $`\theta `$-vacuum angle which have been neglected in . In order to construct the Lagrangian in the baryon-sector, one has to adopt a non-linear representations for baryons. The main ingredient for a non-linear realization is the compensator field $`K(\stackrel{~}{U},R,L)U(3)_V`$, which appears in the chiral $`U(3)_L\times U(3)_R`$ transformation of the left and right coset representatives, $`\stackrel{~}{\xi }_L(\stackrel{~}{U})`$ and $`\stackrel{~}{\xi }_R(\stackrel{~}{U})`$: $`\stackrel{~}{\xi }_L(\stackrel{~}{U})`$ $``$ $`L\stackrel{~}{\xi }_L(\stackrel{~}{U})K^{}(\stackrel{~}{U},R,L)`$ $`\stackrel{~}{\xi }_R(\stackrel{~}{U})`$ $``$ $`R\stackrel{~}{\xi }_R(\stackrel{~}{U})K^{}(\stackrel{~}{U},R,L).`$ (20) The field $`\stackrel{~}{U}`$ from the last section is defined as $$\stackrel{~}{U}=\stackrel{~}{\xi }_R\stackrel{~}{\xi }_L^{}.$$ (21) The baryon octet $`B`$ is given by the matrix $`B=\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\mathrm{\Sigma }^0+\frac{1}{\sqrt{6}}\mathrm{\Lambda }& \mathrm{\Sigma }^+& p\\ \mathrm{\Sigma }^{}& \frac{1}{\sqrt{2}}\mathrm{\Sigma }^0+\frac{1}{\sqrt{6}}\mathrm{\Lambda }& n\\ \mathrm{\Xi }^{}& \mathrm{\Xi }^0& \frac{2}{\sqrt{6}}\mathrm{\Lambda }\end{array}\right)`$ which transforms as a matter field $$BB^{}=KBK^{}.$$ (22) The covariant derivative of the baryon fields reads $$[\stackrel{~}{D}_\mu ,B]=_\mu B+[\stackrel{~}{\mathrm{\Gamma }}_\mu ,B]$$ (23) with $`\stackrel{~}{\mathrm{\Gamma }}_\mu `$ being the chiral connection $$\stackrel{~}{\mathrm{\Gamma }}_\mu =\frac{1}{2}\left[\stackrel{~}{\xi }_R^{}(_\mu ir_\mu )\stackrel{~}{\xi }_R+\stackrel{~}{\xi }_L^{}(_\mu il_\mu )\stackrel{~}{\xi }_L\right]$$ (24) and $`r_\mu =v_\mu +a_\mu `$, $`l_\mu =v_\mu a_\mu `$. For electromagnetic interactions the external fields are $`a_\mu =0`$ and $`v_\mu =eQ𝒜_\mu `$ with the quark charge matrix $`Q=\frac{1}{3}\text{diag}(2,1,1)`$. In order to incorporate the interactions with the mesons into the effective theory it is convenient to form an object of axial-vector type with one derivative $$\stackrel{~}{\xi }_\mu =i\left[\stackrel{~}{\xi }_R^{}(_\mu ir_\mu )\stackrel{~}{\xi }_R\stackrel{~}{\xi }_L^{}(_\mu il_\mu )\stackrel{~}{\xi }_L\right].$$ (25) Further ingredients of the non-linear representation are $$\stackrel{~}{\chi }_\pm =\stackrel{~}{\xi }_L^{}\stackrel{~}{\chi }^{}\stackrel{~}{\xi }_R\pm \stackrel{~}{\xi }_R^{}\stackrel{~}{\chi }\stackrel{~}{\xi }_L$$ (26) and the quantity $$\stackrel{~}{F}_{\mu \nu }^\pm =\stackrel{~}{\xi }_R^{}F_{\mu \nu }^R\stackrel{~}{\xi }_R\pm \stackrel{~}{\xi }_L^{}F_{\mu \nu }^L\stackrel{~}{\xi }_L,$$ (27) where $`F_{\mu \nu }^{R/L}`$ are the field strength tensors of $`r_\mu /l_\mu `$. The most general relativistic effective Lagrangian up to second order in the derivative expansion and contributing to the electric dipole moments of the neutron and $`\mathrm{\Lambda }`$ reads $`_{\varphi B}`$ $`=`$ $`iW_1[\stackrel{~}{D}^\mu ,\overline{B}]\gamma _\mu BiW_1^{}\overline{B}\gamma _\mu [\stackrel{~}{D}^\mu ,B]+W_2\overline{B}B`$ (28) $`+W_3\overline{B}\gamma _\mu \gamma _5\{\stackrel{~}{\xi }^\mu ,B\}+W_4\overline{B}\gamma _\mu \gamma _5[\stackrel{~}{\xi }^\mu ,B]+W_5\overline{B}\gamma _\mu \gamma _5B\stackrel{~}{\xi }^\mu `$ $`+iW_6\overline{B}\gamma _5B+W_7\overline{B}\{\stackrel{~}{\chi }_+,B\}+W_8\overline{B}[\stackrel{~}{\chi }_+,B]+W_9\overline{B}B\stackrel{~}{\chi }_+`$ $`+iW_{10}\overline{B}\{\stackrel{~}{\chi }_{},B\}+iW_{11}\overline{B}[\stackrel{~}{\chi }_{},B]+iW_{12}\overline{B}B\stackrel{~}{\chi }_{}`$ $`+iW_{13}\overline{B}\sigma _{\mu \nu }\gamma _5\{\stackrel{~}{F}_+^{\mu \nu },B\}+iW_{14}\overline{B}\sigma _{\mu \nu }\gamma _5[\stackrel{~}{F}_+^{\mu \nu },B]`$ $`+iW_{15}\overline{B}\sigma _{\mu \nu }\gamma _5B\stackrel{~}{F}_+^{\mu \nu }.`$ The $`W_i`$ are functions of the combination $`\sqrt{6}\eta _0/F_0+\overline{\theta }_0`$. From parity it follows that $`W_{1,\mathrm{},5}`$ and $`W_{7,8,9}`$ are even in this variable, whereas $`W_6`$ and $`W_{10,\mathrm{},15}`$ are odd. The latter have not been taken into account in . One can further reduce the number of independent terms by making the following transformation. By decomposing the baryon fields into their left- and right handed components $$B_{R/L}=\frac{1}{2}(1\pm \gamma _5)B$$ (29) and transforming the left- and right-handed states separately via $`B_{R/L}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{W_2\pm iW_6}}}B_{R/L}`$ $`\overline{B}_{R/L}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{W_2iW_6}}}\overline{B}_{R/L}`$ (30) one can eliminate the $`\overline{B}\gamma _5B`$ term and simplify the coefficient of $`\overline{B}B`$. The details of this calculation are given in . Note that this transformation leads to mixing of the terms of the type $`\overline{B}\stackrel{~}{\chi }_\pm B`$ with $`\overline{B}\gamma _5\stackrel{~}{\chi }_\pm B`$ which are of third chiral order and have been neglected here. Furthermore, the terms $`W_{13,\mathrm{},15}`$ mix with the terms from the Lagrangian $$=W_{16}\overline{B}\sigma _{\mu \nu }\{\stackrel{~}{F}_+^{\mu \nu },B\}+W_{17}\overline{B}\sigma _{\mu \nu }[\stackrel{~}{F}_+^{\mu \nu },B]+W_{18}\overline{B}\sigma _{\mu \nu }B\stackrel{~}{F}_+^{\mu \nu },$$ but in both cases the form of the Lagrangian does not change and we can proceed by neglecting the $`W_6`$ term and setting $`W_1=W_1^{}`$ and $`W_2=\stackrel{}{M}`$ with $`\stackrel{}{M}`$ being the baryon mass in the chiral limit . The expansion of the coefficients in terms of the $`W_i`$ read $`W_1={\displaystyle \frac{1}{2}}+\mathrm{},W_2=\stackrel{}{M},`$ $`W_3={\displaystyle \frac{1}{2}}D+\mathrm{},W_4={\displaystyle \frac{1}{2}}F+\mathrm{},W_5={\displaystyle \frac{1}{2}}\lambda +\mathrm{},`$ $`W_7=b_D+\mathrm{},W_8=b_F+\mathrm{},W_9=b_0+\mathrm{},`$ $`W_i=w_i({\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0+\overline{\theta }_0)+\mathrm{},\text{for}i=10,\mathrm{},15,`$ (31) where the ellipses denote higher orders in $`\sqrt{6}\eta _0/F_0+\overline{\theta }_0`$ and we have only shown terms that can contribute to one-loop order. The axial-vector couplings $`D`$ and $`F`$ can be determined from semileptonic hyperon decays. A fit to the experimental data delivers $`D=0.80\pm 0.01`$ and $`F=0.46\pm 0.01`$ . The third coupling, $`\lambda `$, is specific to the axial flavor-singlet baryonic current. The coefficients $`b_{D,F,0}`$ have been determined from the calculation of the baryon masses and the $`\pi N`$ $`\sigma `$-term up to fourth chiral order and their mean values are, in units of GeV<sup>-1</sup>, $$b_D=0.079,b_F=0.316,b_0=0.606.$$ (32) The numerical values of the parameters $`w_{10,\mathrm{},15}`$ are not known. We leave them undetermined for the time being and will give later an upper bound. Again the correct vacuum has to be chosen. This is done by setting $`\stackrel{~}{\xi }_R`$ $`=`$ $`\sqrt{U_0}\xi _R`$ $`\stackrel{~}{\xi }_L`$ $`=`$ $`\sqrt{U_0^{}}\xi _L,`$ (33) and we will choose the coset representatives such that $$\xi _R=\xi _L^{}=u=\sqrt{U}.$$ (34) For diagonal external fields $`v_\mu `$ and $`a_\mu `$ one can write $$\stackrel{~}{\mathrm{\Gamma }}_\mu =\mathrm{\Gamma }_\mu =\frac{1}{2}\left[u^{}(_\mu ir_\mu )u+u(_\mu il_\mu )u^{}\right]$$ (35) and $$\stackrel{~}{\xi }_\mu =u_\mu =i\left[u^{}(_\mu ir_\mu )uu(_\mu il_\mu )u^{}\right].$$ (36) Furthermore, one obtains $`\stackrel{~}{\chi }_+`$ $`=`$ $`\chi _+i𝒜(UU^{})i\chi _{}`$ $`\stackrel{~}{\chi }_{}`$ $`=`$ $`\chi _{}i𝒜(U+U^{})i\chi _+,`$ (37) where the quark mass matrix enters in the combinations $$\chi _\pm =u\chi ^{}u\pm u^{}\chi u^{}.$$ (38) Finally, $`\stackrel{~}{F}_{\mu \nu }^+`$ simplifies to $$\stackrel{~}{F}_{\mu \nu }^+=F_{\mu \nu }^+=u^{}F_{\mu \nu }^Ru+uF_{\mu \nu }^Lu^{}.$$ (39) The relativistic baryon Lagrangian reads to the order we are working $`_{\varphi B}`$ $`=`$ $`i\overline{B}\gamma _\mu [D^\mu ,B]\stackrel{}{M}\overline{B}B`$ (40) $`{\displaystyle \frac{1}{2}}D\overline{B}\gamma _\mu \gamma _5\{u^\mu ,B\}{\displaystyle \frac{1}{2}}F\overline{B}\gamma _\mu \gamma _5[u^\mu ,B]+{\displaystyle \frac{1}{2}}\lambda \overline{B}\gamma _\mu \gamma _5Bu^\mu `$ $`ib_D𝒜\overline{B}\{UU^{},B\}ib_F𝒜\overline{B}[UU^{},B]ib_0𝒜\overline{B}BUU^{}`$ $`+4𝒜w_{10}{\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0\overline{B}B+6𝒜w_{12}{\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0\overline{B}B`$ $`+i(w_{13}^{}\overline{\theta }_0+w_{13}{\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0)\overline{B}\sigma _{\mu \nu }\gamma _5\{F_+^{\mu \nu },B\}`$ $`+i(w_{14}^{}\overline{\theta }_0+w_{14}{\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0)\overline{B}\sigma _{\mu \nu }\gamma _5[F_+^{\mu \nu },B]`$ $`+i(w_{15}^{}\overline{\theta }_0+w_{15}{\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0)\overline{B}\sigma _{\mu \nu }\gamma _5BF_+^{\mu \nu },`$ where we have neglected meson-baryon interactions with more than one meson field since they do not contribute at one-loop order and terms of $`𝒪(\overline{\theta }_0^2)`$ or higher orders are omitted throughout this work. Note that there exist terms of fourth chiral order of the type $`\overline{B}\sigma _{\mu \nu }\gamma _5\stackrel{~}{\chi }_{}\stackrel{~}{F}_+^{\mu \nu }B`$. Using Eq. (3) for $`\stackrel{~}{\chi }_{}`$ they induce $`CP`$-violating terms of the form $`\overline{\theta }_0\overline{B}\sigma _{\mu \nu }\gamma _5\stackrel{~}{F}_+^{\mu \nu }B`$ which are already accounted for by the terms $`w_{13,14,15}`$. This amounts to a renormalization of the couplings $`\overline{\theta }_0w_{13,14,15}`$. We have therefore introduced the notation $`w_{13,14,15}^{}`$ for these interaction terms, in order to distinguish them from the unrenormalized $`w_{13,14,15}`$ of the interactions proportional to $`\eta _0`$. The drawback of the relativistic framework including baryons is that due to the existence of a new mass scale, namely the baryon mass in the chiral limit $`\stackrel{}{M}`$, there exists no strict chiral counting scheme, i.e. a one-to-one correspondence between the meson loops and the chiral expansion. In order to overcome this problem one integrates out the heavy degrees of freedom of the baryons, similar to a Foldy-Wouthuysen transformation, so that a chiral counting scheme emerges. To this end, one constructs eigenstates of the velocity projection operator $`P_v=(1+v/)/2`$ $$B_v(x)=e^{i\stackrel{}{M}vx}P_vB(x).$$ (41) The Dirac algebra simplifies considerably. It allows to express any Dirac bilinear $`\overline{B_v}\mathrm{\Gamma }_\mu B_v(\mathrm{\Gamma }_\mu =1,\gamma _\mu ,\gamma _5,\mathrm{})`$ in terms of the velocity $`v_\mu `$ and the spin operator $`2S_\mu =i\gamma _5\sigma _{\mu \nu }v^\nu `$. One can rewrite the Dirac bilinears which appear in the present calculation as $`\overline{B_v}\gamma _\mu \gamma _5B_v=2\overline{B_v}S_\mu B_v,\overline{B_v}\sigma _{\mu \nu }\gamma _5B_v=2i(v_\mu \overline{B_v}S_\nu B_vv_\nu \overline{B_v}S_\mu B_v).`$ (42) In the following, we will drop the index $`v`$. The Lagrangian of the heavy baryon formulation reads $`_{\varphi B}`$ $`=`$ $`i\overline{B}[vD,B]D\overline{B}S_\mu \{u^\mu ,B\}F\overline{B}S_\mu [u^\mu ,B]+\lambda \overline{B}S_\mu Bu^\mu `$ (43) $`ib_D𝒜\overline{B}\{UU^{},B\}ib_F𝒜\overline{B}[UU^{},B]ib_0𝒜\overline{B}BUU^{}`$ $`+4𝒜w_{10}{\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0\overline{B}B+6𝒜w_{12}{\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0\overline{B}B`$ $`4(w_{13}^{}\overline{\theta }_0+w_{13}{\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0)\overline{B}v_\mu S_\nu \{F_+^{\mu \nu },B\}`$ $`4(w_{14}^{}\overline{\theta }_0+w_{14}{\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0)\overline{B}v_\mu S_\nu [F_+^{\mu \nu },B]`$ $`4(w_{15}^{}\overline{\theta }_0+w_{15}{\displaystyle \frac{\sqrt{6}}{F_0}}\eta _0)\overline{B}v_\mu S_\nu BF_+^{\mu \nu }.`$ No relativistic corrections are needed to the order we are working. ## 4 The electric dipole moments $`d_n^\gamma `$ and $`d_\mathrm{\Lambda }^\gamma `$ In this section, we will calculate the electric dipole moments of the neutron and the $`\mathrm{\Lambda }`$ at lowest order in chiral perturbation theory, i.e. $`𝒪(p^2)`$. Due to the vacuum alignment the baryon Lagrangian contains interaction terms $`b_{D,F,0}`$ of zeroth chiral order and one-loop diagrams with these vertices will also contribute at $`𝒪(p^2)`$. In the relativistic framework of the electric dipole moment of the neutron $`d_n^\gamma `$ has been defined via $$_{nEDM}=\frac{1}{2}d_n^\gamma \overline{n}i\sigma _{\mu \nu }\gamma _5nF^{\mu \nu },$$ (44) where $`F^{\mu \nu }`$ is the field strength tensor of the photon field $`A^\mu `$. We prefer to rewrite this as a form factor $$D_n^\gamma (q^2)\overline{u}(p^{})\sigma _{\mu \nu }\gamma _5u(p)q^\mu $$ (45) with $`q=p^{}p`$ being the momentum transfer. The electric dipole moment is given by $$d_n^\gamma =D_n^\gamma (q^2=0).$$ (46) For the calculation of the form factor in the heavy baryon approach we set $`v_\mu =(1,\mathrm{𝟎})`$ and use the Breit frame $`vp=vp^{}`$ since it allows a unique translation of Lorentz-covariant matrix elements into non-relativistic ones. In this frame the form factor reads $$2iD_n^\gamma (q^2)\overline{H}v_\nu SqH+\mathrm{},$$ (47) where $`H`$ is the large component of $`u`$ and the ellipsis stands for a similar expression in the small components of $`u`$, which is of higher chiral order and will be omitted. The electric dipole moment $`d_n^\gamma `$ receives contributions from the $`w_{13}`$-term and the loops (we work in the isospin limit $`m_u=m_d=\widehat{m}`$) $$d_n^\gamma =d_n^{\gamma (tree)}+d_n^{\gamma (loop)}$$ (48) with $$d_n^{\gamma (tree)}=8e\overline{\theta }_0\left[\frac{1}{3}w_{13}^r+\frac{16}{F_\pi ^2F_0^2m_{\eta _0}^2}V_3^{(1)}V_0^{(2)}w_{13}\right],$$ (49) where the pertinent diagrams are shown in Figure 1. Diagram 1b) is missing in since the mesonic Lagrangian used within this work does not have a term linear in the singlet field $`\eta _0`$. Such an interaction originates from the term $`V_3\chi _{}`$ which is not considered in . The chiral logarithms of the diagrams shown in Fig. 2 read $$d_n^{\gamma (loop)}=\frac{1}{\pi ^2F_\pi ^4}eV_0^{(2)}\overline{\theta }_0\left[(b_D+b_F)(D+F)\mathrm{ln}\frac{m_\pi ^2}{\mu ^2}+(b_Db_F)(DF)\mathrm{ln}\frac{m_K^2}{\mu ^2}\right]$$ (50) with $`\mu `$ being the scale introduced in dimensional regularization. In the Breit frame only diagrams 2a) and b) contribute to the electric dipole moment. Diagrams 2c), d) vanish and 2e), f) are proportional to $`S_\nu `$ and therefore do not contribute to $`d_n^\gamma `$. The loop integral for diagrams 2a) and b) contains also analytic and divergent pieces which can be absorbed by redefining $`w_{13}^{}`$. The divergent pieces of $`w_{13}^{}`$ cancel the divergencies from the loops and render the final expression finite. We summarize the remaining analytical contributions in $`w_{13}^r`$, so that $`w_{13}^r`$ in Eq. (49) is understood to be finite. The results in Eqs. (49) and (50) are not in contradiction with the fact that the electric dipole moment induced by the $`\theta _0`$-term tends to zero if any of the quark masses vanish. If, e.g., $`m_u=0`$ then a solution for Eq. (11) is given by $`\phi _u=\theta _0`$ and $`\phi _{d,s}=0`$ leading to $`\overline{\theta }_0=\theta _0\phi _u=0`$. Therefore, $`d_n^\gamma `$ vanishes in this case. The results for the $`\mathrm{\Lambda }`$ are $$d_\mathrm{\Lambda }^{\gamma (tree)}=4e\overline{\theta }_0\left[\frac{1}{3}w_{13}^r+\frac{16}{F_\pi ^2F_0^2m_{\eta _0}^2}V_3^{(1)}V_0^{(2)}w_{13}\right]$$ (51) and $$d_\mathrm{\Lambda }^{\gamma (loop)}=\frac{1}{\pi ^2F_\pi ^4}eV_0^{(2)}\overline{\theta }_0[b_DF+b_FD]\mathrm{ln}\frac{m_K^2}{\mu ^2}.$$ (52) Note that the relation $`d_n^\gamma =2d_\mathrm{\Lambda }^\gamma `$ is only valid at tree level and not for the chiral logarithms as claimed in . The discrepancy is due to the lack of pion loops in the present work. Once one accounts for the mistake made in by replacing $`\mathrm{ln}m_\pi `$ by $`\mathrm{ln}m_K`$ in the chiral loop contribution for $`d_\mathrm{\Lambda }^\gamma `$ one obtains our result (52). ## 5 Numerical results and conclusions In order to compute the numerical results for the electric dipole moments shown in the last section, we use the central values for the parameters $`b_D=0.079`$ GeV<sup>-1</sup>, $`b_F=0.316`$ GeV<sup>-1</sup>, $`D=0.80`$ and $`F=0.46`$. To lowest order in the angles $`\phi _q`$ and using $`\widehat{m}m_s`$ one can express $`\overline{\theta }_0`$ in terms of $`\theta _0`$ via $$\theta _0[1+\frac{8V_0^{(2)}}{F_\pi ^2m_\pi ^2}]\overline{\theta }_0.$$ (53) From the calculation of the $`\eta `$ and $`\eta ^{}`$ masses and decay constants one can extract the value for $`V_0^{(2)}`$ $$V_0^{(2)}\frac{27}{4}F_\pi ^45.0\times 10^4\text{GeV}^4,$$ (54) so that $$\overline{\theta }_0\frac{F_\pi ^2m_\pi ^2}{8V_0^{(2)}}\theta _00.04\theta _0.$$ (55) Inserting this into the loop contribution and using $`\mu =1`$ GeV, $`m_{\eta _0}m_\eta ^{}=958`$ MeV we obtain $`d_n^{\gamma (loop)}`$ $`=`$ $`7.5\times 10^{16}\theta _0e\text{cm}`$ $`d_\mathrm{\Lambda }^{\gamma (loop)}`$ $`=`$ $`1.7\times 10^{16}\theta _0e\text{cm}.`$ (56) The numerical result for $`d_n^{\gamma (loop)}`$ is in agreement with the one given in once one accounts for the different values of the parameters $`b_D,b_F,D,F`$ and the scale $`\mu `$ used within that work. The result for the chiral logarithm of $`d_\mathrm{\Lambda }^{\gamma (loop)}`$ is considerably smaller than for $`d_n^{\gamma (loop)}`$ since there is no contribution from the pion loops which dominate in the case of $`d_n^\gamma `$. A precise numerical value for the tree contribution to the electric dipole moments cannot be given since the parameters $`w_{13}`$ and $`w_{13}^r`$ are not known. However, we will give an upper bound for their contribution based on large $`N_c`$ arguments which seem to work well in the purely mesonic sector . In the present investigation we are only interested in an order of magnitude estimate for the vacuum angle $`\theta `$ and for this purpose it is sufficient to give a numerical range for the tree level contribution. We will first estimate the ratio of diagrams 1a) and 1b) using 1/$`N_c`$ arguments. Applying large $`N_c`$ counting rules, see e.g. , both $`w_{13}`$ and $`w_{13}^r`$ are of order $`𝒪(N_c^0)`$ so that we can assume $`\|w_{13}/w_{13}^r\|=𝒪(1)`$. One obtains the ratio $$\frac{\|d_n^{\gamma (1b)}\|}{\|d_n^{\gamma (1a)}\|}=\frac{\|d_\mathrm{\Lambda }^{\gamma (1b)}\|}{\|d_\mathrm{\Lambda }^{\gamma (1a)}\|}\frac{48}{F_\pi ^4m_\eta ^{}^2}\|V_0^{(2)}V_3^{(1)}\|0.12,$$ (57) where we have used $`F_\pi /F_0=1+𝒪(N_c^1)`$ and taken the value for $`V_3^{(1)}`$ from $$V_3^{(1)}0.04F_\pi ^23.5\times 10^4\text{GeV}^2.$$ (58) Diagram 1b) turns out to be insignificant in our estimate. Based on large $`N_c`$ arguments one can also give an upper bound for $`w_{13}^r`$. The contact terms of the Lagrangian in Eq. (3), which contribute to the magnetic moments of the baryons, describe – similar to the $`w_{13}^r`$-term – the coupling of the field strength tensor $`F_{\mu \nu }^+`$ to the baryons. But the leading coefficient of the Taylor expansion of $`W_{16}`$, $`w_{16}`$, is of order $`𝒪(N_c^1)`$, whereas $`\|w_{13}^r\|\|w_{13}\|=𝒪(N_c^0)`$, so that we can assume $`\|w_{13}^{(r)}\|<\|w_{16}\|`$. In a calculation of the baryon magnetic moments to fourth chiral order $`w_{16}`$ has been determined to be $$w_{16}0.4\text{GeV}^1.$$ (59) Inserting this upper limit for $`w_{13}^r`$ one obtains $`\|d_n^{\gamma (tree)}\|`$ $`<`$ $`9.6\times 10^{16}\theta _0e\text{cm}`$ $`\|d_\mathrm{\Lambda }^{\gamma (tree)}\|`$ $`<`$ $`4.8\times 10^{16}\theta _0e\text{cm}.`$ (60) Since we have taken quite conservative limits, the tree level contributions could be dominating if the extreme values are chosen. A more realistic estimate might be obtained by setting $`\|w_{13}^{(r)}/w_{16}\|\frac{1}{3}`$ for $`N_c=3`$, so that $`\|d_n^{\gamma (tree)}\|`$ $``$ $`3.2\times 10^{16}\theta _0e\text{cm}`$ $`\|d_\mathrm{\Lambda }^{\gamma (tree)}\|`$ $``$ $`1.6\times 10^{16}\theta _0e\text{cm}.`$ (61) While the chiral logarithms dominate for $`d_n^\gamma `$, the tree level contribution can be substantial for the $`\mathrm{\Lambda }`$. This leads to $`d_n^\gamma =(7.5\pm 3.2)\times 10^{16}\theta _0e\text{cm}`$ $`d_\mathrm{\Lambda }^\gamma =(1.7\pm 1.6)\times 10^{16}\theta _0e\text{cm},`$ (62) where the theoretical uncertainty is given by the estimate of the tree level contribution in Eq. (5). From a comparison of the central value for $`d_n^\gamma `$ with the experimental upper limit in Eq. (1) we obtain $$\|\theta _0\|<8.4\times 10^{11}.$$ (63) In the case of the $`\mathrm{\Lambda }`$ the experimental constraint is given by $$d_\mathrm{\Lambda }^\gamma <1.5\times 10^{16}e\text{cm}$$ (64) which leads to $$\|\theta _0\|<0.9.$$ (65) We have shown that it is possible to obtain a reliable limit for the vacuum angle $`\theta _0`$ by calculating the electric dipole moment of the neutron within the framework of heavy baryon chiral perturbation theory. To this end, we have constructed the most general effective Lagrangian up to one-loop order in the presence of the vacuum angle $`\theta _0`$ with the method proposed in . The theoretical uncertainty from unknown parameters at tree level has been estimated by using large $`N_c`$ arguments. While the chiral loops dominate for the electric dipole moment of the neutron, the counterterms can be substantial in the case of the $`\mathrm{\Lambda }`$. ## Acknowledgments The author thanks N. Kaiser for helpful discussions and careful reading of the manuscript. ## Figure captions 1. Shown are the tree diagrams for the electric dipole moment. Solid and dashed lines denote baryons and pseudoscalar mesons, respectively. The wavy line represents a photon and the dot is a $`CP`$-violating vertex. 2. Loop diagrams contributing to the electric dipole moment. Solid and dashed lines denote baryons and pseudoscalar mesons, respectively. The wavy line represents a photon and the dot is a $`CP`$-violating vertex. Figure 1 Figure 2
warning/0004/hep-ph0004267.html	ar5iv	text	# DESY 00-068 ISSN 0418-9833 hep-ph/0004267 April 2000 Order 𝛼³⁢ln(1/𝛼) Corrections to Positronium Decays ## Abstract The logarithmically enhanced $`\alpha ^3\mathrm{ln}(1/\alpha )`$ corrections to the para- and orthopositronium decay widths are calculated in the framework of dimensionally regularized nonrelativistic quantum electrodynamics. In the case of parapositronium, the correction is negative, approximately doubles the effect of the leading logarithmic $`\alpha ^3\mathrm{ln}^2(1/\alpha )`$ one, and is comparable to the nonlogarithmic $`O(\alpha ^2)`$ one. As for orthopositronium, the correction is positive and almost cancels the $`\alpha ^3\mathrm{ln}^2(1/\alpha )`$ one. The uncertainties in the theoretical predictions for the decay widths are reduced. PACS numbers: 12.20.Ds, 31.30.Jv, 36.10.Dr Positronium (Ps), which is an electromagnetic bound state of the electron $`e^{}`$ and the positron $`e^+`$, is the lightest known atom. Since its theoretical description is not plagued by strong-interaction uncertainties, thanks to the smallness of the electron mass $`m_e`$ relative to typical hadronic mass scales, its properties can be calculated perturbatively in quantum electrodynamics (QED), as an expansion in Sommerfeld’s fine-structure constant $`\alpha `$, with very high precision. Ps is thus a unique laboratory for testing the QED theory of weakly bound systems. The decay widths of the $`{}_{}{}^{1}S_{0}^{}`$ parapositronium (p-Ps) and $`{}_{}{}^{3}S_{1}^{}`$ orthopositronium (o-Ps) ground states to two and three photons, respectively, have been the subject of a vast number of theoretical and experimental investigations. The present theoretical knowledge may be summarized as: $`\mathrm{\Gamma }_p^{\mathrm{th}}`$ $`=`$ $`\mathrm{\Gamma }_p^{(0)}[1+A_p{\displaystyle \frac{\alpha }{\pi }}+2\alpha ^2\mathrm{ln}{\displaystyle \frac{1}{\alpha }}+B_p\left({\displaystyle \frac{\alpha }{\pi }}\right)^2`$ (1) $`{\displaystyle \frac{3\alpha ^3}{2\pi }}\mathrm{ln}^2{\displaystyle \frac{1}{\alpha }}+C_p{\displaystyle \frac{\alpha ^3}{\pi }}\mathrm{ln}{\displaystyle \frac{1}{\alpha }}+D_p\left({\displaystyle \frac{\alpha }{\pi }}\right)^3+\mathrm{}],`$ $`\mathrm{\Gamma }_o^{\mathrm{th}}`$ $`=`$ $`\mathrm{\Gamma }_o^{(0)}[1+A_o{\displaystyle \frac{\alpha }{\pi }}{\displaystyle \frac{\alpha ^2}{3}}\mathrm{ln}{\displaystyle \frac{1}{\alpha }}+B_o\left({\displaystyle \frac{\alpha }{\pi }}\right)^2`$ (2) $`{\displaystyle \frac{3\alpha ^3}{2\pi }}\mathrm{ln}^2{\displaystyle \frac{1}{\alpha }}+C_o{\displaystyle \frac{\alpha ^3}{\pi }}\mathrm{ln}{\displaystyle \frac{1}{\alpha }}+D_o\left({\displaystyle \frac{\alpha }{\pi }}\right)^3+\mathrm{}],`$ where $`\mathrm{\Gamma }_p^{(0)}`$ $`=`$ $`{\displaystyle \frac{\alpha ^5m_e}{2}},`$ $`\mathrm{\Gamma }_o^{(0)}`$ $`=`$ $`{\displaystyle \frac{2(\pi ^29)\alpha ^6m_e}{9\pi }},`$ (3) are the lowest-order results. The $`O(\alpha )`$ coefficients in Eqs. (1) and (2) read $`A_p`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{4}}5,`$ $`A_o`$ $`=`$ $`\mathrm{10.286\hspace{0.17em}606}(10).`$ (4) The logarithmically enhanced $`\alpha ^2\mathrm{ln}(1/\alpha )`$ terms in Eqs. (1) and (2) have been obtained in Refs. , respectively. Recently, the nonlogarithmic $`O(\alpha ^2)`$ coefficients in Eqs. (1) and (2) have been found to be $`B_p`$ $`=`$ $`5.14(30),`$ $`B_o`$ $`=`$ $`44.52(26).`$ (5) Note that the light-by-light scattering diagrams have been omitted in Ref. . In the p-Ps calculation the diagrams of this type increase the coefficient $`B_p`$ by $`1.28`$, and, therefore, their contribution to the coefficient $`B_o`$ is assumed to be relatively small. The p-Ps (o-Ps) decays into four (five) photons, which are not included in Eq. (5), lead to an increase of the coefficient $`B_p`$ ($`B_o`$) by $`0.274(1)`$ ($`0.19(1)`$) . In $`O(\alpha ^3)`$, only the leading logarithmic $`\alpha ^3\mathrm{ln}^2(1/\alpha )`$ terms are known . Including all the terms known so far, we obtain for the p-Ps and o-Ps total decay widths $`\mathrm{\Gamma }_p^{\mathrm{th}}`$ $`=`$ $`7989.512(13)\mu \mathrm{s}^1,`$ (6) $`\mathrm{\Gamma }_o^{\mathrm{th}}`$ $`=`$ $`7.039943(10)\mu \mathrm{s}^1,`$ (7) where the given errors stem only from the coefficients $`B_p`$ and $`B_o`$ respectively and we postpone the discussion of total uncertainty of theoretical estimates to the end of the paper. The purpose of this letter is complete our knowledge of the logarithmically enhanced terms of $`O(\alpha ^3)`$ by providing the coefficients $`C_p`$ and $`C_o`$, in analytic form. We also give order-of-magnitude estimates of the unknown coefficients $`D_p`$ and $`D_o`$. On the experimental side, the present situation is not entirely clear. Recently, the Ann Arbor group measured the p-Ps width to be $$\mathrm{\Gamma }_p^{\mathrm{exp}}=7990.9(1.7)\mu \mathrm{s}^1,$$ (8) which agrees with Eq. (6) within the experimental error. However, in the case of o-Ps, their measurements , $`\mathrm{\Gamma }_o^{\mathrm{exp}}(\text{gas})`$ $`=`$ $`7.0514(14)\mu \mathrm{s}^1,`$ $`\mathrm{\Gamma }_o^{\mathrm{exp}}(\text{vacuum})`$ $`=`$ $`7.0482(16)\mu \mathrm{s}^1,`$ (9) exceed Eq. (7) by 8 and 5 experimental standard deviations, respectively. This apparent contradiction is known as the o-Ps lifetime puzzle. On the other hand, the Tokyo group found $`\mathrm{\Gamma }_o^{\mathrm{exp}}(\text{SiO}\text{2})`$ $`=`$ $`7.0398(29)\mu \mathrm{s}^1,`$ (10) which agrees with Eq. (7) within the experimental error. Leaving this aside, the o-Ps results from Ann Arbor could be considered as a signal of new physics beyond the standard model. However, a large number of exotic decay modes have already been ruled out . No conclusion on the o-Ps lifetime puzzle can be drawn until the experimental precision increases and the data become unambiguous. On the theoretical side, it is an urgent matter to improve the predictions of the Ps lifetimes as much as possible. Thus, one is faced with the task of analyzing the $`O(\alpha ^3)`$ corrections, which is extremely difficult, especially for o-Ps. However, there is a special subclass of the $`O(\alpha ^3)`$ corrections which can be analyzed separately, namely those which are enhanced by powers of $`\mathrm{ln}(1/\alpha )5`$. They may reasonably be expected to provide an essential part of the $`O(\alpha ^3)`$ corrections. This may be substantiated by considering Eqs. (1) and (2) in $`O(\alpha ^2)`$, where logarithmic terms enter for the first time. In the case of p-Ps, 98% of the $`O(\alpha ^2)`$ correction stems from the logarithmic term. In the case of o-Ps the logarithmic term is not so dominant but still gives about $`1/4`$ of the total $`O(\alpha ^2)`$ correction. The origin of the logarithmic corrections is the presence of several scales in the bound-state problem. The dynamics of the nonrelativistic (NR) $`e^+e^{}`$ pair near threshold involves four different scales : (i) the hard scale (energy and momentum scale like $`m_e`$); (ii) the soft scale (energy and momentum scale like $`\beta m_e`$); (iii) the potential scale (energy scales like $`\beta ^2m_e`$, while momentum scales like $`\beta m_e`$); and (iv) the ultrasoft (US) scale (energy and momentum scale like $`\beta ^2m_e`$). Here $`\beta `$ denotes the electron velocity in the center-of-mass frame. The logarithmic integration over a loop momentum between different scales yields a power of $`\mathrm{ln}(1/\beta )`$. Since Ps is approximately a Coulomb system, we have $`\beta \alpha `$. This explains the appearance of powers of $`\mathrm{ln}(1/\alpha )`$ in Eqs. (1) and (2). The leading logarithmic corrections may be obtained straightforwardly by identifying the regions of logarithmic integration . The calculation of the subleading logarithms is much more involved because certain loop integrations must be performed exactly beyond the logarithmic accuracy. In the following, we briefly outline the main features of our analysis. We work in NR QED (NRQED) , which is the effective field theory that emerges by expanding the QED Lagrangian in $`\beta `$ and integrating out the hard modes. If we also integrate out the soft modes and the potential photons, we arrive at the effective theory of potential NRQED (pNRQED) , which contains potential electrons and US photons as active particles. Thus, the dynamics of the NR $`e^+e^{}`$ pair is governed by the effective Schrödinger equation and by its multipole interaction with the US photons. The corrections from harder scales are contained in the higher-dimensional operators of the NR Hamiltonian, corresponding to an expansion in $`\beta `$, and in the Wilson coefficients, which are expanded in $`\alpha `$. In the process of scale separation, spurious infrared (IR) and ultraviolet (UV) divergences arise, which endow the operators in the NR Hamiltonian with anomalous dimensions. We use dimensional regularization (DR), with $`d=42ϵ`$ space-time dimensions, to handle these divergences . This has the advantage that contributions from different scales are matched automatically. The logarithmic corrections are closely related to the anomalous dimensions and can be found by analyzing the divergences of the NR effective theory. In this way, we have obtained the leading logarithmic third-order corrections to the energy levels and wave functions at the origin of heavy quark-antiquark bound states , which includes the QED result as a special case. Here, we extend this approach to the subleading logarithms in QED. Note that the NRQED approach, endowed with an explicit momentum cutoff and a fictitious photon mass to regulate the UV and IR divergences, has also been applied to find the third-order correction, including subleading logarithms, to the hyperfine splitting in muonium . The annihilation of Ps is the hard process which gives rise to imaginary parts in the local operators of the NR Hamiltonian . The decay width can be obtained by averaging these operators over the bound-state wave function. The hard-scale corrections, which require fully relativistic QED calculations and are most difficult to find, do not depend on $`\beta `$ and do not lead to logarithmic contributions by themselves. However, they can interfere with the logarithmic corrections from the softer scales. Thus, the only results from relativistic perturbation theory that enter our analysis are (i) the one-loop hard renormalizations of the imaginary parts of the leading four-fermion operators, i.e., the Born decay amplitudes, which are given by the coefficients $`A_p`$ and $`A_o`$, and (ii) the hard parts of the one-loop $`O(\alpha \beta ^2)`$ operators . The missing ingredients can all be obtained in the NR approximation. These include (i) the correction to the Ps ground-state wave function at the origin due to the $`O(\alpha \beta ^2)`$ terms in the NR Hamiltonian, (ii) the $`O(\alpha \beta ^2)`$ and $`O(\alpha ^2\beta )`$ corrections to the leading four-fermion operators, and (iii) the correction due to the emission and absorption of US photons by the Ps bound state. The value of the ground-state ($`n=1`$) wave function at the origin $`\psi _1(0)`$ may be extracted from the NR Green function $`G(𝐱,𝐲,E)`$, which satisfies the equation $$\left(_C+\mathrm{\Delta }E\right)G(𝐱,𝐲,E)=\delta ^{(3)}(𝐱𝐲),$$ (11) where $`_C`$ is the Coulomb Hamiltonian and $`\mathrm{\Delta }`$ stands for the terms of higher orders in $`\alpha `$ and $`\beta `$. The solution of Eq. (11) can be found in time-independent perturbation theory as an expansion in $`\alpha `$ around the leading-order Coulomb Green function. We thus obtain the correction $`\mathrm{\Delta }\psi _1^2`$ in the relationship $`\|\psi _1(0)\|^2=\left\|\psi _1^C(0)\right\|^2\left(1+\mathrm{\Delta }\psi _1^2\right)`$, where $`\psi _1^C(0)`$ is the ground-state wave function at the origin in the Coulomb approximation. As mentioned above, this analysis may be enormously simplified by the use of DR. Proceeding along the lines of Ref. , we thus recover with ease the well-known $`\alpha ^2\mathrm{ln}(1/\alpha )`$ terms in Eqs. (1) and (2) , $$\mathrm{\Delta }^{}\psi _1^2=\alpha ^2\mathrm{ln}\frac{1}{\alpha }\left[2\frac{7}{6}S(S+1)\right],$$ (12) where $`S`$ is the eigenvalue of the total-spin operator $`𝐒`$. As mentioned above, $`\mathrm{\Delta }^{}\psi _1^2`$ interferes with the one-loop hard renormalization of the Born amplitudes, $`A_p`$ and $`A_o`$, to produce $`\alpha ^3\mathrm{ln}(1/\alpha )`$ terms. The resulting contributions to the coefficients $`C_p`$ and $`C_o`$ read $`2A_p`$ and $`A_o/3`$, respectively. The generic logarithmic $`O(\alpha ^3)`$ correction $`\mathrm{\Delta }^{\prime \prime }\psi _1^2`$ to $`\|\psi _1(0)\|^2`$ is generated by the following one-loop operators, given in the momentum representation with $`O(ϵ)`$ accuracy, $`\mathrm{\Delta }_\mathrm{h}`$ $`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle \frac{\alpha ^2}{m_e^2}}[{\displaystyle \frac{1}{\widehat{ϵ}}}\left({\displaystyle \frac{\mu ^2}{m_e^2}}\right)^ϵ+{\displaystyle \frac{39}{5}}12\mathrm{ln}2`$ (13) $`+({\displaystyle \frac{32}{3}}+6\mathrm{ln}2)𝐒^2],`$ $`\mathrm{\Delta }_\mathrm{s}`$ $`=`$ $`{\displaystyle \frac{7}{3}}{\displaystyle \frac{\alpha ^2}{m_e^2}}\left[{\displaystyle \frac{1}{\widehat{ϵ}}}\left({\displaystyle \frac{\mu ^2}{𝐪^2}}\right)^ϵ{\displaystyle \frac{1}{7}}\right],`$ (14) $`\mathrm{\Delta }_{\mathrm{us}}`$ $`=`$ $`{\displaystyle \frac{8}{3}}{\displaystyle \frac{\alpha ^2}{m_e^2}}\left[{\displaystyle \frac{1}{\widehat{ϵ}}}\left({\displaystyle \frac{\mu }{𝐩^2/m_eE_1^C}}\right)^{2ϵ}+{\displaystyle \frac{5}{3}}2\mathrm{ln}2\right],`$ (15) where $`1/\widehat{ϵ}=1/ϵ\gamma _E+\mathrm{ln}(4\pi )`$, with $`\gamma _E`$ being Euler’s constant, $`\mu `$ is the ’t Hooft mass scale of DR, $`𝐪`$ is the three-momentum transfer, and $`E_1^C=\alpha ^2m_e/4`$ is the Coulomb ground-state energy. Equations (13) and (14) give the hard and soft $`O(\alpha \beta ^2)`$ contributions to the NR Hamiltonian, respectively. The US contribution given in Eq. (15) arises from the emission and absorption of an US photon, which converts the on-shell Ps ground state into some off-shell state of the Coulomb spectrum, with three-momentum $`𝐩`$, before it decays. It is the only US contribution which can be represented by an operator of instantaneous interaction and thus give rise to logarithmic corrections. It has been found with the help of the method developed for the more complicated case of quantum chromodynamics in Ref. , where it was applied to the on-shell renormalization of the heavy-quarkonium wave function at the origin. The singularities of the operators in Eqs. (13)–(15) yield the logarithmic corrections which we are interested in. Up to their logarithmic dependences on $`𝐪^2`$ and $`𝐩^2`$, these operators are of the $`\delta `$-function type in coordinate space and, therefore, lead to additional singularities in the Coulomb Green function at the origin . As a consequence, in the evaluation of the Green function in time-independent perturbation theory from Eq. (11) with $`\mathrm{\Delta }_\mathrm{h}`$, $`\mathrm{\Delta }_\mathrm{s}`$, and $`\mathrm{\Delta }_{\mathrm{us}}`$, overlapping logarithmic divergences appear in the part of the first-order term which corresponds to the interference of the one-photon contribution to the Coulomb Green function and the first terms of Eqs. (13)–(15). This results in the double-logarithmic contribution, which can be directly extracted from the coefficient of the leading double-pole singularity . Since we are interested in the single-logarithmic contribution, we also have to keep the subleading terms in this analysis. The logarithmic corrections which are generated by the non-overlapping singularities can be obtained by putting $`\mu =m_e`$ in the Coulomb Green function at the origin and $`𝐪^2=𝐩^2=m_eE_1^C`$ in Eqs. (13)–(15), and proceeding as in the evaluation of Eq. (12). We thus obtain $`\mathrm{\Delta }^{\prime \prime }\psi _1^2`$ $`=`$ $`{\displaystyle \frac{\alpha ^3}{\pi }}\{{\displaystyle \frac{3}{2}}\mathrm{ln}^2{\displaystyle \frac{1}{\alpha }}+[{\displaystyle \frac{184}{45}}+{\displaystyle \frac{2}{3}}\mathrm{ln}2`$ (16) $`+({\displaystyle \frac{16}{9}}+\mathrm{ln}2)S(S+1)]\mathrm{ln}{\displaystyle \frac{1}{\alpha }}\}.`$ The first term herein agrees with the corresponding terms in Eqs. (1) and (2) , while the second one represents a new result. Another source of $`\alpha ^3\mathrm{ln}(1/\alpha )`$ terms is the $`O(\alpha \beta ^2)`$ corrections to the leading four-fermion operators. Since they do not involve the singular Coulomb Green function at the origin, there are no overlapping divergences, and we may simply read off the resulting $`\alpha ^3\mathrm{ln}(1/\alpha )`$ terms from the poles of their US parts, which are given by the operator $$\frac{1}{ϵ}\frac{2\alpha }{3\pi }\frac{𝐩^2+𝐩_{}^{}{}_{}{}^{2}}{m_e^2}V_4(𝐩,𝐩^{},𝐒).$$ (17) Here $`V_4(𝐩,𝐩^{},𝐒)`$ is the local four-fermion operator which generates the leading-order decay widths. Taking the expectation value of Eq. (17) w.r.t. the ground-state wave function, one encounters power-divergent integrals . They can be consistently treated within DR . This leads to the substitution $`𝐩^2,𝐩_{}^{}{}_{}{}^{2}m_eE_1^C`$ in the matrix element. The UV-pole contribution of Eq. (17) is then canceled by the IR pole of the hard contribution . This implies that the logarithmic integration ranges from the US scale $`\alpha ^2m_e`$ up to the hard scale $`m_e`$, so that $`1/ϵ`$ should be replaced by $`4\mathrm{ln}(1/\alpha )`$ . The resulting $`\alpha ^3\mathrm{ln}(1/\alpha )`$ corrections to the decay widths are spin-independent and read $$\mathrm{\Delta }\mathrm{\Gamma }_{p,o}=\mathrm{\Gamma }_{p,o}^{(0)}\frac{4\alpha ^3}{3\pi }\mathrm{ln}\frac{1}{\alpha }.$$ (18) The last source of $`\alpha ^3\mathrm{ln}(1/\alpha )`$ terms is the $`O(\alpha ^2\beta )`$ corrections to the leading four-fermion operators. These corrections are non-analytic in $`𝐩^2`$ and of the form $$\left[\frac{7}{6}+\left(8\frac{32}{3}\mathrm{ln}2\right)\right]\frac{\alpha ^2}{4}\frac{\|𝐩\|+\|𝐩^{}\|}{m_e}V_4(𝐩,𝐩^{},𝐒),$$ (19) where the first and second terms contained within the square brackets are the soft and US contributions, respectively. Although the coefficient of $`V_4(𝐩,𝐩^{},𝐒)`$ in Eq. (19) is finite, the matrix element of Eq. (19) between Ps bound-state wave functions is logarithmically divergent in the UV region. In DR, the divergent part is given by the matrix element of Eq. (19) with $`\|𝐩\|,\|𝐩^{}\|`$ replaced by $`m_e\alpha /(\pi ϵ)`$. In contrast to Eq. (17), the logarithmic integration now ranges from the soft scale up to the hard one, and $`1/ϵ`$ should be replaced by $`2\mathrm{ln}(1/\alpha )`$. The resulting contributions to the decay widths read $$\mathrm{\Delta }\mathrm{\Gamma }_{p,o}=\mathrm{\Gamma }_{p,o}^{(0)}\frac{\alpha ^3}{\pi }\mathrm{ln}\frac{1}{\alpha }\left(\frac{41}{6}\frac{32}{3}\mathrm{ln}2\right).$$ (20) Summing up the various $`\alpha ^3\mathrm{ln}(1/\alpha )`$ terms derived above, we obtain $`C_p`$ $`=`$ $`2A_p+{\displaystyle \frac{367}{90}}10\mathrm{ln}27.919,`$ $`C_o`$ $`=`$ $`{\displaystyle \frac{A_o}{3}}+{\displaystyle \frac{229}{30}}8\mathrm{ln}25.517.`$ (21) In the case of p-Ps, the new $`\alpha ^3\mathrm{ln}(1/\alpha )`$ term has the same sign as the $`\alpha ^3\mathrm{ln}^2(1/\alpha )`$ one and exceeds the latter in magnitude. The sum of these two terms compensates approximately 1/3 of the positive contribution from the nonlogarithmic $`O(\alpha ^2)`$ term. As for o-Ps, the new $`\alpha ^3\mathrm{ln}(1/\alpha )`$ term cancels approximately 3/4 of the $`\alpha ^3\mathrm{ln}^2(1/\alpha )`$ contribution. Our final predictions for the p-Ps and o-Ps total decay widths, including the multi-photon channels, read $`\mathrm{\Gamma }_p^{\mathrm{th}}`$ $`=`$ $`7989.620(13)\mu \mathrm{s}^1,`$ (22) $`\mathrm{\Gamma }_o^{\mathrm{th}}`$ $`=`$ $`7.039968(10)\mu \mathrm{s}^1,`$ (23) which has to be compared with Eqs. (6) and (7). As before, Eq. (22) agrees with the Ann Arbor measurement (8), and Eq. (23) favours the Tokyo measurement (10), while it significantly undershoots the Ann Arbor measurements (9). The missing nonlogarithmic $`O(\alpha ^3)`$ corrections in Eqs. (1) and (2) receive contributions from three-loop QED diagrams with a considerable number of external lines, which are far beyond the reach of presently available computational techniques. In this sense, we expect Eqs. (22) and (23) to remain the best predictions for the forseeable future. However, we may speculate about the magnitudes of the coefficients $`D_p`$ and $`D_o`$. Two powers of $`\alpha `$ in these terms can be of NR origin. Each of them should be accompanied by the characteristic factor $`\pi `$, which happens for the logarithmic terms. Thus, we estimate the coefficients $`D_p`$ and $`D_o`$ to be a few units times $`\pi ^2`$. This rule of thumb is in reasonable agreement with the situation at $`O(\alpha ^2)`$, where we have $`B_p\pi ^2/2`$ and $`B_o4\pi ^2`$. If the coefficients $`D_p`$ and $`D_o`$ do not have magnitudes in excess of 100, then the uncertainties due the lack of their knowledge falls within the errors quoted in Eqs. (22) and (23). Then, our new results reduce the uncertainties in the predicted p-Ps decay width to $`10^2\mu \mathrm{s}^1`$. The main remaining theoretical uncertainty in o-Ps decay width is related to the unknown $`O(\alpha ^2)`$ contribution of the light-by-light scattering diagrams which can be estimated as a few units times $`10^5\mu \mathrm{s}^1`$ on the basis of p-Ps result. Calculation of this contribution along with our present result will reduce the uncertainties in the predicted o-Ps decay width to $`10^5\mu \mathrm{s}^1`$. Further progress in our understanding the Ps lifetime problem crucially depends also on the reduction of the experimental errors, which now greatly exceed the theoretical ones. Finally, we note that the technique developed in this letter can also be applied to the calculation of the subleading logarithmic $`\alpha ^7\mathrm{ln}(1/\alpha )`$ terms for the Ps hyperfine splitting. This problem is of special interest because of the apparent discrepancy between the latest experimental data and the best theoretical predictions, which include the $`O(\alpha ^6)`$ corrections (see Ref. and references cited therein) and the leading logarithmic $`\alpha ^7\mathrm{ln}^2\alpha `$ term . We are indebted to R. Hill and G. P. Lepage for useful discussions and to K. Melnikov and A. Yelkhovsky for pointing out the relevance of the $`O(\alpha ^2\beta )`$ contribution from Eq. (19). This work was supported in part by the Deutsche Forschungsgemeinschaft under Contract No. KN 365/1-1, by the Bundesministerium für Bildung und Forschung under Contract No. 05 HT9GUA 3, and by the European Commission through the Research Training Network Quantum Chromodynamics and the Deep Structure of Elementary Particles under Contract No. ERBFMRXCT980194. The work of A.A.P. was supported in part by the Volkswagen Foundation under Contract No. I/73611. Note added After including the $`O(\alpha ^2\beta )`$ contribution from Eq. (19), which was missed in the previous version of this letter, we find agreement with the analytical results for the p-Ps and o-Ps decay widths of Ref. and with the numerical result for the o-Ps decay width of Ref. . However, this correction is of insignificant size and immaterial for our numerical estimates of the Ps decay widths.
warning/0004/cond-mat0004410.html	ar5iv	text	# The spatial correlations in the velocities arising from a random distribution of point vortices. ## I Introduction Recently, several papers have focused on the statistics of velocity and velocity gradients produced by a random distribution of point vortices in two dimensions . This problem was first considered by Min et al. and independently by Jiménez . They showed, using different methods, that the velocity p.d.f. are Gaussian (but with a slow convergence) while the distribution of velocity gradients follow a Cauchy law. Their theoretical results were confirmed by numerical simulations of point vortices and by Direct Navier Stokes simulations of 2D decaying turbulence when the flow becomes dominated by a large number of coherent vortices . Additional simulations of point vortex systems were performed by Weiss et al. who emphasized the importance of vortex pairs in the dynamics and their role in the tail of the velocity p.d.f. They showed in particular that these pairs are responsible for anomalous diffusion and proposed to model the motion of the vortices by an Ornstein-Uhlenbeck stochastic process combined with Lévy walks. A formal analogy exists between the statistics of the velocity fluctuations due to a collection of point vortices and the statistics of the gravitational field produced by a random distribution of stars. The only difference, of great importance, is the space dimensionality $`D=2`$ instead of $`D=3`$. In a long series of papers, Chandrasekhar and von Neumann analyzed in detail the distribution of the gravitational field arising from a random distribution of stars. Their study was inspired by the work of Holtsmark concerning the fluctuations of the electric field in a gas composed of simple ions and by the work of von Smoluchowski concerning the persistence of fluctuations in the Brownian theory. Chandrasekhar and Von Neumann rederived the Holtsmark distribution for $`𝐅`$, the gravitational field, and determined many other statistical measures for the correlations of $`𝐅`$ or the joint distribution of $`𝐅`$ and $`d𝐅/dt`$. Their motivation was to derive the expression for the speed of fluctuations and the diffusion coefficient of stars in a purely stochastic framework. Likewise in the vortex problem, the formation of binary stars alters the results of the statistical analysis at large field strengths. The beautiful work of Chandrasekhar and von Neumann is an imposing “tour de force” and provides detailed mathematical methods for analyzing the statistics of fluctuations in physics and astronomy. The close connexion between stellar and vortex systems has been investigated by one of us in a series of papers and it was natural to consider the extension of the Chandrasekhar-von Neumann analysis to the case of point vortices. Therefore, in Ref. we analyzed in some details the statistical features of the stochastic velocity field produced by a random distribution of point vortices in 2D turbulence. We rederived the results of Min et al. and Jiménez for the distribution of velocity $`𝐕`$ and acceleration $`𝐀`$ (closely related to the velocity gradients) using the methods of Chandrasekhar and von Neumann. This formalism is very powerful and allowed us to obtain generalizations of previous works: we considered the distributions of the vectors $`𝐕`$ and $`𝐀`$ (not only the projection along an axis), we gave expressions for these distributions valid for an arbitrary number of vortices (for finite $`N`$ these distributions can be evaluated numerically) and we extended the analysis to an arbitrary spectrum of circulations among the vortices and to the case of non singular vortex “blobs”. We also considered for the first time the joint distribution of velocity and acceleration $`W(𝐕,𝐀)`$, determined the typical duration of the velocity fluctuations $`T(V)`$ and proposed an expression for the diffusion coefficient $`D`$ of point vortices thereby justifying the phenomenological result of Weiss et al. . These theoretical results were found to be in good agreement with numerical simulations and experiments of 2D decaying turbulence. They are also of great importance to build up a rational kinetic theory of point vortices . The motivation of the present paper is to characterize the spatial correlations in the velocities occurring at two points separated by an arbitrary distance. This problem is customary in turbulence but, in general, exact results are difficult to obtain. It is an interest of our model to allow for nice analytical solutions. Like in Ref. , we consider a collection of $`N`$ point vortices with circulation $`\gamma `$ randomly distributed in a disk of radius $`R`$. We assume that the vortices have a spatial Poisson distribution, i.e. their positions are independent and uniformly distributed over the entire domain. We are particularly interested in the “thermodynamical limit” in which the number of vortices and the size of the domain go to infinity ($`N\mathrm{}`$, $`R\mathrm{}`$) in such a way that the vortex density $`n=\frac{N}{\pi R^2}`$ remains finite. In this limit, the Poisson distribution is shown to be stationary . Therefore, if the vortices are initially uniformly distributed they will remain uniformly distributed in average during all the subsequent evolution. This property has been checked numerically by Jiménez using Direct Navier Stokes simulations of 2D decaying turbulence. In the statistical theory of point vortices initiated by Onsager and further developed by Joyce & Montgomery and Lundgren & Pointin , the Poisson distribution corresponds to a structureless equilibrium state with inverse temperature $`\beta =0`$ (when the angular momentum is not conserved) or $`\beta +\mathrm{}`$ (when the angular momentum is conserved). Of course, more general initial conditions are possible and lead to equilibrium states with $`\beta <0`$ in which the vortices are clustered in “macrovortices” . The statistics of fluctuations remain the same in these more general situations but the vortices are expected to experience a systematic drift (Chavanis, 1998) in addition to their diffusive motion, due to the inhomogeneity of the vortex cloud. It is in this context that a kinetic theory of point vortices, consistent with Onsager approach, can be constructed . At equilibrium the drift balances the scattering and maintains nontrivial density distributions. We shall restrict ourselves, however, in the present article to the case of a uniform distribution of point vortices for which the drift cancels out. In section II, we recall known results concerning the distribution of velocities $`W(𝐕)`$ occuring at a fixed point (equations (10)(11) and (12)). This is essentially to set the notations that will be used in the sequel. Since the velocity distribution created by point vortices is intermediate between Gaussian and Lévy laws, we shall call it the “marginal Gaussian distribution”. We also consider the case of non singular “vortex blobs” with a core of size $`a`$. We write down the exact characteristic function of $`W(𝐕)`$ valid for all velocities and all core sizes (equation (20)). For “extended” vortices, we prove that the velocity distribution is Gaussian (with no tail) in agreement with the numerical observations of Jiménez and Bracco et al. . For $`a=0`$ we recover the point vortex limit for which the velocity distribution has a Gaussian core and an algebraic tail. For “small” non singular vortices ($`a0`$ but $`a0`$), we argue that the form of our characteristic function can explain the occurence of almost exponential tails observed by Jiménez and Bracco et al. . In section III we analyze the joint distribution of velocity and velocity gradient $`W(𝐕,\delta 𝐕)`$. In subsection III B, we rederive the Cauchy distribution for $`\delta 𝐕`$ using the method of Chandrasekhar and von Neumann (equation (52)). In subsection III C, we generalize this method to the case of “vortex blobs” and find the explicit characteristic function of $`W(\delta 𝐕)`$ valid for all velocities and core sizes (equation (57)). We find that the distribution is Cauchy for small fluctuations and Gaussian for large fluctuations. This is in agreement with the asymptotic behaviours found by Min et al.. It is likely that in between the distribution passes through an exponential tail as observed numerically in Ref. . In subsection III D, we determine an exact expression for the moment $`\delta 𝐕_𝐕`$, the average value of $`\delta 𝐕`$ for a given velocity $`𝐕`$ (equations (85) and (86)). When $`V+\mathrm{}`$, our results have a clear physical meaning in the nearest neighbor approximation. In this approximation, we determine the $`n`$-th conditional moments $`(\delta 𝐕)^n_𝐕`$ of the velocity gradients. In section IV, we turn to the more difficult problem concerning the velocity correlations between two points separated by a finite distance. Special attention is devoted to the bivariate distribution $`W(𝐕_0,𝐕_1)`$ that a velocity $`𝐕_0`$ occurs at $`O`$ and that, simultaneously, a velocity $`𝐕_1`$ occurs at a point separated from the first by a distance $`𝐫_1`$. In subsection IV B, we determine an exact but complicated expression for the moment $`𝐕_1_{𝐕_0}`$ (equations (111)(118) and (122)). More explicit results are obtained in section V when this quantity is averaged over all mutual orientations of $`𝐫_1`$ and $`𝐕_0`$. An exact expression is found for $`𝐕_1_{V_0}`$ in subsection V B (equations (140)(142)). In subsection V C, we give the asymptotic behaviour of this quantity when $`r_1+\mathrm{}`$ (equations (156)(157)). In section VI, we consider the projection of $`𝐕_1_{𝐕_0}`$ in the direction of $`𝐕_0`$ and its average over all mutual orientations of $`𝐫_1`$ and $`𝐕_0`$. This defines a quantity $`V_{1\|\|}_{V_0}`$ (subsection VI A) whose expression is given by equation (165). In subsection VI B, we give the asymptotic behaviours of this quantity in the limits $`r_10`$ (equation (175)) and $`r_1+\mathrm{}`$ (equations (178)(179)). When $`V_{1\|\|}_{V_0}`$ is further averaged over all possible values of $`V_0`$, this determines a function of $`r_1`$ alone which characterizes the correlations in the velocities occurring simultaneously at two points separated by a distance $`𝐫_1`$ (subsection VI C). The exact expression of $`V_{1\|\|}`$ is given by equation (191) and its asymptotic behaviours for $`r_10`$ and $`r_1+\mathrm{}`$ by equations (195)(196). In subsection VI D we calculate the correlation function $`𝐕_0𝐕_1/V_0^2`$. Its general expression is given by equation (200) and its asymptotic behaviours by equations (206)(207). Finally, in section VII, we calculate the spatial velocity autocorrelation function and the energy spectrum of point vortices. We find in subsection VII A that the spatial velocity correlation function decays extremely slowly with the distance (equation (216)). The simple result $`𝐕_0𝐕_1=\frac{n\gamma ^2}{2\pi }\mathrm{ln}(R/r_1)`$, also derived more directly in Appendix D, does not seem to have been given previously. In subsection VII B we use this correlation function to determine the energy spectrum $`E(k)`$ of point vortices (equation (223)). This offers an alternative to the method considered by Novikov . When $`k+\mathrm{}`$, we recover the classical result $`E(k)k^1`$ of Novikov and for $`k0`$ we find $`E(k)k`$. Note that the spatial velocity autocorrelation function $`𝐕_0𝐕_1`$ diverges logarithmically at small separations since the variance of the velocity is not defined. This is the reason why we consider more regular quantities such as $`V_{1\|\|}𝐕_0𝐕_1/V_0`$ and $`𝐕_0𝐕_1/V_0^2`$ that are finite at small separations. The computation of these quantities is difficult and is a important part of our work. ## II The statistics of velocity fluctuations occurring at a fixed point We recall in this section the results concerning the distribution of velocity fluctuations occuring at a fixed point. This problem was first considered by Min et al. and Jiménez and later on by Weiss et al. , Chukbar , Kuvshinov & Schep and Chavanis & Sire . Min et al. and Weiss et al. focus on the mathematical limit $`N+\mathrm{}`$ and claim that the velocity distribution is Gaussian by invoking a generalization of the Central Limit Theorem due to Ibragimov & Linnik (the proof is not straightforward because the variance of the velocity created by a single vortex diverges logarithmically). However, they stress that the convergence to the Gaussian is extremely slow and that this slow convergence is responsible for discrepencies at large velocities. This is confirmed by direct numerical simulations of point vortices . Jiménez and Chavanis & Sire adopt another point of view by treating $`N`$ as a large number but not $`\mathrm{ln}N`$. This seems to be more appropriate to physical situations where the typical number of point vortices does not exceed $`10^4`$. They proved in that limit that the velocity p.d.f. has a Gaussian core and an algebraic tail (see equations (10)(11) and (12)). When $`\mathrm{ln}N+\mathrm{}`$, the tail is rejected to infinity and the distribution is “purely” Gaussian in agreement with the theorem of Ibragimov and Linnik. All authors understood that the algebraic tail is produced by the nearest neighbor and found the correct exponent by a phenomenological approach. A mathematical justification of this result was given by Jiménez and Chavanis & Sire who explicitly derived the characteristic function (7) associated with the velocity distribution. The $`\rho ^2`$ factor in the characteristic function gives rise to the Gaussian core while the term in $`\mathrm{ln}\rho `$ gives rise to the algebraic tail . This last result is not straightforward and the rigorous proof requires to evaluate the integrals in the complex plane as done by Chavanis & Sire . The characteristic function (7) was also determined by Chukbar but he arbitrarily replaced the logarithmic factor appearing in (7) by a constant which is not justified. Therefore, he could not derive the algebraic tail. Kuvshinov & Schep (2000) wrote the right distribution but justified the algebraic tail by the phenomenological nearest neighbor argument not by a rigorous calculation. Below, we recall the main results concerning the velocity p.d.f. following the presentation and notations of Ref. . The velocity $`𝐕`$ occurring at the center $`O`$ of the domain is the sum of the velocities $`𝚽_i`$ $`(i=1,\mathrm{},N)`$ produced by the $`N`$ vortices: $$𝐕=\underset{i=1}{\overset{N}{}}𝚽_i$$ (1) $$𝚽_i=\frac{\gamma }{2\pi }\frac{𝐫_i}{r_i^2}$$ (2) where $`𝐫_i`$ denotes the position of the $`i^{th}`$ vortex and, by definition, $`𝐫_i`$ is the vector $`𝐫_i`$ rotated by $`+\frac{\pi }{2}`$. Since the vortices are randomly distributed, the velocity $`𝐕`$ fluctuates. In the limit when the number of vortices and the size of the domain go to infinity $`N,R\mathrm{}`$ in such a way that the density $`n=\frac{N}{\pi R^2}`$ remains finite, the distribution of the velocity $`𝐕`$ is given by $$W(𝐕)=\frac{1}{4\pi ^2}A(\rho \rho \rho )e^{i\rho \rho \rho 𝐕}d^2\rho \rho \rho $$ (3) with $$A(\rho \rho \rho )=e^{nC(\rho \rho \rho )}$$ (4) and $$C(\rho \rho \rho )=\frac{\gamma ^2}{4\pi ^2}_{\|\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\|=\frac{\gamma }{2\pi R}}^+\mathrm{}(1e^{i\rho \rho \rho \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }})\frac{1}{\mathrm{\Phi }^4}d^2\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }$$ (5) Introducing polar coordinates and performing the angular integration, equation (5) can be reduced to $$C(\rho \rho \rho )=\frac{\gamma ^2\rho ^2}{2\pi }_{\frac{\gamma \rho }{2\pi R}}^+\mathrm{}(1J_0(x))\frac{dx}{x^3}$$ (6) Because of logarithmic divergences when $`R+\mathrm{}`$, the “thermodynamical limit” does not exist. Therefore, (3) must be considered as an equivalent of $`W_N(𝐕)`$ for large $`N`$’s not a true limit. In Ref. , it was found that $$C(\rho \rho \rho )=\frac{\gamma ^2\rho ^2}{16\pi }\mathrm{ln}\left(\frac{4\pi N}{n\gamma ^2\rho ^2}\right)$$ (7) For $`\rho >0`$ and $`N\mathrm{}`$, we have $$C(\rho \rho \rho )=\frac{\gamma ^2\rho ^2}{16\pi }\mathrm{ln}N$$ (8) and for $`\rho 0`$, we obtain $$C(\rho \rho \rho )=\frac{\gamma ^2\rho ^2}{8\pi }\mathrm{ln}\rho $$ (9) The velocity distribution (3) is just the Fourier transform of (4) with (8)-(9). It is found that (see Ref. for detailed calculations): $$W(𝐕)=\frac{4}{n\gamma ^2\mathrm{ln}N}e^{\frac{4\pi }{n\gamma ^2\mathrm{ln}N}V^2}(VV_{crit}(N))$$ (10) $$W(𝐕)\frac{n\gamma ^2}{4\pi ^2V^4}(VV_{crit}(N))$$ (11) where $$V_{crit}(N)\left(\frac{n\gamma ^2}{4\pi }\mathrm{ln}N\right)^{1/2}\mathrm{ln}^{1/2}(\mathrm{ln}N)$$ (12) As discussed in Ref. , the velocity p.d.f behaves in a manner which is intermediate between Gaussian and Lévy laws. This is because the variance $`\mathrm{\Phi }^2`$ of the velocity created by a single vortex diverges logarithmically. Therefore, the Central Limit Theorem is only marginally applicable: the core of the velocity distribution is Gaussian while the tail decays algebraically as for a Lévy law. Therefore, the distribution function (10)(11) can be called a “marginal Gaussian distribution”. The velocity $`V_{crit}(N)`$ marks the transition between the Gaussian core and the algebraic tail. This algebraic tail can be given a simple interpretation in the nearest neighbor approximation (see, e.g., Ref. ). It is remarkable that the velocity created by the nearest neighbor $$V_{n.n.}^2\left(\frac{\gamma }{2\pi d}\right)^2\frac{\gamma ^2}{4\pi ^2}\frac{N}{\pi R^2}$$ (13) is precisely of the same order (up to some logarithmic corrections) as the typical velocity created by the rest of the system $$V^2N\frac{\gamma ^2}{4\pi ^2r^2}N_{\|𝐫\|=d}^R\tau (𝐫)\frac{\gamma ^2}{4\pi ^2r^2}d^2𝐫\frac{\gamma ^2}{4\pi }\frac{N}{\pi R^2}\mathrm{ln}N$$ (14) where $`dn^{1/2}`$ denotes the interparticle distance and $`\tau (𝐫)=1/\pi R^2`$ the probability of occurrence of a vortex in $`𝐫`$. In a sense, we can consider that the velocity is dominated by the contribution of the nearest neighbor and that collective effects are responsible for logarithmic corrections. The variance of the velocity can be written as $$V^2=\frac{N}{\pi R^2}_0^+\mathrm{}\frac{\gamma ^2}{4\pi ^2r^2}2\pi r𝑑r$$ (15) The integral (15) diverges logarithmically at both small and large separations. This is consistent with the distribution (10)(11) and the formula $$V^2=_0^+\mathrm{}W(𝐕)V^22\pi V𝑑V$$ (16) Because of the algebraic tail, the integral (16) diverges logarithmically as $`V+\mathrm{}`$ (i.e $`r0`$). In addition, if we were to extend the Gaussian distribution (10) to all velocities we would find that the “variance” $$V^2=\frac{n\gamma ^2}{4\pi }\mathrm{ln}N$$ (17) diverges logarithmically when $`N+\mathrm{}`$ (i.e $`R+\mathrm{}`$). The first moment of the velocity defined by $$\|𝐕\|=_0^+\mathrm{}W(𝐕)V2\pi V𝑑V$$ (18) diverges logarithmically as $`N\mathrm{}`$ but, unlike $`V^2`$, converges for $`V+\mathrm{}`$. To leading order in $`\mathrm{ln}N`$, we have: $$\|𝐕\|=\left(\frac{n\gamma ^2}{16}\mathrm{ln}N\right)^{\frac{1}{2}}$$ (19) If we account for a spectrum of circulations among the vortices, the previous results are maintained with $`\overline{\gamma ^2}`$ in place of $`\gamma ^2`$ . In particular, for a neutral system consisting in an equal number of vortices with circulation $`+\gamma `$ and $`\gamma `$, the results are unchanged. This is to be expected as a vortex with circulation $`\gamma `$ located in $`𝐫`$ produces the same velocity as a vortex with circulation $`+\gamma `$ located in $`𝐫`$. Since the vortices are randomly distributed with uniform probability, the two populations are statistically equivalent. When the system is non neutral (but still homogeneous), the previous results derived at the center of the domain remain valid at any point provided that the velocity is replaced by the fluctuating velocity $`𝒱𝒱𝒱=𝐕𝐕`$ where $`𝐕=\frac{1}{2}n\overline{\gamma }𝐚_{}`$ is the average velocity in $`𝐚`$ corresponding to a solid rotation . The case of inhomogeneous distribution of point vortices is considered in . In reality, the vortices have a finite radius $`a`$ which is not necessarily small (vortex “blobs”). If we consider that the core radius $`a`$ acts as a lower cut-off, equation (6) must be replaced by $$C(\rho \rho \rho )=\frac{\gamma ^2\rho ^2}{2\pi }_{\frac{\gamma \rho }{2\pi R}}^{\frac{\gamma \rho }{2\pi a}}(1J_0(x))\frac{dx}{x^3}$$ (20) This expression can be further simplified by part integrations. When $`a`$ is sufficiently large, we can make the approximation $$C(\rho \rho \rho )\frac{\gamma ^2\rho ^2}{8\pi }\mathrm{ln}\left(\frac{R}{a}\right)$$ (21) and this proves that for “extended vortices” the velocity distribution is exactly Gaussian (with no tail) . This result is consistent with the numerical observations of Jiménez and Bracco et al.. For $`a=0`$, we recover the marginal Gaussian distribution (10)(11) with an algebraic tail. For “small” non singular vortices ($`a0`$ but $`a0`$), the characteristic function (20) is quadratic for $`\rho +\mathrm{}`$ and $`\rho 0`$ corresponding to a Gaussian velocity distribution for $`V0`$ and $`V+\mathrm{}`$ (this large velocity limit is purely formal since, physically, the distribution must be cut at the maximum allowable velocity $`V_{max}\gamma /4\pi a`$ achieved when two vortices are at distance $`2a`$ from each other). However, for intermediate values of $`\rho `$, $`C(\rho \rho \rho )`$ is not quadratic and depends on the value of $`a`$. This fact can probably explain the occurence of almost exponential tails observed by Jiménez and Bracco et al. in their simulations. ## III The velocity correlations between two neighboring points In this section, we analyze the spatial correlations of the velocity fluctuations. We first consider the correlations between two neighboring points separated by an infinitesimal distance $`\delta 𝐫`$. ### A The formal solution of the problem The difference between the velocities occurring at two points distant $`\delta 𝐫`$ from each other is given by: $$\delta 𝐕=\underset{i=1}{\overset{N}{}}\psi \psi \psi _i$$ (22) with $$\psi \psi \psi _i=\frac{\gamma }{2\pi }\left\{\frac{\delta 𝐫_{}}{r_i^2}\frac{2(𝐫_i\delta 𝐫)𝐫_i}{r_i^4}\right\}$$ (23) where we have assumed that one of the points is at the center of the domain. This is similar to the expression for the acceleration $`𝐀`$ with $`\delta 𝐫`$ replacing $`𝐯`$ (see Ref. , equations (4)(5)). The correlations in the velocities occurring between these points can be specified by the function $`W(𝐕,\delta 𝐕)`$ which gives the simultaneous probability of the velocity $`𝐕`$ and the velocity increment $`\delta 𝐕`$. A general expression for the bivariate probability $`W(𝐕,\delta 𝐕)`$ can be readily written down following Markov’s method outlined in Ref. , section II.A. We have: $$W_N(𝐕,\delta 𝐕)=\frac{1}{16\pi ^4}A_N(\rho \rho \rho ,\sigma \sigma \sigma )e^{i(\rho \rho \rho 𝐕+\sigma \sigma \sigma \delta 𝐕)}d^2\rho \rho \rho d^2\sigma \sigma \sigma $$ (24) with: $$A_N(\rho \rho \rho ,\sigma \sigma \sigma )=\left(_{\|𝐫\|=0}^R\tau (𝐫)e^{i(\rho \rho \rho \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }+\sigma \sigma \sigma \psi \psi \psi )}d^2𝐫\right)^N$$ (25) and $$\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }=\frac{\gamma }{2\pi }\frac{𝐫_{}}{r^2}$$ (26) $$\psi \psi \psi =\frac{\gamma }{2\pi }\left(\frac{\delta 𝐫_{}}{r^2}\frac{2(𝐫\delta 𝐫)𝐫_{}}{r^4}\right)$$ (27) If we now suppose that the vortices are uniformly distributed on average, then $$\tau (𝐫)=\frac{1}{\pi R^2}$$ (28) and equation (25) becomes $$A_N(\rho \rho \rho ,\sigma \sigma \sigma )=\left(\frac{1}{\pi R^2}_{\|𝐫\|=0}^Re^{i(\rho \rho \rho \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }+\sigma \sigma \sigma \psi \psi \psi )}d^2𝐫\right)^N$$ (29) Since $$\frac{1}{\pi R^2}_{\|𝐫\|=0}^Rd^2𝐫=1$$ (30) we can rewrite our expression for $`A_N(\rho \rho \rho ,\sigma \sigma \sigma )`$ in the form $$A_N(\rho \rho \rho ,\sigma \sigma \sigma )=\left(1\frac{1}{\pi R^2}_{\|𝐫\|=0}^R(1e^{i(\rho \rho \rho \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }+\sigma \sigma \sigma \psi \psi \psi )})d^2𝐫\right)^N$$ (31) We now consider the limit when the number of vortices and the size of the domain go to infinity in such a way that the density remains finite $$N\mathrm{},R\mathrm{},n=\frac{N}{\pi R^2}\mathrm{finite}$$ If the integral occurring in equation (31) increases less rapidly than $`N`$, then $$A(\rho \rho \rho ,\sigma \sigma \sigma )=e^{nC(\rho \rho \rho ,\sigma \sigma \sigma )}$$ (32) with $$C(\rho \rho \rho ,\sigma \sigma \sigma )=_{\|𝐫\|=0}^R(1e^{i(\rho \rho \rho \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }+\sigma \sigma \sigma \psi \psi \psi )})d^2𝐫$$ (33) To make progress in the evaluation of (33), we find it more convenient to take $`\mathrm{\Phi }`$$`\mathrm{\Phi }`$$`\mathrm{\Phi }`$ as a variable of integration instead of $`𝐫`$. The Jacobian of the transformation $`\{𝐫\}\{𝚽\}`$ is $$\left\|\left\|\frac{(𝐫)}{(𝚽)}\right\|\right\|=\frac{\gamma ^2}{4\pi ^2\mathrm{\Phi }^4}$$ (34) so that $$C(\rho \rho \rho ,\sigma \sigma \sigma )=\frac{\gamma ^2}{4\pi ^2}_{\|𝚽\|=\frac{\gamma }{2\pi R}}^+\mathrm{}(1e^{i(\rho \rho \rho \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }+\sigma \sigma \sigma \psi \psi \psi )})\frac{1}{\mathrm{\Phi }^4}d^2𝚽$$ (35) In equation (35), $`\psi `$$`\psi `$$`\psi `$ must be expressed in terms of $`\mathrm{\Phi }`$$`\mathrm{\Phi }`$$`\mathrm{\Phi }`$. Combining (26) and (27) we find that $$\psi \psi \psi =\frac{2\pi }{\gamma }\left\{\mathrm{\Phi }^2\delta 𝐫_{}+2(\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }_{}\delta 𝐫)\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\right\}$$ (36) ### B The Cauchy distribution for $`\delta 𝐕`$ In this subsection, we show that the unconditional distribution of the velocity fluctuations $`\delta 𝐕`$ is governed by a 2D Cauchy law. This result was previously given by Jiménez and Min et al. using the methods of Ibragimov & Linnik and Feller. On a technical point of view, our calculations differ from the previous authors in the sense that we consider the distribution of the vector $`\delta 𝐕`$ instead of its projection along a particular direction (the same remark applies to the distribution of $`𝐕`$ in section II). In addition, the method presented here can be extended to an arbitrary spectrum of circulations among the vortices and to the case of vortex “blobs” (subsection III C). According to equation (24), we clearly have $$W(\delta 𝐕)=\frac{1}{16\pi ^4}A(\rho \rho \rho ,\sigma \sigma \sigma )e^{i(\rho \rho \rho 𝐕+\sigma \sigma \sigma \delta 𝐕)}d^2\rho \rho \rho d^2\sigma \sigma \sigma d^2𝐕$$ (37) This gives the distribution of the velocity increment $`\delta 𝐕`$ with no condition on $`𝐕`$. Using the identity $$\delta (𝐱)=\frac{1}{(2\pi )^2}e^{i\rho \rho \rho 𝐱}d^2\rho \rho \rho $$ (38) the foregoing expression for $`W(\delta 𝐕)`$ reduces to $$W(\delta 𝐕)=\frac{1}{4\pi ^2}A(\sigma \sigma \sigma )e^{i\sigma \sigma \sigma \delta 𝐕}d^2\sigma \sigma \sigma $$ (39) where we have written $`A(\sigma \sigma \sigma )`$ for $`A(\mathrm{𝟎},\sigma \sigma \sigma )`$. Hence, according to equations (32), (35) and (36) we obtain $$A(\sigma \sigma \sigma )=e^{nC(\sigma \sigma \sigma )}$$ (40) with $$C(\sigma \sigma \sigma )=\frac{\gamma ^2}{4\pi ^2}_{\|\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\|=0}^+\mathrm{}(1e^{i\sigma \sigma \sigma \psi \psi \psi })\frac{1}{\mathrm{\Phi }^4}d^2\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }$$ (41) and $$\sigma \sigma \sigma \psi \psi \psi =\frac{2\pi }{\gamma }\left\{\mathrm{\Phi }^2\sigma \sigma \sigma \delta 𝐫_{}+2(\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }_{}\delta 𝐫)(\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\sigma \sigma \sigma )\right\}$$ (42) Following the usual prescription, we have let $`R\mathrm{}`$ in (41) since the integral is convergent when $`\mathrm{\Phi }0`$. To evaluate this integral, we shall first introduce a system of coordinates with the $`x`$-axis in the direction of $`\delta 𝐫`$. Let us denote by $`\theta `$ and $`\beta `$ the angles that $`\mathrm{\Phi }`$$`\mathrm{\Phi }`$$`\mathrm{\Phi }`$ and $`\sigma `$$`\sigma `$$`\sigma `$ form with $`\delta 𝐫`$. Equation (42) now becomes $$\sigma \sigma \sigma \psi \psi \psi =\frac{2\pi }{\gamma }\sigma \mathrm{\Phi }^2\|\delta 𝐫\|\mathrm{sin}(2\theta \beta )$$ (43) and we obtain $$C(\sigma \sigma \sigma )=\frac{\gamma ^2}{4\pi ^2}_0^{2\pi }_0^+\mathrm{}\left(1\mathrm{cos}\left[\frac{2\pi }{\gamma }\sigma \mathrm{\Phi }^2\|\delta 𝐫\|\mathrm{sin}(2\theta )\right]\right)\frac{1}{\mathrm{\Phi }^3}𝑑\mathrm{\Phi }𝑑\theta $$ (44) With the identity $$_0^{2\pi }\mathrm{cos}(x\mathrm{sin}(2\theta ))𝑑\theta =2\pi J_0(x)$$ (45) the expression (44) for $`C(\sigma \sigma \sigma )`$ reduces to $$C(\sigma \sigma \sigma )=\frac{\gamma ^2}{2\pi }_0^+\mathrm{}\left(1J_0\left(\frac{2\pi }{\gamma }\sigma \mathrm{\Phi }^2\|\delta 𝐫\|\right)\right)\frac{1}{\mathrm{\Phi }^3}𝑑\mathrm{\Phi }$$ (46) With the change of variables $`x=\frac{2\pi }{\gamma }\sigma \|\delta 𝐫\|\mathrm{\Phi }^2`$, we obtain $$C(\sigma \sigma \sigma )=\frac{1}{2}\gamma \sigma \|\delta 𝐫\|_0^+\mathrm{}(1J_0(x))\frac{dx}{x^2}$$ (47) Integrating by parts and using the identities $$J_0^{}(x)=J_1(x)$$ (48) and $$_0^+\mathrm{}\frac{J_1(x)}{x}𝑑x=1$$ (49) we find $$C(\sigma \sigma \sigma )=\frac{1}{2}\gamma \|\delta 𝐫\|\sigma $$ (50) Hence $$A(\sigma \sigma \sigma )=e^{\frac{n\gamma }{2}\|\delta 𝐫\|\sigma }$$ (51) The distribution $`W(\delta 𝐕)`$ is just the Fourier transform of (51). This is the 2D-Cauchy distribution $$W(\delta 𝐕)=\frac{2}{\pi n^2\gamma ^2\|\delta 𝐫\|^2}\frac{1}{\left(1+\frac{4\|\delta 𝐕\|^2}{n^2\gamma ^2\|\delta 𝐫\|^2}\right)^{3/2}}$$ (52) The tail of the Cauchy distribution decreases algebraically like $$W(\delta 𝐕)\frac{n\gamma \|\delta 𝐫\|}{4\pi \|\delta 𝐕\|^3}$$ (53) The distribution of velocity increments $`W(\delta 𝐕)`$ is the same as the distribution of accelerations $`W(𝐀)`$. The only difference is the occurrence of the displacement $`\|\delta 𝐫\|`$ in place of the average velocity $`\overline{\|𝐯\|}`$ (see Ref. , equation (B30)). The Cauchy distribution is a particular Lévy law. Accordingly, the typical velocity increment produced by all the vortices $$\delta V^2N_{\|𝐫\|=d}^R\tau (𝐫)\frac{\gamma ^2}{4\pi ^2}\|\delta 𝐫\|^2\frac{1}{r^4}d^2𝐫\frac{\gamma ^2}{4\pi }\|\delta 𝐫\|^2\left(\frac{N}{\pi R^2}\right)^2$$ (54) is dominated by the contribution of the nearest neighbor $$\delta V_{n.n.}^2\frac{\gamma ^2}{4\pi ^2}\frac{\|\delta 𝐫\|^2}{d^4}\frac{\gamma ^2}{4\pi ^2}\|\delta 𝐫\|^2\left(\frac{N}{\pi R^2}\right)^2$$ (55) In addition, the algebraic tail (53) has a clear physical interpretation in the nearest neighbor approximation (see, e.g., Ref. , section V). Accounting for a spectrum of circulations among the vortices and following a procedure similar to that adopted in the Appendix B of Ref. , we find that the previous results are maintained provided that $`\gamma `$ is replaced by $`\overline{\|\gamma \|}`$. In particular, for a neutral system consisting in an equal number of vortices with circulation $`+\gamma `$ and $`\gamma `$, the results are unchanged. This is to be expected since a vortex with circulation $`\gamma `$ located in $`𝐫_{}`$ produces the same velocity increment as a vortex with circulation $`+\gamma `$ located in $`𝐫`$. ### C The distribution $`W(\delta 𝐕)`$ for vortex “blobs” If the vortices are not point-like but rather “blobs”, their size $`a`$ acts as a lower cut-off and equation (46) must be replaced by: $$C(\sigma \sigma \sigma )=\frac{\gamma ^2}{2\pi }_0^{\frac{\gamma }{2\pi a}}\left(1J_0\left(\frac{2\pi }{\gamma }\sigma \mathrm{\Phi }^2\|\delta 𝐫\|\right)\right)\frac{1}{\mathrm{\Phi }^3}𝑑\mathrm{\Phi }$$ (56) After integrating by parts, we obtain $$C(\sigma \sigma \sigma )=\pi a^2\left[J_0\left(\frac{\gamma }{2\pi a^2}\|\delta 𝐫\|\sigma \right)1\right]+\frac{1}{2}\gamma \|\delta 𝐫\|H\left(\frac{\gamma }{2\pi a^2}\sigma \|\delta 𝐫\|\right)\sigma $$ (57) where $`H(x)`$ denotes the function $$H(x)=_0^x\frac{J_1(t)}{t}𝑑t$$ (58) For $`\sigma 0`$, the function $`C(\sigma \sigma \sigma )`$ is quadratic: $$C(\sigma \sigma \sigma )\frac{\gamma ^2}{16\pi }\frac{\|\delta 𝐫\|^2}{a^2}\sigma ^2$$ (59) implying that $`W(\delta 𝐕)`$ is Gaussian for large velocity increments $$W(\delta 𝐕)\frac{4}{n\gamma ^2}\frac{a^2}{\|\delta 𝐫\|^2}e^{\frac{4\pi }{n\gamma ^2}\frac{a^2}{\|\delta 𝐫\|^2}\|\delta 𝐕\|^2}(\|\delta 𝐕\|+\mathrm{})$$ (60) The variance of the distribution (60) is $$(\delta 𝐕)^2=\frac{n\gamma ^2}{4\pi }\frac{\|\delta 𝐫\|^2}{a^2}$$ (61) as can be seen directly from equation (54) introducing a cut-off at $`ra`$ instead of $`d`$. For $`\sigma +\mathrm{}`$, the function $`C(\sigma \sigma \sigma )`$ is linear $$C(\sigma \sigma \sigma )\frac{1}{2}\gamma \|\delta 𝐫\|\sigma $$ (62) and we recover the Cauchy distribution (52) for small velocity increments. Therefore the distribution $`W(\delta 𝐕)`$ makes a smooth transition from Cauchy (concave in a semi-log plot) for small fluctuations to Gaussian (convex in a semi-log plot) for large fluctuations. These asymptotic limits were found by Min et al. , but the general expression for the characteristic function (57) is new. It is likely that the two distributions are connected by an exponential tail as predicted and observed numerically by Min et al. Of course, when $`a`$ is reduced, the transition happens at larger fluctuations (see equation (61)) and we have typically a Cauchy distribution. Inversely, for “extended” vortices (large $`a`$), the distribution of velocity increments tends to a Gaussian. Similar results were obtained in Ref. for the distribution of accelerations (note, however, that the functional form of the function $`C(\sigma \sigma \sigma )`$, equation (57), is not exactly the same). ### D The moment $`\delta 𝐕_𝐕`$ The average value of $`\delta 𝐕`$ for a given velocity $`𝐕`$ is given by $$\delta 𝐕_𝐕=\frac{1}{W(𝐕)}W(𝐕,\delta 𝐕)\delta 𝐕d^2(\delta 𝐕)$$ (63) According to equation (24) it can be rewritten $$W(𝐕)\delta 𝐕_𝐕=\frac{1}{(2\pi )^4}A(\rho \rho \rho ,\sigma \sigma \sigma )e^{i(\rho \rho \rho 𝐕+\sigma \sigma \sigma \delta 𝐕)}\delta 𝐕d^2\rho \rho \rho d^2\sigma \sigma \sigma d^2(\delta 𝐕)$$ (64) or, equivalently, $$W(𝐕)\delta 𝐕_𝐕=i\frac{1}{(2\pi )^4}A(\rho \rho \rho ,\sigma \sigma \sigma )\frac{}{\sigma \sigma \sigma }\left\{e^{i(\rho \rho \rho 𝐕+\sigma \sigma \sigma \delta 𝐕)}\right\}d^2\rho \rho \rho d^2\sigma \sigma \sigma d^2(\delta 𝐕)$$ (65) Integrating by parts, we obtain $$W(𝐕)\delta 𝐕_𝐕=i\frac{1}{(2\pi )^4}\frac{A}{\sigma \sigma \sigma }(\rho \rho \rho ,\sigma \sigma \sigma )e^{i(\rho \rho \rho 𝐕+\sigma \sigma \sigma \delta 𝐕)}d^2\rho \rho \rho d^2\sigma \sigma \sigma d^2(\delta 𝐕)$$ (66) Using the identity (38), we can readily carry out the integration on $`\delta 𝐕`$ and $`\sigma `$$`\sigma `$$`\sigma `$ to finally get: $$W(𝐕)\delta 𝐕_𝐕=\frac{i}{(2\pi )^2}\frac{A}{\sigma \sigma \sigma }(\rho \rho \rho ,\mathrm{𝟎})e^{i\rho \rho \rho 𝐕}d^2\rho \rho \rho $$ (67) To go further in the evaluation of the integral, we need to determine the behaviour of $`A(\rho \rho \rho ,\sigma \sigma \sigma )`$ for $`\|\sigma \sigma \sigma \|0`$. In Appendix A it is found that $$A(\rho \rho \rho ,\sigma \sigma \sigma )=e^{nC(\rho \rho \rho )+i\frac{\gamma n}{2}\|\delta 𝐫\|\{\sigma _y\mathrm{cos}(2\theta )\sigma _x\mathrm{sin}(2\theta )\}+o(\|\sigma \sigma \sigma \|^2)}$$ (68) where $`(\sigma _x,\sigma _y)`$ are the components of $`\sigma `$$`\sigma `$$`\sigma `$ in a system of coordinates where the $`x`$-axis coincides with the direction of $`\delta 𝐫`$ and $`\theta `$ denotes the angle that $`\rho `$$`\rho `$$`\rho `$ forms with $`\delta 𝐫`$. We have therefore $$\frac{A}{\sigma _x}(\rho \rho \rho ,\mathrm{𝟎})=i\frac{\gamma n}{2}\|\delta 𝐫\|\mathrm{sin}(2\theta )e^{nC(\rho \rho \rho )}$$ (69) $$\frac{A}{\sigma _y}(\rho \rho \rho ,\mathrm{𝟎})=i\frac{\gamma n}{2}\|\delta 𝐫\|\mathrm{cos}(2\theta )e^{nC(\rho \rho \rho )}$$ (70) The $`x`$-component of equation (67) now becomes $$W(𝐕)\delta V_x_𝐕=\frac{\gamma n}{8\pi ^2}\|\delta 𝐫\|_0^{2\pi }𝑑\theta _0^+\mathrm{}\rho 𝑑\rho \mathrm{cos}\{\rho V\mathrm{cos}(\chi \theta )\}\mathrm{sin}(2\theta )e^{nC(\rho )}$$ (71) where $`\chi `$ denotes the angle that $`𝐕`$ forms with $`\delta 𝐫`$. Using the expansion formula (A7) and the identities (A9)(A10) we can carry out the angular integration to finally obtain $$W(𝐕)\delta V_x_𝐕=\frac{\gamma n}{4\pi }\|\delta 𝐫\|\mathrm{sin}(2\chi )_0^+\mathrm{}J_2(\rho V)e^{nC(\rho )}\rho 𝑑\rho $$ (72) For $`VV_{crit}(N)`$, the contribution of small $`\rho `$’s in the integral (72) is negligible and we can replace $`C(\rho )`$ by its approximate value (8). Therefore $$W(𝐕)\delta V_x_𝐕=\frac{\gamma n}{4\pi }\|\delta 𝐫\|\mathrm{sin}(2\chi )_0^+\mathrm{}J_2(\rho V)e^{\frac{n\gamma ^2}{16\pi }\mathrm{ln}N\rho ^2}\rho 𝑑\rho $$ (73) Using (10) and the identity $$_0^+\mathrm{}J_2(x)e^{\alpha x^2}x𝑑x=22\left(1+\frac{1}{4\alpha }\right)e^{\frac{1}{4\alpha }}$$ (74) we obtain $$\delta V_x_𝐕=\frac{n\gamma }{2}B\left(\frac{4\pi V^2}{n\gamma ^2\mathrm{ln}N}\right)\|\delta 𝐫\|\mathrm{sin}(2\chi )(VV_{crit}(N))$$ (75) where $`B(x)`$ denotes the function $$B(x)=\frac{1}{x}(e^x1x)$$ (76) Similarly, we find $$\delta V_y_𝐕=\frac{n\gamma }{2}B\left(\frac{4\pi V^2}{n\gamma ^2\mathrm{ln}N}\right)\|\delta 𝐫\|\mathrm{cos}(2\chi )(VV_{crit}(N))$$ (77) For $`VV_{crit}(N)`$, the integral (72) is dominated by small values of $`\rho `$ and one must use the general expression (7) for $`C(\rho )`$. With the change of variables $`z=\rho V`$, equation (72) becomes $$W(𝐕)\delta V_x_𝐕=\frac{\gamma n}{4\pi V^2}\|\delta 𝐫\|\mathrm{sin}(2\chi )_0^+\mathrm{}J_2(z)e^{nC(\frac{z}{V})}z𝑑z$$ (78) Using the recursion formula $$J_2(z)=\frac{2}{z}J_1(z)J_0(z)$$ (79) we obtain $$W(𝐕)\delta V_x_𝐕=\frac{\gamma n}{4\pi V^2}\|\delta 𝐫\|\mathrm{sin}(2\chi )\left\{_0^+\mathrm{}2J_1(z)e^{nC(\frac{z}{V})}𝑑z_0^+\mathrm{}J_0(z)e^{nC(\frac{z}{V})}z𝑑z\right\}$$ (80) When $`V\mathrm{}`$, the first integral is convergent while the second, equal to $`2\pi V^2W(𝐕)`$, decreases like $`V^2`$ \[see equation (11)\]. Therefore, to leading order in $`1/V`$: $$W(𝐕)\delta V_x_𝐕\frac{\gamma n}{2\pi V^2}\|\delta 𝐫\|\mathrm{sin}(2\chi )_0^+\mathrm{}J_1(z)𝑑z$$ (81) With (11) and the identity $$_0^+\mathrm{}J_n(x)𝑑x=1$$ (82) we find $$\delta V_x_𝐕=\frac{2\pi V^2}{\gamma }\|\delta 𝐫\|\mathrm{sin}(2\chi )(VV_{crit}(N))$$ (83) Similarly $$\delta V_y_𝐕=\frac{2\pi V^2}{\gamma }\|\delta 𝐫\|\mathrm{cos}(2\chi )(VV_{crit}(N))$$ (84) Equations (75) (77) and (83) (84) can be written more compactly in the form $$\delta 𝐕_𝐕=\frac{n\gamma }{2}B\left(\frac{4\pi V^2}{n\gamma ^2\mathrm{ln}N}\right)\left\{\delta 𝐫_{}+2\frac{(𝐕_{}\delta 𝐫)}{V^2}𝐕\right\}(VV_{crit}(N))$$ (85) $$\delta 𝐕_𝐕=\frac{2\pi }{\gamma }V^2\left\{\delta 𝐫_{}+2\frac{(𝐕_{}\delta 𝐫)}{V^2}𝐕\right\}(VV_{crit}(N))$$ (86) where $`B(x)`$ is defined by equation (76). Equations (85) (86) can be compared with equation (172) of Chandrasekhar & von Neumann for the increment of the gravitational field between two neighboring points. Equation (86) has a clear physical interpretation in the nearest neighbor approximation. In Ref. it was shown that the high velocity tail of the distribution $`W(𝐕)`$ is produced solely by the nearest neighbor. Now, the velocity and the velocity difference are related to the position $`𝐫`$ of the nearest neighbor by \[see equations (2) and (23)\]: $$𝐕=\frac{\gamma }{2\pi }\frac{𝐫_{}}{r^2}$$ (87) and $$\delta 𝐕=\frac{\gamma }{2\pi }\left\{\frac{\delta 𝐫_{}}{r^2}\frac{2(𝐫\delta 𝐫)𝐫_{}}{r^4}\right\}$$ (88) Eliminating $`𝐫`$ between these two expressions, we obtain $$(\delta 𝐕)_𝐕=\frac{2\pi }{\gamma }\left\{\delta 𝐫_{}V^2+2(𝐕_{}\delta 𝐫)𝐕\right\}$$ (89) which coincides with the expression (86) of $`\delta 𝐕_𝐕`$ for large $`V`$’s. Note that $$(\delta 𝐕)_𝐕^2=\frac{4\pi ^2}{\gamma ^2}\|\delta 𝐫\|^2V^4$$ (90) From (89) and (90) we can determine all conditional moments of $`\delta 𝐕`$ in the nearest neighbor approximation (i.e., valid for large $`V`$’s). If we account for a spectrum of circulations among the vortices, we have in place of equations (85)(86): $$\delta 𝐕_𝐕=\frac{n\overline{\gamma }}{2}B\left(\frac{4\pi V^2}{n\overline{\gamma ^2}\mathrm{ln}N}\right)\left\{\delta 𝐫_{}+2\frac{(𝐕_{}\delta 𝐫)}{V^2}𝐕\right\}(VV_{crit}(N))$$ (91) $$\delta 𝐕_𝐕=\frac{2\pi \overline{\gamma }}{\overline{\gamma ^2}}V^2\left\{\delta 𝐫_{}+2\frac{(𝐕_{}\delta 𝐫)}{V^2}𝐕\right\}(VV_{crit}(N))$$ (92) In particular, for a neutral system $`\delta 𝐕_𝐕=\mathrm{𝟎}`$. This is to be expected (at least for a symmetrical distribution of circulations) since a vortex with circulation $`\gamma `$ located in $`𝐫`$ produces the same velocity but an opposite velocity increment as a vortex with circulation $`+\gamma `$ located in $`𝐫`$. Physically, we must realize that the calculations developed in this section are not valid for large values of $`\|𝐕\|`$ (for a fixed $`\|\delta 𝐫\|`$). Indeed, large velocities are produced by a vortex very close to the origin. In that case, $`\|𝐕\|=\frac{\gamma }{2\pi r}\mathrm{}`$ while the velocity produced in $`\delta 𝐫`$ is bound to the value $`\|𝐕_1\|\frac{\gamma }{2\pi \|\delta 𝐫\|}`$. Therefore, the velocity difference $`\delta 𝐕=𝐕𝐕_1`$ should behave like $`\|𝐕\|`$. Now, according to formula (86), we have $`\|\delta 𝐕\|V^2`$. This inconsistency is related to the fact that, when $`\|𝐕\|+\mathrm{}`$, the velocity difference $`\delta 𝐕`$ cannot be represented to any degree of accuracy by the Taylor expansion (23) which assumes $`\|\delta 𝐫\|\|𝐫\|`$. We expect therefore that the results of this section will be valid only for velocities $`\|𝐕\|<\frac{\gamma }{2\pi \|\delta 𝐫\|}`$. In view of these remarks, it is necessary to consider now the more general problem of the velocity correlations between two points separated by a finite distance. ## IV The correlations in the velocities occurring at two points separated by a finite distance ### A A general formula for $`W(𝐕_0,𝐕_1)`$ The general expression for $`W_N(𝐕_0,𝐕_1)`$, the bivariate probability that a velocity $$𝐕_0=\frac{\gamma }{2\pi }\underset{i=1}{\overset{N}{}}\frac{𝐫_i}{r_i^2}$$ (93) occurs at the center of the domain and that, simultaneously, a velocity $$𝐕_1=\frac{\gamma }{2\pi }\underset{i=1}{\overset{N}{}}\frac{(𝐫_i𝐫_1)_{}}{\|𝐫_i𝐫_1\|^2}$$ (94) occurs at a point $`M_1`$ separated from the first by a distance $`𝐫_1`$ can be readily written down following Markov’s method outlined in Ref. , section II.A. We have $$W_N(𝐕_0,𝐕_1)=\frac{1}{(2\pi )^4}A_N(\rho \rho \rho ,\sigma \sigma \sigma )e^{i(\rho \rho \rho 𝐕_0+\sigma \sigma \sigma 𝐕_1)}d^2\rho \rho \rho d^2\sigma \sigma \sigma $$ (95) with: $$A_N(\rho \rho \rho ,\sigma \sigma \sigma )=\left(1\frac{1}{\pi R^2}_{\|𝐫\|=0}^R\left(1e^{i\frac{\gamma }{2\pi }(\frac{\rho \rho \rho 𝐫_{}}{r^2}+\frac{\sigma \sigma \sigma (𝐫𝐫_1)_{}}{\|𝐫𝐫_1\|^2})}\right)d^2𝐫\right)^N$$ (96) We now consider the limit when the number of vortices and the size of the domain go to infinity in such a way that the density remains finite. If the integral occurring in equation (96) increases less rapidly than $`N`$, then $$A(\rho \rho \rho ,\sigma \sigma \sigma )=e^{nC(\rho \rho \rho ,\sigma \sigma \sigma )}$$ (97) with $$C(\rho \rho \rho ,\sigma \sigma \sigma )=_{\|𝐫\|=0}^R(1e^{i\frac{\gamma }{2\pi }(\frac{\rho \rho \rho 𝐫_{}}{r^2}+\frac{\sigma \sigma \sigma (𝐫𝐫_1)_{}}{\|𝐫𝐫_1\|^2})})d^2𝐫$$ (98) Writing $$C(\rho \rho \rho ,\sigma \sigma \sigma )=C(\rho \rho \rho )+D(\rho \rho \rho ,\sigma \sigma \sigma )$$ (99) where $`C(\rho \rho \rho )=C(\rho \rho \rho ,\mathrm{𝟎})`$ is given by equation (7), we have $$D(\rho \rho \rho ,\sigma \sigma \sigma )=_{\|𝐫\|=0}^Re^{i\frac{\gamma }{2\pi }\frac{\rho \rho \rho 𝐫_{}}{r^2}}\left(1e^{i\frac{\gamma }{2\pi }\frac{\sigma \sigma \sigma (𝐫𝐫_1)_{}}{\|𝐫𝐫_1\|^2})}\right)d^2𝐫$$ (100) Equations (95) (97) (99) and (100) formally solve the problem but it does not seem possible to obtain an explicit expression for $`D(\rho \rho \rho ,\sigma \sigma \sigma )`$. However, like in section III D, if we are interested only in the moments of $`𝐕_1`$ for a given value of $`𝐕_0`$, we need only the behaviour of $`D(\rho \rho \rho ,\sigma \sigma \sigma )`$ for $`\|\sigma \sigma \sigma \|0`$. Expanding the exponential term which appears under the integral sign in a power series in $`\|\sigma \sigma \sigma \|`$, we obtain to first order $$D(\rho \rho \rho ,\sigma \sigma \sigma )=D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )+o(\|\sigma \sigma \sigma \|^2)$$ (101) with $$D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )=\frac{i\gamma }{2\pi }_{\|𝐫\|=0}^Re^{i\frac{\gamma }{2\pi }\frac{\rho \rho \rho 𝐫_{}}{r^2}}\frac{\sigma \sigma \sigma (𝐫𝐫_1)_{}}{\|𝐫𝐫_1\|^2}d^2𝐫$$ (102) According to equations (97) (99) and (101), $`A(\rho \rho \rho ,\sigma \sigma \sigma )`$ can be expressed as: $$A(\rho \rho \rho ,\sigma \sigma \sigma )=e^{nC(\rho \rho \rho )nD^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )+o(\|\sigma \sigma \sigma \|^2)}$$ (103) In Appendix B, it is found that $`D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )={\displaystyle \frac{\gamma ^2}{2\pi }}\rho (\sigma _x\mathrm{sin}\theta _1+\sigma _y\mathrm{cos}\theta _1){\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dr}{r}}J_2\left({\displaystyle \frac{\gamma \rho }{2\pi r}}\right)\mathrm{ln}r_>`$ (104) $`(\sigma _x\mathrm{sin}\theta _1+\sigma _y\mathrm{cos}\theta _1){\displaystyle \underset{n=1}{\overset{+\mathrm{}}{}}}A_n(\rho ,r_1)\mathrm{cos}(n\theta _1)`$ (105) $`(\sigma _x\mathrm{cos}\theta _1\sigma _y\mathrm{sin}\theta _1){\displaystyle \underset{n=1}{\overset{+\mathrm{}}{}}}B_n(\rho ,r_1)\mathrm{sin}(n\theta _1)`$ (106) $`{\displaystyle \frac{\gamma ^2}{4\pi }}\rho \mathrm{ln}R(\sigma _x\mathrm{sin}\theta _1+\sigma _y\mathrm{cos}\theta _1)i{\displaystyle \frac{\gamma }{2}}r_1\sigma _y`$ (107) where $`r_>`$ (resp. $`r_<`$) is the larger (resp. smaller) of $`(r,r_1)`$, $`A_n(\rho ,r_1)`$ and $`B_n(\rho ,r_1)`$ are defined by equations (B15) and (B16), $`(\sigma _x,\sigma _y)`$ are the components of $`\sigma `$$`\sigma `$$`\sigma `$ in a system of coordinates where the $`x`$-axis is in the direction of $`𝐫_1`$ and $`\theta _1`$ is the angle that $`\rho \rho \rho _{}`$ forms with $`𝐫_1`$. ### B The first moment $`𝐕_1_{𝐕_0}`$ The average value of $`𝐕_1`$ for a given $`𝐕_0`$ is given by: $$𝐕_1_{𝐕_0}=\frac{1}{W(𝐕_0)}W(𝐕_0,𝐕_1)𝐕_1d^2𝐕_1$$ (108) By a procedure similar to that adopted in section III D, we find that an equivalent expression for $`𝐕_1_{𝐕_0}`$ is $$W(𝐕_0)𝐕_1_{𝐕_0}=\frac{i}{(2\pi )^2}\frac{A}{\sigma \sigma \sigma }(\rho \rho \rho ,\mathrm{𝟎})e^{i\rho \rho \rho 𝐕_0}d^2\rho \rho \rho $$ (109) or, using equation (103) $$W(𝐕_0)𝐕_1_{𝐕_0}=\frac{in}{(2\pi )^2}\frac{D^{(1)}}{\sigma \sigma \sigma }(\rho \rho \rho )e^{nC(\rho \rho \rho )}e^{i\rho \rho \rho 𝐕_0}d^2\rho \rho \rho $$ (110) Introducing polar coordinates, it can be written: $$W(𝐕_0)𝐕_1_{𝐕_0}=\frac{in}{(2\pi )^2}_0^+\mathrm{}𝐙(\rho ,𝐕_0)e^{nC(\rho )}\rho 𝑑\rho $$ (111) where $$𝐙(\rho ,𝐕_0)=_0^{2\pi }\frac{D^{(1)}}{\sigma \sigma \sigma }(\rho \rho \rho )e^{i\rho V_0\mathrm{cos}\varphi }𝑑\varphi $$ (112) In equation (112), $`\varphi `$ denotes the angle that $`\rho `$$`\rho `$$`\rho `$ forms with $`𝐕_0`$. It is related to $`\chi `$ and $`\theta _1`$, the angles that $`𝐕_0`$ and $`\rho \rho \rho _{}`$ form with $`𝐫_1`$ by: $$\varphi =\chi \theta _1\frac{\pi }{2}$$ (113) Therefore, an alternative expression for $`𝐙(\rho ,𝐕_0)`$ is: $$𝐙(\rho ,𝐕_0)=_0^{2\pi }\frac{D^{(1)}}{\sigma \sigma \sigma }(\rho \rho \rho )e^{i\rho V_0\mathrm{sin}(\chi +\theta _1)}𝑑\theta _1$$ (114) Using (107) and the identity $$e^{ix\mathrm{sin}\theta }=J_0(x)+2\underset{n=1}{\overset{+\mathrm{}}{}}J_{2n}(x)\mathrm{cos}(2n\theta )+2i\underset{n=0}{\overset{+\mathrm{}}{}}J_{2n+1}(x)\mathrm{sin}((2n+1)\theta )$$ (115) we find after lengthy calculations that $`Z_x(\rho ,𝐕_0)=i\gamma ^2\rho J_1(\rho V_0)\mathrm{cos}\chi {\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dr}{r}}J_2\left({\displaystyle \frac{\gamma \rho }{2\pi r}}\right)\mathrm{ln}r_>i{\displaystyle \frac{\gamma ^2}{2}}\rho J_1(\rho V_0)\mathrm{cos}\chi \mathrm{ln}R`$ (116) $`+\pi {\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}[C_{2n1}(\rho ,r_1)D_{2n+1}(\rho ,r_1)]J_{2n}(\rho V_0)\mathrm{sin}(2n\chi )`$ (117) $`i\pi {\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}[C_{2n}(\rho ,r_1)D_{2n+2}(\rho ,r_1)]J_{2n+1}(\rho V_0)\mathrm{cos}((2n+1)\chi )`$ (118) and $`Z_y(\rho ,𝐕_0)=i\gamma ^2\rho J_1(\rho V_0)\mathrm{sin}\chi {\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dr}{r}}J_2\left({\displaystyle \frac{\gamma \rho }{2\pi r}}\right)\mathrm{ln}r_>`$ (119) $`i{\displaystyle \frac{\gamma ^2}{2}}\rho J_1(\rho V_0)\mathrm{sin}\chi \mathrm{ln}Ri\pi \gamma r_1J_0(\rho V_0)`$ (120) $`\pi {\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}[C_{2n1}(\rho ,r_1)+D_{2n+1}(\rho ,r_1)]J_{2n}(\rho V_0)\mathrm{cos}(2n\chi )`$ (121) $`i\pi {\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}[C_{2n}(\rho ,r_1)+D_{2n+2}(\rho ,r_1)]J_{2n+1}(\rho V_0)\mathrm{sin}((2n+1)\chi )`$ (122) where we have introduced the notations: $`C_n(\rho ,r_1)=A_n(\rho ,r_1)+B_n(\rho ,r_1)`$ (123) $`D_n(\rho ,r_1)=A_n(\rho ,r_1)B_n(\rho ,r_1)`$ (124) for $`n>0`$ and $`C_n,D_n=0`$ otherwise. Evaluating $`C_n`$ and $`D_n`$ as defined in equations (123) (124), we find explicitly that: $`C_n(\rho ,r_1)={\displaystyle \frac{\gamma ^2}{2\pi }}\rho {\displaystyle \frac{i^n}{n}}{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dr}{r}}\left({\displaystyle \frac{r_<}{r_>}}\right)^nJ_{n+2}\left({\displaystyle \frac{\gamma \rho }{2\pi r}}\right)`$ (125) $`D_n(\rho ,r_1)={\displaystyle \frac{\gamma ^2}{2\pi }}\rho {\displaystyle \frac{i^n}{n}}{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dr}{r}}\left({\displaystyle \frac{r_<}{r_>}}\right)^nJ_{n2}\left({\displaystyle \frac{\gamma \rho }{2\pi r}}\right)`$ (126) and we recall that $`J_n(z)=(1)^nJ_n(z)`$. ## V The average value of $`𝐕_1`$ and the function $`𝐕_1_{V_0}`$ ### A The average velocity $`𝐕_1`$ The moment $`𝐕_1`$ is simply obtained by averaging $`𝐕_1_{𝐕_0}`$ over all possible values of $`𝐕_0`$: $$𝐕_1=W(𝐕_0)𝐕_1_{𝐕_0}d^2𝐕_0$$ (127) Substituting equation (111) in equation (127) and introducing polar coordinates, this can be rewritten: $`𝐕_1={\displaystyle \frac{in}{(2\pi )^2}}{\displaystyle _0^{2\pi }}𝑑\chi {\displaystyle _0^+\mathrm{}}V_0𝑑V_0{\displaystyle _0^+\mathrm{}}\rho 𝑑\rho e^{nC(\rho )}𝐙(\rho ,V_0,\chi )`$ (128) where $`\chi `$ denotes the angle that $`𝐕_0`$ forms with $`𝐫_1`$. When the integration is performed over $`\chi `$, using equations (118)(122), it is seen that the components of $`𝐕_1`$ reduce to: $$V_{1x}=0$$ (129) and: $$V_{1y}=i\frac{n}{2}_0^+\mathrm{}V_0𝑑V_0_0^+\mathrm{}\rho 𝑑\rho e^{nC(\rho )}(i\gamma r_1+D_1(\rho ,r_1))J_0(\rho V_0)$$ (130) where, according to (126): $$D_1(\rho ,r_1)=i\frac{\gamma ^2}{2\pi }\rho _0^+\mathrm{}J_1\left(\frac{\gamma \rho }{2\pi r}\right)\frac{r_<}{r_>}\frac{dr}{r}$$ (131) Under this form, it is not possible to interchange the order of integration in (130). However, an alternative expression for $`V_{1y}`$ can be obtained along the following lines. Writing $$V_{1y}=_0^+\mathrm{}𝑑V_0_0^+\mathrm{}\rho V_0J_0(\rho V_0)\psi (\rho )𝑑\rho $$ (132) where $$\psi (\rho )=i\frac{n}{2}e^{nC(\rho )}(i\gamma r_1+D_1(\rho ,r_1))$$ (133) and integrating by parts with the identity $$xJ_0(x)=\frac{d}{dx}(xJ_1(x))$$ (134) we obtain: $$V_{1y}=_0^+\mathrm{}𝑑V_0_0^+\mathrm{}\rho J_1(\rho V_0)\psi ^{}(\rho )𝑑\rho $$ (135) It is now possible to interchange the order of integration. Using formula (82), we find $$V_{1y}=_0^+\mathrm{}\psi ^{}(\rho )𝑑\rho =\psi (0)\psi (+\mathrm{})=\frac{n}{2}\gamma r_1$$ (136) Hence $$𝐕_1=\frac{1}{2}n\gamma 𝐫_1$$ (137) This result is of course to be expected. It corresponds to the solid rotation produced by a uniform distribution of point vortices with density $`n\gamma `$. ### B The function $`𝐕_1_{V_0}`$ A simple result is obtained when $`𝐕_1_{𝐕_0}`$ is averaged over all mutual orientations of $`𝐫_1`$ and $`𝐕_0`$. Define $$𝐕_1_{V_0}=\frac{1}{2\pi }_0^{2\pi }𝐕_1_{𝐕_0}𝑑\chi $$ (138) The integral of $`𝐕_1_{V_0}`$ over $`V_0`$ is just the average velocity in $`𝐫_1`$: $$𝐕_1=_0^+\mathrm{}W(𝐕_0)𝐕_1_{V_0}2\pi V_0𝑑V_0$$ (139) Therefore, comparing equation (139) with equation (128) we directly obtain $$V_{1x}_{V_0}=0$$ (140) $$W(𝐕_0)V_{1y}_{V_0}=i\frac{n}{4\pi }_0^+\mathrm{}\rho 𝑑\rho e^{nC(\rho )}(i\gamma r_1+D_1(\rho ,r_1))J_0(\rho V_0)$$ (141) Using (131), we can rewrite our expression for $`V_{1y}_{V_0}`$ in the form $$W(𝐕_0)V_{1y}_{V_0}=\frac{n\gamma ^2}{8\pi ^2}_0^+\mathrm{}\rho ^2𝑑\rho e^{nC(\rho )}J_0(\rho V_0)R(\rho ,r_1)$$ (142) where $`R(\rho ,r_1)`$ denotes the function $$R(\rho ,r_1)=\frac{2\pi r_1}{\gamma \rho }+_0^+\mathrm{}J_1\left(\frac{\gamma \rho }{2\pi r}\right)\frac{r_<}{r_>}\frac{dr}{r}$$ (143) Remembering that the range of integration in equation (143) has to be broken at $`r_1`$ with the prescription given at the end of section IV A, the equation defining $`R(\rho ,r_1)`$ has explicitly the form $$R(\rho ,r_1)=\frac{2\pi r_1}{\gamma \rho }+\frac{1}{r_1}_0^{r_1}J_1\left(\frac{\gamma \rho }{2\pi r}\right)𝑑r+r_1_{r_1}^+\mathrm{}J_1\left(\frac{\gamma \rho }{2\pi r}\right)\frac{dr}{r^2}$$ (144) With the change of variables $`z=\frac{\gamma \rho }{2\pi r}`$, it becomes $$R(\rho ,r_1)=\frac{2\pi r_1}{\gamma \rho }+\frac{2\pi r_1}{\gamma \rho }_0^{\frac{\gamma \rho }{2\pi r_1}}J_1(z)𝑑z+\frac{\gamma \rho }{2\pi r_1}_{\frac{\gamma \rho }{2\pi r_1}}^+\mathrm{}\frac{J_1(z)}{z^2}𝑑z$$ (145) Using the identity $$\frac{d}{dz}\left(\frac{J_n(z)}{z^n}\right)=\frac{J_{n+1}(z)}{z^n}$$ (146) for $`n=0`$ and integrating by parts, we obtain: $$R(\rho ,r_1)=\frac{2\pi r_1}{\gamma \rho }J_0\left(\frac{\gamma \rho }{2\pi r_1}\right)+\frac{\gamma \rho }{2\pi r_1}_{\frac{\gamma \rho }{2\pi r_1}}^+\mathrm{}\frac{J_1(z)}{z^2}𝑑z$$ (147) From equation (147) it is readily found that $$R(\rho ,r_1)0(r_10)$$ (148) $$R(\rho ,r_1)=\frac{2\pi r_1}{\gamma \rho }+\frac{\gamma \rho }{4\pi r_1}\mathrm{ln}\left(\frac{2\pi r_1}{\gamma \rho }\right)(r_1+\mathrm{})$$ (149) where we recall that $`J_1(z)\frac{z}{2}`$ for $`z0`$. ### C The asymptotic behaviour of $`𝐕_1_{V_0}`$ for $`r_1+\mathrm{}`$ Substituting for the asymptotic expansion of $`R(\rho ,r_1)`$ from equation (149) in equation (142), we obtain for $`r_1+\mathrm{}`$: $`W(𝐕_0)V_{1y}_{V_0}={\displaystyle \frac{n\gamma r_1}{4\pi }}{\displaystyle _0^+\mathrm{}}e^{nC(\rho )}J_0(\rho V_0)\rho 𝑑\rho `$ (150) $`+{\displaystyle \frac{n\gamma ^3}{32\pi ^3r_1}}{\displaystyle _0^+\mathrm{}}e^{nC(\rho )}J_0(\rho V_0)\mathrm{ln}\left({\displaystyle \frac{\gamma \rho }{2\pi r_1}}\right)\rho ^3𝑑\rho `$ (151) The first integral is equal to $`2\pi W(𝐕_0)`$ so the first term is just the average velocity $`V_{1y}`$ \[see section V A\]. Defining $$\mathrm{\Delta }𝐕_1_{V_0}=𝐕_1_{V_0}𝐕_1$$ (152) we have $$W(𝐕_0)\mathrm{\Delta }𝐕_1_{V_0}=\frac{n\gamma ^3}{32\pi ^3}\frac{𝐫_1}{r_1^2}_0^+\mathrm{}e^{nC(\rho )}J_0(\rho V_0)\mathrm{ln}\left(\frac{\gamma \rho }{2\pi r_1}\right)\rho ^3𝑑\rho $$ (153) For $`V_0V_{crit}(N)`$, the integral is dominated by large values of $`\rho `$ and we can make the approximation $$W(𝐕_0)\mathrm{\Delta }𝐕_1_{V_0}\frac{n\gamma ^3}{64\pi ^3}\frac{𝐫_1}{r_1^2}\mathrm{ln}N_0^+\mathrm{}e^{\frac{n\gamma ^2}{16\pi }\mathrm{ln}N\rho ^2}J_0(\rho V_0)\rho ^3𝑑\rho $$ (154) Using (10) and the identity $$_0^+\mathrm{}e^{\alpha x^2}J_0(\beta x)x^3𝑑x=\frac{1}{2\alpha ^2}\left(1\frac{\beta ^2}{4\alpha }\right)e^{\frac{\beta ^2}{4\alpha }}$$ (155) we obtain $$\mathrm{\Delta }𝐕_1_{V_0}=\frac{\gamma }{2\pi }\frac{𝐫_1}{r_1^2}\left(1\frac{4\pi V_0^2}{n\gamma ^2\mathrm{ln}N}\right)(V_0V_{crit}(N))$$ (156) For $`V_0V_{crit}(N)`$, we can follow the procedure outlined in Ref. , section II.C., and we find after some calculations $$\mathrm{\Delta }𝐕_1_{V_0}=\frac{\gamma }{2\pi }\frac{𝐫_1}{r_1^2}(V_0V_{crit}(N))$$ (157) According to these equations $$\mathrm{\Delta }𝐕_1_{V_0}=\frac{\gamma }{2\pi }\frac{𝐫_1}{r_1^2}(V_00)$$ (158) $$\mathrm{\Delta }𝐕_1_{V_0}=\frac{\gamma }{2\pi }\frac{𝐫_1}{r_1^2}(V_0+\mathrm{})$$ (159) This is similar to the velocity produced by a fictitious point vortex located in $`O`$ with circulation $`+\gamma `$ (when $`V_00`$) and $`\gamma `$ (when $`V_0+\mathrm{}`$). It is not clear whether these results have a deeper significance than is apparent at first sight. We observe that $`\mathrm{\Delta }𝐕_1_{V_0}`$ changes sign at $$V_c=\left(\frac{n\gamma ^2}{4\pi }\mathrm{ln}N\right)^{1/2}$$ (160) We can also check that $$_0^+\mathrm{}W(𝐕_0)\mathrm{\Delta }𝐕_1_{V_0}2\pi V_0𝑑V_0=0$$ (161) in agreement with (137). ## VI The first moment of $`𝐕_1`$ in the direction of $`𝐕_0`$ and its average ### A The moment $`V_{1\|\|}_{V_0}`$ A quantity of physical interest is the average value of $`𝐕_1`$ in the direction of $`𝐕_0`$. Let us define $$V_{1\|\|}=𝐕_1111_0=V_{1x}\mathrm{cos}\chi +V_{1y}\mathrm{sin}\chi $$ (162) where $`111_0`$ is the unit vector in the direction of $`𝐕_0`$. Then $$V_{1\|\|}_{𝐕_0}=V_{1x}_{𝐕_0}\mathrm{cos}\chi +V_{1y}_{𝐕_0}\mathrm{sin}\chi $$ (163) A relatively simple formula can be obtained when the foregoing expression for $`V_{1\|\|}_{𝐕_0}`$ is averaged over all possible orientations of $`𝐕_0`$. Define $$V_{1\|\|}_{V_0}=\frac{1}{2\pi }_0^{2\pi }V_{1\|\|}_{𝐕_0}𝑑\chi $$ (164) When equation (111) is substituted in equation (164) and the integration is performed over $`\chi `$, only the logarithmic terms survive in the series (118) (122). We are thus left with $$W(𝐕_0)V_{1\|\|}_{V_0}=\frac{n\gamma ^2}{4\pi ^2}_0^+\mathrm{}\rho ^2𝑑\rho e^{nC(\rho )}J_1(\rho V_0)Q(\rho ,r_1)$$ (165) where we have written $$Q(\rho ,r_1)=_0^+\mathrm{}J_2\left(\frac{\gamma \rho }{2\pi r}\right)\mathrm{ln}r_>\frac{dr}{r}\frac{1}{2}\mathrm{ln}R$$ (166) With the usual interpretation of the notation $`r_>`$, see section IV A, the equation defining $`Q(\rho ,r_1)`$ has explicitly the form $$Q(\rho ,r_1)=\mathrm{ln}r_1_0^{r_1}J_2\left(\frac{\gamma \rho }{2\pi r}\right)\frac{dr}{r}+_{r_1}^+\mathrm{}J_2\left(\frac{\gamma \rho }{2\pi r}\right)\mathrm{ln}r\frac{dr}{r}\frac{1}{2}\mathrm{ln}R$$ (167) With the change of variables $`z=\frac{\gamma \rho }{2\pi r}`$, it becomes: $$Q(\rho ,r_1)=\mathrm{ln}r_1_{\frac{\gamma \rho }{2\pi r_1}}^+\mathrm{}J_2(z)\frac{dz}{z}+_0^{\frac{\gamma \rho }{2\pi r_1}}J_2(z)\mathrm{ln}\left(\frac{\gamma \rho }{2\pi z}\right)\frac{dz}{z}\frac{1}{2}\mathrm{ln}R$$ (168) Using the identity (146) with $`n=1`$ and integrating by parts, we obtain: $$Q(\rho ,r_1)=\frac{1}{2}\mathrm{ln}\left(\frac{\gamma \rho }{2\pi R}\right)\frac{1}{2}\mathrm{ln}ϵ_ϵ^{\frac{\gamma \rho }{2\pi r_1}}\frac{J_1(z)}{z^2}𝑑z$$ (169) where it is understood that $`ϵ0`$. From equation (169), it is readily found that $$Q(\rho ,r_1)=\frac{1}{2}\mathrm{ln}\left(\frac{\gamma \rho }{2\pi R}\right)(r_10)$$ (170) $$Q(\rho ,r_1)=\frac{1}{2}\mathrm{ln}\left(\frac{R}{r_1}\right)(r_1+\mathrm{})$$ (171) ### B The asymptotic behaviour of $`V_{1\|\|}_{V_0}`$ for $`r_10`$ and $`r_1+\mathrm{}`$ First considering the limit of $`V_{1\|\|}_{V_0}`$ for $`r_10`$, we have according to (165) and (170) $$W(𝐕_0)V_{1\|\|}_{V_0}=\frac{n\gamma ^2}{8\pi ^2}_0^+\mathrm{}e^{nC(\rho )}J_1(\rho V_0)\mathrm{ln}\left(\frac{\gamma \rho }{2\pi R}\right)\rho ^2𝑑\rho $$ (172) With the expression (7) for $`C(\rho )`$, the foregoing expression can be rewritten: $$W(𝐕_0)V_{1\|\|}_{V_0}=\frac{n}{2\pi }_0^+\mathrm{}C^{}(\rho )e^{nC(\rho )}J_1(\rho V_0)\rho 𝑑\rho $$ (173) Using (134) and integrating by parts, we find $$W(𝐕_0)V_{1\|\|}_{V_0}=\frac{V_0}{2\pi }_0^+\mathrm{}e^{nC(\rho )}J_0(\rho V_0)\rho 𝑑\rho $$ (174) The integral is just $`2\pi W(𝐕_0)`$. Hence $$V_{1\|\|}_{V_0}=V_0(r_10)$$ (175) a result of course to be expected. Considering now the behaviour of $`V_{1\|\|}_{V_0}`$ for $`r_1+\mathrm{}`$, we have according to (165) and (171): $$W(𝐕_0)V_{1\|\|}_{V_0}=\frac{n\gamma ^2}{8\pi ^2}\mathrm{ln}\left(\frac{R}{r_1}\right)_0^+\mathrm{}e^{nC(\rho )}J_1(\rho V_0)\rho ^2𝑑\rho $$ (176) For $`V_0V_{crit}(N)`$, we can make the approximation $$W(𝐕_0)V_{1\|\|}_{V_0}\frac{n\gamma ^2}{8\pi ^2}\mathrm{ln}\left(\frac{R}{r_1}\right)_0^+\mathrm{}e^{\frac{n\gamma ^2}{16\pi }\mathrm{ln}N\rho ^2}J_1(\rho V_0)\rho ^2𝑑\rho $$ (177) Integrating by parts, we obtain $$V_{1\|\|}_{V_0}=\frac{2}{\mathrm{ln}N}V_0\mathrm{ln}\left(\frac{R}{r_1}\right)(V_0V_{crit}(N))$$ (178) For $`V_0V_{crit}(N)`$, we can adapt the procedure outlined in Ref. section II.C and we find after some calculations $$V_{1\|\|}_{V_0}=\frac{n\gamma ^2}{\pi V_0}\mathrm{ln}\left(\frac{R}{r_1}\right)(V_0V_{crit}(N))$$ (179) ### C The moment $`V_{1\|\|}`$ If we now average $`V_{1\|\|}_{V_0}`$ over $`V_0`$, we shall obtain a function of $`\|𝐫_1\|`$ only: $$V_{1\|\|}=_0^+\mathrm{}W(𝐕_0)V_{1\|\|}_{V_0}2\pi V_0𝑑V_0$$ (180) which will describe the correlations in the velocities occurring simultaneously at two points separated by a distance $`\|𝐫_1\|`$. Equation (180) is clearly the same as $$V_{1\|\|}=\frac{𝐕_0𝐕_1}{V_0}=W(𝐕_0,𝐕_1)𝐕_1111_0d^2𝐕_0d^2𝐕_1$$ (181) Substituting for $`V_{1\|\|}_{V_0}`$ from equation (165) in equation (180), we have $$V_{1\|\|}=\frac{n\gamma ^2}{2\pi }_0^+\mathrm{}V_0𝑑V_0_0^+\mathrm{}\rho ^2𝑑\rho e^{nC(\rho )}J_1(\rho V_0)Q(\rho ,r_1)$$ (182) It is apparent that we cannot interchange the order of integration in equation (182). Following a method similar to that adopted in section V A, we rewrite equation (182) in the form: $$V_{1\|\|}=_0^+\mathrm{}𝑑V_0_0^+\mathrm{}V_0J_1(\rho V_0)\mathrm{\Phi }(\rho )𝑑\rho $$ (183) where $$\mathrm{\Phi }(\rho )=\frac{n\gamma ^2}{2\pi }\rho ^2e^{nC(\rho )}Q(\rho ,r_1)$$ (184) Integrating by parts the second integral with the identity (48) we obtain $$V_{1\|\|}=_0^+\mathrm{}𝑑V_0_0^+\mathrm{}J_0(\rho V_0)\mathrm{\Phi }^{}(\rho )𝑑\rho $$ (185) Under this form, it is now possible to interchange the order of integration. Using formula (82), we find that $$V_{1\|\|}=_0^+\mathrm{}\frac{\mathrm{\Phi }^{}(\rho )}{\rho }𝑑\rho $$ (186) Again integrating by parts we obtain $$V_{1\|\|}=_0^+\mathrm{}\frac{\mathrm{\Phi }(\rho )}{\rho ^2}𝑑\rho =\frac{n\gamma ^2}{2\pi }_0^+\mathrm{}e^{nC(\rho )}Q(\rho ,r_1)𝑑\rho $$ (187) In Appendix C 1, we derive an alternative form of equation (187) in terms of $`\mathrm{\Gamma }`$-functions. Indeed, we show that $$V_{1\|\|}=\|𝐕_0\|\left(\frac{n\gamma ^2}{\pi \mathrm{ln}N}\right)^{1/2}_0^+\mathrm{}𝑑z\frac{J_1(z)}{z^2}_0^{s^2z^2}e^tt^{1/2}𝑑t$$ (188) where we have introduced the notation $$s=\left(\frac{\pi n\mathrm{ln}N}{4}\right)^{1/2}r_1$$ (189) to measure the distance $`r_1`$ in terms of the average distance between vortices. Defining the incomplete $`\mathrm{\Gamma }`$-function by: $$\mathrm{\Gamma }_x(p+1)=_0^xe^tt^p𝑑t$$ (190) we have $$V_{1\|\|}=\|𝐕_0\|\left(\frac{n\gamma ^2}{\pi \mathrm{ln}N}\right)^{1/2}_0^+\mathrm{}\frac{J_1(z)}{z^2}\mathrm{\Gamma }_{s^2z^2}\left(\frac{1}{2}\right)𝑑z$$ (191) This result can be compared with formula (117) of Chandrasekhar for the correlation in the gravitational forces acting at two points separated by a finite distance. The dependence upon $`r_1`$, through the variable $`s`$ is encapsulated in the function $$f(s)=_0^+\mathrm{}\frac{J_1(z)}{z^2}\mathrm{\Gamma }_{s^2z^2}\left(\frac{1}{2}\right)𝑑z$$ (192) In Appendix C 2, it is shown that $$f(s)2s(s0)$$ (193) $$f(s)1.41548\mathrm{}+\frac{\sqrt{\pi }}{2}\mathrm{ln}s(s+\mathrm{})$$ (194) This leads to the asymptotic behaviours $$V_{1\|\|}=\|𝐕_0\|n\gamma r_1+\mathrm{}(r_10)$$ (195) $$V_{1\|\|}=\frac{2}{\mathrm{ln}N}\|𝐕_0\|\mathrm{ln}\left(\frac{R}{r_1}\right)(r_1+\mathrm{})$$ (196) consistent with formulae (175) and (178)-(179). We observe that the velocity correlations decay extremely slowly with the distance $`r_1`$ while the initial slope of the function $`V_{1\|\|}`$ goes rapidly to zero at $`dn^{1/2}`$, the interparticle distance (see equation (195)). ### D The function $`\frac{𝐕_0𝐕_1}{V_0^2}`$ We now wish to evaluate the function $$K(r_1)=\frac{𝐕_0𝐕_1}{V_0^2}$$ (197) According to equation (182), we immediately have $$K(r_1)=\frac{n\gamma ^2}{2\pi }_0^+\mathrm{}𝑑V_0_0^+\mathrm{}\rho ^2𝑑\rho e^{nC(\rho )}J_1(\rho V_0)Q(\rho ,r_1)$$ (198) Interchanging the order of integrations and using the identity (82), we obtain $$K(r_1)=\frac{n\gamma ^2}{2\pi }_0^+\mathrm{}\rho 𝑑\rho e^{nC(\rho )}Q(\rho ,r_1)$$ (199) This function is related to equation (187) simply by introducing $`\rho `$ in the integral. We can therefore repeat the steps leading to equation (191) and we obtain instead: $$K(s)=1\frac{4}{\mathrm{ln}N}_0^+\mathrm{}\frac{J_1(z)}{z^2}\mathrm{\Gamma }_{s^2z^2}(1)𝑑z$$ (200) where $`s`$ is defined by equation (189). Explicitly, it has the form: $$K(s)=1\frac{4}{\mathrm{ln}N}_0^+\mathrm{}\frac{J_1(z)}{z^2}(1e^{s^2z^2})𝑑z$$ (201) It turns out that this integral can be expressed in terms of known functions. Indeed: $$K(s)=1\frac{4s^2}{\mathrm{ln}N}\left(1e^{\frac{1}{4s^2}}+\frac{1}{4s^2}E_1\left(\frac{1}{4s^2}\right)\right)$$ (202) where $`E_1(z)`$ denotes the exponential integral $$E_1(z)=_z^+\mathrm{}\frac{e^t}{t}𝑑t$$ (203) It has the series expansion $$E_1(z)=C\mathrm{ln}z\underset{n=1}{\overset{+\mathrm{}}{}}\frac{(1)^nz^n}{n!n}$$ (204) where $`C=0.5772\mathrm{}`$ is the Euler constant. For large values of the argument $`z+\mathrm{}`$, we have: $$E_1(z)\frac{e^z}{z}\left(1\frac{1}{z}+\frac{2}{z^2}\frac{3!}{z^3}+\mathrm{}\right)$$ (205) Using these formulae, we obtain the following behaviour of $`K(r_1)`$ for small and large separations: $$K(r_1)=1\pi nr_1^2+\mathrm{}(r_10)$$ (206) $$K(r_1)=\frac{2}{\mathrm{ln}N}\mathrm{ln}\left(\frac{R}{r_1}\right)(r_1+\mathrm{})$$ (207) ## VII The spatial velocity correlation function and the energy spectrum ### A The correlation function $`𝐕_0𝐕_1`$ The spatial correlation function of the velocity is defined by: $$𝐕_0𝐕_1=W(𝐕_0,𝐕_1)𝐕_0𝐕_1d^2𝐕_0d^2𝐕_1$$ (208) It is related to equation (181) by simply introducing $`V_0`$ in the integral. Therefore, according to equation (185), we have immediately: $$𝐕_0𝐕_1=_0^+\mathrm{}𝑑V_0_0^+\mathrm{}\rho V_0J_0(\rho V_0)\frac{\mathrm{\Phi }^{}(\rho )}{\rho }𝑑\rho $$ (209) where $`\mathrm{\Phi }(\rho )`$ is defined by equation (184). Using the identity (134) and integrating by parts the second integral, we obtain: $$𝐕_0𝐕_1=_0^+\mathrm{}𝑑V_0_0^+\mathrm{}\rho J_1(\rho V_0)\frac{d}{d\rho }\left(\frac{\mathrm{\Phi }^{}(\rho )}{\rho }\right)𝑑\rho $$ (210) Inverting the order of integration and using (82), we find $$𝐕_0𝐕_1=_0^+\mathrm{}\frac{d}{d\rho }\left(\frac{\mathrm{\Phi }^{}(\rho )}{\rho }\right)𝑑\rho $$ (211) The correlation function (208) therefore takes the relatively simple form: $$𝐕_0𝐕_1=\underset{\rho 0}{lim}\frac{\mathrm{\Phi }^{}(\rho )}{\rho }$$ (212) Using equations (184) and (7), the derivative of $`\mathrm{\Phi }`$ is $`\mathrm{\Phi }^{}(\rho )={\displaystyle \frac{n^2\gamma ^4}{16\pi ^2}}\rho ^3\mathrm{ln}\left({\displaystyle \frac{4\pi N}{n\gamma ^2\rho ^2}}\right)e^{nC(\rho )}Q(\rho ,r_1)`$ (213) $`{\displaystyle \frac{n\gamma ^2}{\pi }}\rho e^{nC(\rho )}Q(\rho ,r_1){\displaystyle \frac{n\gamma ^2}{2\pi }}e^{nC(\rho )}Q^{}(\rho ,r_1)\rho ^2`$ (214) For $`\rho 0`$ equation (169) reduces to: $$Q(\rho ,r_1)=\frac{1}{2}\mathrm{ln}\left(\frac{R}{r_1}\right)+O(\rho ^2)$$ (215) Therefore, according to equation (212), we obtain: $$𝐕_0𝐕_1=\frac{n\gamma ^2}{2\pi }\mathrm{ln}\left(\frac{R}{r_1}\right)$$ (216) Like the quantities (181) (197), the velocity autocorrelation function (216) decays extremely slowly at large distances. However, unlike (181) and (197), the velocity autocorrelation function diverges logarithmically at small separations since the variance $`V^2`$ does not exist (see section II). In Appendix D, we give a more direct derivation of formula (216). ### B The energy spectrum of a random distribution of point vortices The correlation function (208) is related to the energy spectrum of the random distribution of point vortices by the Fourier transform $$𝐕(0)𝐕(𝐫)=\frac{1}{4\pi ^2}e^{i\mathrm{𝐤𝐫}}\frac{4\pi E(k)}{k}d^2𝐤$$ (217) Inversely, $$\frac{4\pi E(k)}{k}=e^{i\mathrm{𝐤𝐫}}𝐕(0)𝐕(𝐫)d^2𝐫$$ (218) Introducing polar coordinates and carrying out the integration over the angular variable, we find $$E(k)=\frac{1}{2}_0^+\mathrm{}J_0(kr)𝐕(0)𝐕(𝐫)kr𝑑r$$ (219) According to equation (216), we have for $`rR`$: $$𝐕(0)𝐕(r)=\frac{n\gamma ^2}{2\pi }\mathrm{ln}\left(\frac{R}{r}\right)$$ (220) Hence $$E(k)=\frac{n\gamma ^2}{4\pi }_0^RJ_0(kr)\mathrm{ln}\left(\frac{R}{r}\right)kr𝑑r$$ (221) Using (134) and integrating by parts, we get $$E(k)=\frac{n\gamma ^2}{4\pi }_0^RJ_1(kr)𝑑r$$ (222) Again integrating by parts with the identity (48), we obtain $$E(k)=\frac{n\gamma ^2}{4\pi k}(1J_0(kR))$$ (223) For large $`k`$ the energy spectrum reduces to $$E(k)\frac{n\gamma ^2}{4\pi k}(k+\mathrm{})$$ (224) which is Novikov result. At small $`k`$, using equation (223), we find $$E(k)\frac{N\gamma ^2}{16\pi ^2}k(k0)$$ (225) ## VIII Conclusion In this paper, we have analyzed in some details the statistical features of the stochastic velocity field produced by a random distribution of point vortices in two dimensions. In particular, we have obtained exact results characterizing the correlations in the velocities occuring at two points separated by an arbitrary distance. We have derived an explicit expression for the spatial velocity autocorrelation function and found that the correlations decay extremely slowly with the distance. The other quantities computed in this article are less standard quantities in turbulence but we do not see any reason why they should not be considered in details. They could be measured in simulations of point vortex dynamics and this would provide a direct confrontation between our theoretical model and a situation in which the temporal dynamics of the point vortices is explicitly taken into account. One interest of our model is to provide exact results for the velocity correlations, which is not so frequent in turbulence. Our calculations could be relevant to the context of decaying two dimensional turbulence when the flow becomes dominated by a large number of coherent vortices. Indeed, our model is based on the same assumptions as in Ref. and these assumptions have been vindicated by Direct Navier Stokes simulations and laboratory experiments. In particular, the assumption that the Poisson distribution is stationary is vindicated both by theoretical arguments and by the simulations of Jiménez. Of course, in decaying turbulence the density of vortices varies with time but it should be possible to integrate this dependance in the theory to make our results useful in more general situations. Also, in two dimensional turbulence the vortices have a finite core (vortex “blobs”) and this can severely alter the predictions of the point vortex model. However, our formalism is general and we can introduce a lower cut off in the theory to take into account finite size effects. For the first time, we have obtained explicit expressions for the distribution of velocity and velocity increments taking into account the finite size of the vortices. For “extended” vortices we have proved that the p.d.f. of both velocity and velocity differences are Gaussian. Extended vortices occur in the early stage of 2D decaying turbulence (when the area covered by the vortices is still large ) or when the Reynolds number is low . Gaussian p.d.f. are indeed observed in these situations and have been discussed in Ref. . At the late stages of the decline (when the area covered by the vortices has decreased ) or for high Reynolds numbers the point vortex model should be more and more accurate. It would predict a marginal Gaussian distribution (with an algebraic tail) for the velocity p.d.f. and a Cauchy law for the velocity differences. However, our analytic results show that for “small” but non singular vortices the p.d.f can substantially deviate from the case $`a=0`$ and this can possibly explain (in a quasi-analytical framework) the occurence of exponential tails observed and discussed in Ref. (this point will be developed elsewhere). In the present study, we have exclusively considered the spatial correlations of the velocity occuring between two points at the same time. Another quantity of fundamental interest is the temporal correlation function whose integral determines the diffusion coefficient of point vortices through a Kubo formula . It is not possible, however, to relate the spatial and temporal correlations functions by considering that the point vortices follow linear trajectories with uniform velocity as is commonly done in the case of stars or electric charges . Physically, the linear trajectory approximation made in plasma physics or in stellar dynamics is not applicable here because the vortices, unlike material particles, do not have inertia so they do not move on their own. There is, however, a situation where such an assumption can be implemented. It concerns the motion of point vortices in the presence of a strong background shear as investigated by Chavanis . In that situation we can consider, to a first approximation, that the point vortices follow the streamlines of the shear and evaluate the temporal correlation function within this assumption. The correlation function is found to decay as $`t^2`$ for $`t+\mathrm{}`$. This is a slow decay but sufficient to insure the convergence of the diffusion coefficient. In addition, when the vortices move in a background shear they experience a systematic drift normal to their mean field velocity. The drift coefficient is proportional to the diffusion coefficient (hence to the velocity correlation function) through an Einstein relation, like in Brownian theory. At equilibrium the drift balances the scattering and maintains an inhomogeneous vortex distribution. This drift is of great importance to understand the organization of point vortices at negative temperatures and has been discussed in details elsewhere . In the absence of background shear, we cannot evaluate the temporal correlation function simply but we can estimate the diffusion coefficient of point vortices through the analysis of the velocity fluctuations performed by Chavanis & Sire . Therefore, these results are of great importance to build up a rational kinetic theory of point vortices . ###### Acknowledgements. One of us (PHC) would like to acknowledge interesting discussions with I. Mezic during the program on Hydrodynamic Turbulence at the Institute of Theoretical Physics, Santa Barbara. ## A Detailed calculations of subsection III D In this Appendix, we determine the asymptotic behaviour for $`\|\sigma \sigma \sigma \|0`$ of the function $`A(\rho \rho \rho ,\sigma \sigma \sigma )`$ defined in subsection III. According to equation (32), we need the behaviour of $`C(\rho \rho \rho ,\sigma \sigma \sigma )`$ for $`\|\sigma \sigma \sigma \|0`$. For $`\sigma \sigma \sigma =000`$, $`C(\rho \rho \rho ,\sigma \sigma \sigma )`$ reduces to the function $`C(\rho \rho \rho )`$ introduced in section II. Writing $$C(\rho \rho \rho ,\sigma \sigma \sigma )=C(\rho \rho \rho )+D(\rho \rho \rho ,\sigma \sigma \sigma )$$ (A1) we have $$D(\rho \rho \rho ,\sigma \sigma \sigma )=\frac{\gamma ^2}{4\pi ^2}_{\|𝚽\|=0}^+\mathrm{}e^{i\rho \rho \rho \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }}(1e^{i\sigma \sigma \sigma \psi \psi \psi })\frac{1}{\mathrm{\Phi }^4}d^2𝚽$$ (A2) We have let $`R+\mathrm{}`$ since the integral (A2) is convergent when $`\|\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\|0`$. For $`\|\sigma \sigma \sigma \|0`$, we can expand the exponential term $`e^{i\sigma \sigma \sigma \psi \psi \psi }`$ which occurs under the integral sign in equation (A2) in a power series in $`\sigma `$$`\sigma `$$`\sigma `$. Retaining only the first term in this expansion, we have $$D(\rho \rho \rho ,\sigma \sigma \sigma )=\frac{\gamma ^2}{4\pi ^2}i_{\|𝚽\|=0}^+\mathrm{}e^{i\rho \rho \rho \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }}(\sigma \sigma \sigma \psi \psi \psi )\frac{1}{\mathrm{\Phi }^4}d^2𝚽+o(\|\sigma \sigma \sigma \|^2)$$ (A3) Substituting for $`\psi `$$`\psi `$$`\psi `$ from equation (36) in equation (A3), we obtain $$D(\rho \rho \rho ,\sigma \sigma \sigma )=i\frac{\gamma }{2\pi }_{\|𝚽\|=0}^+\mathrm{}\mathrm{cos}(\rho \rho \rho \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi })\left\{\mathrm{\Phi }^2\sigma \sigma \sigma \delta 𝐫_{}+2(\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }_{}\delta 𝐫)(\sigma \sigma \sigma \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi })\right\}\frac{d^2\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }}{\mathrm{\Phi }^4}+o(\|\sigma \sigma \sigma \|^2)$$ (A4) To evaluate this integral, we introduce a Cartesian system of coordinates where the $`x`$-axis is in the direction of $`\rho `$$`\rho `$$`\rho `$. We denote by $`(\delta r_1,\delta r_2)`$ and $`(\sigma _1,\sigma _2)`$ the components of $`\delta 𝐫`$ and $`\sigma `$$`\sigma `$$`\sigma `$ in this system of coordinates and we introduce $`\theta `$, the angle that $`\mathrm{\Phi }`$$`\mathrm{\Phi }`$$`\mathrm{\Phi }`$ forms with $`\rho `$$`\rho `$$`\rho `$. Transforming to polar coordinates and setting $`x=\rho \mathrm{\Phi }`$, we obtain after some rearrangements: $`D(\rho \rho \rho ,\sigma \sigma \sigma )=i{\displaystyle \frac{\gamma }{2\pi }}{\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dx}{x}}\mathrm{cos}(x\mathrm{cos}\theta )`$ (A5) $`\times \left\{(\sigma _1\delta r_2+\sigma _2\delta r_1)\mathrm{cos}(2\theta )+(\sigma _2\delta r_2\sigma _1\delta r_1)\mathrm{sin}(2\theta )\right\}+o(\|\sigma \sigma \sigma \|^2)`$ (A6) Using the expansion $$\mathrm{cos}(x\mathrm{cos}\theta )=J_0(x)+2\underset{n=1}{\overset{+\mathrm{}}{}}(1)^nJ_{2n}(x)\mathrm{cos}(2n\theta )$$ (A7) and the identities $$_0^{2\pi }\mathrm{cos}(m\theta )\mathrm{cos}(n\theta )𝑑\theta =\pi \delta _{nm}$$ (A8) $$_0^{2\pi }\mathrm{sin}(m\theta )\mathrm{sin}(n\theta )𝑑\theta =\pi \delta _{nm}$$ (A9) $$_0^{2\pi }\mathrm{cos}(m\theta )\mathrm{sin}(n\theta )𝑑\theta =0$$ (A10) we find that $$D(\rho \rho \rho ,\sigma \sigma \sigma )=i\gamma (\sigma _1\delta r_2+\sigma _2\delta r_1)_0^+\mathrm{}\frac{J_2(x)}{x}𝑑x+o(\|\sigma \sigma \sigma \|^2)$$ (A11) Since $$_0^+\mathrm{}J_2(x)\frac{dx}{x}=\frac{1}{2}$$ (A12) we finally obtain $$D(\rho \rho \rho ,\sigma \sigma \sigma )=i\frac{\gamma }{2}(\sigma _2\delta r_1+\sigma _1\delta r_2)+o(\|\sigma \sigma \sigma \|^2)$$ (A13) Now, according to (32)(A1) and (A13), we have: $$A(\rho \rho \rho ,\sigma \sigma \sigma )=e^{nC(\rho \rho \rho )+i\frac{\gamma n}{2}(\sigma _2\delta r_1+\sigma _1\delta r_2)+o(\|\sigma \sigma \sigma \|^2)}$$ (A14) where we recall that $`(\delta r_1,\delta r_2)`$ and $`(\sigma _1,\sigma _2)`$ are the components of $`\delta 𝐫`$ and $`\sigma `$$`\sigma `$$`\sigma `$ in a system of coordinates where the $`x`$-axis is in the direction of $`\rho `$$`\rho `$$`\rho `$. Now, to evaluate the integral (67) we must express $`\delta 𝐫`$ and $`\sigma `$$`\sigma `$$`\sigma `$ in a fixed system of coordinates independent on $`\rho `$$`\rho `$$`\rho `$. We choose this system such that the $`x`$-axis coincides with the direction of $`\delta 𝐫`$. If $`\theta `$ denotes the angle that $`\rho `$$`\rho `$$`\rho `$ forms with $`\delta 𝐫`$, the components $`(\sigma _x,\sigma _y)`$ of $`\sigma `$$`\sigma `$$`\sigma `$ in this system are related to $`(\sigma _1,\sigma _2)`$ by the transformations: $$\sigma _1=\sigma _x\mathrm{cos}\theta +\sigma _y\mathrm{sin}\theta $$ (A15) $$\sigma _2=\sigma _y\mathrm{cos}\theta \sigma _x\mathrm{sin}\theta $$ (A16) Using in addition $`\delta r_1=\|\delta 𝐫\|\mathrm{cos}\theta `$ and $`\delta r_2=\|\delta 𝐫\|\mathrm{sin}\theta `$, $`A(\rho \rho \rho ,\sigma \sigma \sigma )`$ can be rewritten in the form (68). ## B Detailed calculations of subsection IV A In this Appendix, we determine the expression of the function $`D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )`$ defined by equation (102). This function can be written alternatively: $$D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )=\frac{i\gamma }{2\pi }_{\|𝐫\|=0}^Re^{i\frac{\gamma }{2\pi }\frac{\rho \rho \rho _{}𝐫}{r^2}}\sigma \sigma \sigma _{}(\mathrm{ln}\|𝐫𝐫_1\|)d^2𝐫$$ (B1) Integrating by parts, we find $`D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )={\displaystyle \frac{i\gamma }{2\pi }}{\displaystyle _{\|𝐫\|=0}^R}\left\{e^{i\frac{\gamma }{2\pi }\frac{\rho \rho \rho _{}𝐫}{r^2}}\mathrm{ln}\|𝐫𝐫_1\|\sigma \sigma \sigma _{}\right\}d^2𝐫`$ (B2) $`+{\displaystyle \frac{i\gamma }{2\pi }}{\displaystyle _{\|𝐫\|=0}^R}\mathrm{ln}\|𝐫𝐫_1\|\sigma \sigma \sigma _{}\left(e^{i\frac{\gamma }{2\pi }\frac{\rho \rho \rho _{}𝐫}{r^2}}\right)d^2𝐫`$ (B3) Explicitly, it has the form $`D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )={\displaystyle \frac{i\gamma }{2\pi }}{\displaystyle _{\|𝐫\|=R}}e^{i\frac{\gamma }{2\pi }\frac{\rho \rho \rho _{}𝐫}{r^2}}\mathrm{ln}\|𝐫𝐫_1\|\sigma \sigma \sigma _{}dl\widehat{𝐞}_r`$ (B4) $`{\displaystyle \frac{\gamma ^2}{4\pi ^2}}{\displaystyle _{\|𝐫\|=0}^R}e^{i\frac{\gamma }{2\pi }\frac{\rho \rho \rho _{}𝐫}{r^2}}\mathrm{ln}\|𝐫𝐫_1\|\left\{{\displaystyle \frac{\rho \rho \rho _{}\sigma \sigma \sigma _{}}{r^2}}{\displaystyle \frac{2(\rho \rho \rho _{}𝐫)(\sigma \sigma \sigma _{}𝐫)}{r^4}}\right\}d^2𝐫`$ (B5) The first integral is a boundary term which behaves like $`\mathrm{ln}R`$ when $`R\mathrm{}`$; we call it: $$D_{boundary}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )=\frac{i\gamma }{2\pi }_{\|𝐫\|=R}e^{i\frac{\gamma }{2\pi }\frac{\rho \rho \rho _{}𝐫}{r^2}}\mathrm{ln}\|𝐫𝐫_1\|\sigma \sigma \sigma _{}dl\widehat{𝐞}_r$$ (B6) The second integral is convergent when $`R\mathrm{}`$; we call it: $$D_{int}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )=\frac{\gamma ^2}{4\pi ^2}_{\|𝐫\|=0}^+\mathrm{}e^{i\frac{\gamma }{2\pi }\frac{\rho \rho \rho _{}𝐫}{r^2}}\mathrm{ln}\|𝐫𝐫_1\|\left\{\frac{\rho \rho \rho _{}\sigma \sigma \sigma _{}}{r^2}\frac{2(\rho \rho \rho _{}𝐫)(\sigma \sigma \sigma _{}𝐫)}{r^4}\right\}d^2𝐫$$ (B7) Then $$D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )=D_{boundary}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )+D_{int}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )$$ (B8) ### 1 The evaluation of $`D_{int}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )`$ To evaluate $`D_{int}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )`$, we introduce a Cartesian system of coordinates with the $`x`$-axis is the direction of $`\rho \rho \rho _{}`$. We denote by $`(\sigma _1,\sigma _2)`$ the components of $`\sigma `$$`\sigma `$$`\sigma `$ in this system. Introducing $`\theta `$, the angle that $`𝐫`$ forms with $`\rho \rho \rho _{}`$ and transforming to polar coordinates, we obtain: $$D_{int}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )=\frac{\gamma ^2}{4\pi ^2}\rho _0^{2\pi }𝑑\theta _0^+\mathrm{}\frac{dr}{r}e^{i\frac{\gamma \rho \mathrm{cos}\theta }{2\pi r}}\mathrm{ln}\|𝐫𝐫_1\|\{\sigma _2+2\mathrm{cos}\theta (\sigma _1\mathrm{sin}\theta \sigma _2\mathrm{cos}\theta )\}$$ (B9) or, equivalently, $$D_{int}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )=\frac{\gamma ^2}{4\pi ^2}\rho _0^{2\pi }𝑑\theta _0^+\mathrm{}\frac{dr}{r}e^{i\frac{\gamma \rho \mathrm{cos}\theta }{2\pi r}}\mathrm{ln}\|𝐫𝐫_1\|\{\sigma _1\mathrm{sin}(2\theta )\sigma _2\mathrm{cos}(2\theta )\}$$ (B10) We now expand $`e^{i\frac{\gamma \rho \mathrm{cos}\theta }{2\pi r}}`$ and $`\mathrm{ln}\|𝐫𝐫_1\|`$ which occur under the integral sign in terms of sinusoidal functions using the identities: $$e^{ix\mathrm{cos}\theta }=J_0(x)+2\underset{n=1}{\overset{+\mathrm{}}{}}i^nJ_n(x)\mathrm{cos}(n\theta )$$ (B11) $$\mathrm{ln}\|𝐫𝐫_1\|=\mathrm{ln}r_>\underset{m=1}{\overset{+\mathrm{}}{}}\frac{1}{m}\left(\frac{r_<}{r_>}\right)^m\mathrm{cos}[m(\theta \theta _1)]$$ (B12) where $`J_n(x)`$ is the Bessel function of order $`n`$, $`r_>`$ (resp. $`r_<`$) is the larger (resp. smaller) of $`(r,r_1)`$ and $`\theta _1`$ is the angle that $`𝐫_1`$ forms with $`\rho \rho \rho _{}`$. Using the orthogonality properties (A8) (A9) (A10), we can express $`D_{int}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )`$ in the form: $`D_{int}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )={\displaystyle \frac{\gamma ^2}{2\pi }}\rho \sigma _2{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dr}{r}}J_2\left({\displaystyle \frac{\gamma \rho }{2\pi r}}\right)\mathrm{ln}r_>`$ (B13) $`\sigma _2{\displaystyle \underset{n=1}{\overset{+\mathrm{}}{}}}A_n(\rho ,r_1)\mathrm{cos}(n\theta _1)\sigma _1{\displaystyle \underset{n=1}{\overset{+\mathrm{}}{}}}B_n(\rho ,r_1)\mathrm{sin}(n\theta _1)`$ (B14) where: $$A_n(\rho ,r_1)=\frac{\gamma ^2}{4\pi }\rho \frac{i^n}{n}_0^+\mathrm{}\frac{dr}{r}\left(\frac{r_<}{r_>}\right)^n\left[J_{n+2}\left(\frac{\gamma \rho }{2\pi r}\right)+J_{n2}\left(\frac{\gamma \rho }{2\pi r}\right)\right]$$ (B15) $$B_n(\rho ,r_1)=\frac{\gamma ^2}{4\pi }\rho \frac{i^n}{n}_0^+\mathrm{}\frac{dr}{r}\left(\frac{r_<}{r_>}\right)^n\left[J_{n+2}\left(\frac{\gamma \rho }{2\pi r}\right)J_{n2}\left(\frac{\gamma \rho }{2\pi r}\right)\right]$$ (B16) ### 2 The evaluation of $`D_{boundary}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )`$ Using the same system of coordinates, we can rewrite our expression for $`D_{boundary}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )`$ in the form: $$D_{boundary}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )=i\frac{\gamma }{2\pi }R_0^{2\pi }𝑑\theta e^{i\frac{\gamma \rho \mathrm{cos}\theta }{2\pi r}}\mathrm{ln}\|𝐫𝐫_1\|(\sigma _1\mathrm{sin}\theta \sigma _2\mathrm{cos}\theta )$$ (B17) Substituting for (B11) and (B12) in equation (B17) and carrying out the angular integration, we obtain $`D_{boundary}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )=\gamma R\sigma _2J_1\left({\displaystyle \frac{\gamma \rho }{2\pi R}}\right)\mathrm{ln}R`$ (B18) $`{\displaystyle \frac{1}{2}}\gamma R\sigma _2{\displaystyle \underset{n=1}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{i^n}{n}}\left({\displaystyle \frac{r_1}{R}}\right)^n\left[J_{n1}\left({\displaystyle \frac{\gamma \rho }{2\pi R}}\right)J_{n+1}\left({\displaystyle \frac{\gamma \rho }{2\pi R}}\right)\right]\mathrm{cos}(n\theta _1)`$ (B19) $`+{\displaystyle \frac{1}{2}}\gamma R\sigma _1{\displaystyle \underset{n=1}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{i^n}{n}}\left({\displaystyle \frac{r_1}{R}}\right)^n\left[J_{n1}\left({\displaystyle \frac{\gamma \rho }{2\pi R}}\right)+J_{n+1}\left({\displaystyle \frac{\gamma \rho }{2\pi R}}\right)\right]\mathrm{sin}(n\theta _1)`$ (B20) When $`R\mathrm{}`$, the function $`D_{boundary}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )`$ takes the relatively simple form: $$D_{boundary}^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )=\frac{\gamma ^2}{4\pi }\rho \sigma _2\mathrm{ln}R+i\frac{\gamma }{2}r_1(\sigma _2\mathrm{cos}\theta _1+\sigma _1\mathrm{sin}\theta _1)$$ (B21) ### 3 Final expression of $`D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )`$ In equations (B14) and (B21), $`\sigma \sigma \sigma =(\sigma _1,\sigma _2)`$ is refered to a variable system of coordinates depending on the direction of $`\rho \rho \rho _{}`$. Since $`W(𝐕_0,𝐕_1)`$ is the Fourier transform of $`A(\rho \rho \rho ,\sigma \sigma \sigma )`$, related to $`D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )`$ by equation (103), we need to express the vectors $`\rho `$$`\rho `$$`\rho `$ and $`\sigma `$$`\sigma `$$`\sigma `$ in a fixed system of coordinates. We choose this system such that the $`x`$-axis is in the direction of $`𝐫_1`$. The components $`(\sigma _x,\sigma _y)`$ of $`\sigma `$$`\sigma `$$`\sigma `$ in this system are related to $`(\sigma _1,\sigma _2)`$ by the transformation: $`\sigma _1=\sigma _x\mathrm{cos}\theta _1\sigma _y\mathrm{sin}\theta _1`$ (B22) $`\sigma _2=\sigma _x\mathrm{sin}\theta _1+\sigma _y\mathrm{cos}\theta _1`$ (B23) where it might be recalled that $`\theta _1`$ is the angle that $`\rho \rho \rho _{}`$ forms with $`𝐫_1`$. Thus, in this new system of coordinates, $`D^{(1)}(\rho \rho \rho ,\sigma \sigma \sigma )`$ has the form (107). ## C Detailed calculations of subsection VI C ### 1 Alternative form of equation (187) In this Appendix, we derive an alternative form of equation (187) in terms of $`\mathrm{\Gamma }`$ functions. Substituting for $`Q(\rho ,r_1)`$ from equation (169) in equation (187), we obtain $$V_{1\|\|}=\frac{n\gamma ^2}{4\pi }_0^+\mathrm{}𝑑\rho e^{nC(\rho )}\left\{\mathrm{ln}\left(\frac{\gamma \rho }{2\pi R}\right)\mathrm{ln}ϵ\right\}+I^{(ϵ)}$$ (C1) where $$I^{(ϵ)}=\frac{n\gamma ^2}{2\pi }_{\frac{2\pi r_1}{\gamma }ϵ}^+\mathrm{}𝑑\rho e^{nC(\rho )}_ϵ^{\frac{\gamma \rho }{2\pi r_1}}\frac{J_1(z)}{z^2}𝑑z$$ (C2) and we recall that $`ϵ0`$. Inverting the order of integration, we have $$I^{(ϵ)}=\frac{n\gamma ^2}{2\pi }_ϵ^+\mathrm{}𝑑z\frac{J_1(z)}{z^2}_{\frac{2\pi r_1}{\gamma }z}^+\mathrm{}e^{nC(\rho )}𝑑\rho $$ (C3) or, alternatively $$I^{(ϵ)}=\frac{n\gamma ^2}{2\pi }_ϵ^+\mathrm{}\frac{J_1(z)}{z^2}𝑑z_{\frac{2\pi r_1}{\gamma }ϵ}^+\mathrm{}e^{nC(\rho )}𝑑\rho \frac{n\gamma ^2}{2\pi }_ϵ^+\mathrm{}𝑑z\frac{J_1(z)}{z^2}_{\frac{2\pi r_1}{\gamma }ϵ}^{\frac{2\pi r_1}{\gamma }z}e^{nC(\rho )}𝑑\rho $$ (C4) Using the identity (146), integrating by parts the first integral in equation (C4) and taking the limit $`ϵ0`$, we obtain $`I^{(ϵ)}={\displaystyle \frac{n\gamma ^2}{2\pi }}\left\{{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{J_2(z)}{z}}\mathrm{ln}zdz{\displaystyle \frac{1}{2}}\mathrm{ln}ϵ\right\}{\displaystyle _0^+\mathrm{}}e^{nC(\rho )}𝑑\rho `$ (C5) $`{\displaystyle \frac{n\gamma ^2}{2\pi }}{\displaystyle _0^+\mathrm{}}𝑑z{\displaystyle \frac{J_1(z)}{z^2}}{\displaystyle _0^{\frac{2\pi r_1}{\gamma }z}}e^{nC(\rho )}𝑑\rho `$ (C6) According to equations (C1) and (C6) we therefore have $`V_{1\|\|}={\displaystyle \frac{n\gamma ^2}{4\pi }}{\displaystyle _0^+\mathrm{}}e^{nC(\rho )}\mathrm{ln}\left({\displaystyle \frac{\gamma \rho }{2\pi R}}\right)𝑑\rho {\displaystyle \frac{n\gamma ^2}{2\pi }}{\displaystyle _0^+\mathrm{}}𝑑z{\displaystyle \frac{J_1(z)}{z^2}}{\displaystyle _0^{\frac{2\pi r_1}{\gamma }z}}𝑑\rho e^{nC(\rho )}`$ (C7) For $`r_1=0`$, there remains $$V_{0\|\|}=\frac{n\gamma ^2}{4\pi }_0^+\mathrm{}e^{nC(\rho )}\mathrm{ln}\left(\frac{\gamma \rho }{2\pi R}\right)𝑑\rho =\left(\frac{n\gamma ^2}{16}\mathrm{ln}N\right)^{1/2}$$ (C8) This is precisely the value of $`\|𝐕_0\|`$, see equation (19), as expected from equation (181). Therefore, we can write $$V_{1\|\|}=\|𝐕_0\|\frac{n\gamma ^2}{2\pi }_0^+\mathrm{}𝑑z\frac{J_1(z)}{z^2}_0^{\frac{2\pi r_1}{\gamma }z}e^{nC(\rho )}𝑑\rho $$ (C9) In the last integral, we can replace the function $`C(\rho )`$ by its approximate expression (8) without introducing any significant error. Therefore $$V_{1\|\|}=\|𝐕_0\|\frac{n\gamma ^2}{2\pi }_0^+\mathrm{}𝑑z\frac{J_1(z)}{z^2}_0^{\frac{2\pi r_1}{\gamma }z}e^{\frac{n\gamma ^2}{16\pi }\mathrm{ln}N\rho ^2}𝑑\rho $$ (C10) With the change of variables $`t=\frac{n\gamma ^2}{16\pi }\mathrm{ln}N\rho ^2`$, it takes the form (188). ### 2 The function $`f(s)`$ The behaviour of $`f(s)`$ for $`s0`$ can be derived from equation (192) by expanding the incomplete $`\mathrm{\Gamma }`$-function which occurs under the integral sign as a power series in $`s`$. To first order in $`s`$, we have $`\mathrm{\Gamma }_{s^2z^2}\left({\displaystyle \frac{1}{2}}\right)={\displaystyle _0^{s^2z^2}}e^tt^{1/2}𝑑t`$ (C11) $`={\displaystyle _0^{s^2z^2}}\left(1t+{\displaystyle \frac{t^2}{2}}+\mathrm{}\right)t^{1/2}𝑑t=2sz+o(s^3)`$ (C12) Thus $$f(s)=_0^+\mathrm{}\frac{J_1(z)}{z^2}\mathrm{\Gamma }_{s^2z^2}\left(\frac{1}{2}\right)𝑑z=2s_0^+\mathrm{}𝑑z\frac{J_1(z)}{z}+\mathrm{}$$ (C13) With the identity (49), we obtain $$f(s)2s(s0)$$ (C14) To obtain the behaviour of $`f(s)`$ for $`s+\mathrm{}`$, we write equation (192) in the form $$f(s)=_0^+\mathrm{}𝑑z\frac{J_1(z)}{z^2}_0^{s^2z^2}e^tt^{1/2}𝑑t$$ (C15) and take the derivative with respect to $`s`$: $$f^{}(s)=2_0^+\mathrm{}\frac{J_1(z)}{z}e^{s^2z^2}𝑑z$$ (C16) Replacing $`J_1(z)`$ in the foregoing equation by its series expansion $$J_1(z)=\frac{1}{2}z\underset{n=0}{\overset{+\mathrm{}}{}}\frac{(1)^n}{n!(n+1)!}\left(\frac{z}{2}\right)^{2n}$$ (C17) and inverting the order of the integration and the summation, we obtain: $$f^{}(s)=\underset{n=0}{\overset{+\mathrm{}}{}}\frac{(1)^n}{n!(n+1)!}_0^+\mathrm{}\left(\frac{z}{2}\right)^{2n}e^{s^2z^2}𝑑z$$ (C18) If we now introduce the variable $`t=s^2z^2`$, equation (C18) becomes $$f^{}(s)=\underset{n=0}{\overset{+\mathrm{}}{}}\frac{(1)^n}{n!(n+1)!}\frac{1}{(2s)^{2n+1}}_0^+\mathrm{}t^{n\frac{1}{2}}e^t𝑑t$$ (C19) or, alternatively $$f^{}(s)=\underset{n=0}{\overset{+\mathrm{}}{}}\frac{(1)^n}{n!(n+1)!}\frac{1}{(2s)^{2n+1}}\mathrm{\Gamma }\left(n+\frac{1}{2}\right)$$ (C20) Integrating this series term by term, we get $$f(s)=K+\frac{\sqrt{\pi }}{2}\mathrm{ln}s\frac{1}{4}\underset{n=1}{\overset{+\mathrm{}}{}}\frac{(1)^n}{n!(n+1)!}\frac{1}{n}\frac{1}{(2s)^{2n}}\mathrm{\Gamma }\left(n+\frac{1}{2}\right)$$ (C21) where $`K`$ is a constant. Numerically, we find $`K=1.41548\mathrm{}`$. Therefore, to leading order in $`s`$, we can write $$f(s)=1.41548\mathrm{}+\frac{\sqrt{\pi }}{2}\mathrm{ln}s(s+\mathrm{})$$ (C22) ## D An alternative derivation of equation (216) The average velocity occurring in $`O`$ is given by the expression $$𝐕_0=W(𝐕_0,𝐕_1)𝐕_0d^2𝐕_0d^2𝐕_1$$ (D1) Using equation (95) and integrating by parts, we obtain $$𝐕_0=i\frac{A}{\rho \rho \rho }(\mathrm{𝟎},\mathrm{𝟎})=in\frac{C}{\rho \rho \rho }(\mathrm{𝟎},\mathrm{𝟎})$$ (D2) Similarly, $$𝐕_1=i\frac{A}{\sigma \sigma \sigma }(\mathrm{𝟎},\mathrm{𝟎})=in\frac{C}{\sigma \sigma \sigma }(\mathrm{𝟎},\mathrm{𝟎})$$ (D3) The same method applied to the correlation function $$𝐕_0𝐕_1=W(𝐕_0,𝐕_1)𝐕_0𝐕_1d^2𝐕_0d^2𝐕_1$$ (D4) yields $$𝐕_0𝐕_1=\frac{^2A}{\rho \rho \rho \sigma \sigma \sigma }(\mathrm{𝟎},\mathrm{𝟎})=n^2\frac{C}{\rho \rho \rho }(\mathrm{𝟎},\mathrm{𝟎})\frac{C}{\sigma \sigma \sigma }(\mathrm{𝟎},\mathrm{𝟎})+n\frac{^2C}{\rho \rho \rho \sigma \sigma \sigma }(\mathrm{𝟎},\mathrm{𝟎})$$ (D5) Hence, we can write $$𝐕_0𝐕_1=𝐕_0𝐕_1+n\frac{^2C}{\rho \rho \rho \sigma \sigma \sigma }(\mathrm{𝟎},\mathrm{𝟎})$$ (D6) Since $`𝐕_0=\mathrm{𝟎}`$, there remains $$𝐕_0𝐕_1=n\frac{^2C}{\rho \rho \rho \sigma \sigma \sigma }(\mathrm{𝟎},\mathrm{𝟎})$$ (D7) Substituting explicitly for $`C(\rho \rho \rho ,\sigma \sigma \sigma )`$ from equation (98) in equations (D3) and (D7), we find $$𝐕_1=\frac{n\gamma }{2\pi }\frac{(𝐫𝐫_1)_{}}{\|𝐫𝐫_1\|^2}d^2𝐫$$ (D8) and $$𝐕_0𝐕_1=\frac{n\gamma ^2}{4\pi ^2}\frac{𝐫_{}}{r^2}\frac{(𝐫𝐫_1)_{}}{\|𝐫𝐫_1\|^2}d^2𝐫$$ (D9) By part integration, the first integral can be rewritten $$𝐕_1=\frac{n\gamma }{2\pi }\widehat{𝐳}\mathrm{ln}\|𝐫𝐫_1\|\widehat{𝐞}_rdl$$ (D10) where $`\widehat{𝐳}`$ is a unit vector normal to the flow and $`\widehat{𝐞}_r`$ is a unit vector normal to the disk boundary at $`r=R`$. Using the expansion (B12) for $`\mathrm{ln}\|𝐫𝐫_1\|`$, we obtain easily $$𝐕_1=\frac{1}{2}n\gamma 𝐫_1$$ (D11) which corresponds to the solid rotation (137). Integrating by parts the second integral, we get $$𝐕_0𝐕_1=\frac{n\gamma ^2}{4\pi ^2}\left\{\mathrm{ln}\|𝐫𝐫_1\|\frac{𝐫}{r^2}\widehat{𝐞}_rdl2\pi \delta (𝐫)\mathrm{ln}\|𝐫𝐫_1\|d^2𝐫\right\}$$ (D12) Using equation (B12), we obtain $$𝐕_0𝐕_1=\frac{n\gamma ^2}{2\pi }\mathrm{ln}\left(\frac{R}{r_1}\right)$$ (D13) which coincides with formula (216). Of course, the relations (D8) and (D9) can be obtained independently of the formalism introduced in this article. Using equations (93) and (94), we can write $$𝐕_0𝐕_1=\frac{\gamma ^2}{4\pi ^2}\underset{i,j}{}\frac{𝐫_i}{r_i^2}\frac{(𝐫_j𝐫_1)_{}}{\|𝐫_j𝐫_1\|^2}$$ (D14) where $`X`$ denotes the average value $`\tau (𝐫_1)\mathrm{}\tau (𝐫_N)Xd^2𝐫_1\mathrm{}d^2𝐫_N`$ associated with a decorrelated distribution of point vortices. If we split the sum in two parts, distinguishing the contribution of vortex pairs from the contributions of individual vortices, we obtain: $$𝐕_0𝐕_1=\frac{\gamma ^2}{4\pi ^2}\underset{i=1}{\overset{N}{}}\frac{𝐫_i}{r_i^2}\frac{(𝐫_i𝐫_1)_{}}{\|𝐫_i𝐫_1\|^2}+\frac{\gamma ^2}{4\pi ^2}\underset{i=1}{\overset{N}{}}\underset{ji}{}\frac{𝐫_i}{r_i^2}\frac{(𝐫_j𝐫_1)_{}}{\|𝐫_j𝐫_1\|^2}$$ (D15) Since all vortices are identical, we find $`𝐕_0𝐕_1={\displaystyle \frac{\gamma ^2}{4\pi ^2}}N{\displaystyle \tau (𝐫)\frac{𝐫_{}}{r^2}\frac{(𝐫𝐫_1)_{}}{\|𝐫𝐫_1\|^2}d^2𝐫}`$ (D16) $`+{\displaystyle \frac{\gamma ^2}{4\pi ^2}}N(N1){\displaystyle \tau (𝐫)\frac{𝐫_{}}{r^2}d^2𝐫\tau (𝐫)\frac{(𝐫𝐫_1)_{}}{\|𝐫𝐫_1\|^2}d^2𝐫}`$ (D17) Using $`\tau (𝐫)=\frac{1}{\pi R^2}`$ and making the approximation $`N(N1)N^2`$ for large $`N`$’s, we get $$𝐕_0𝐕_1=\frac{n\gamma ^2}{4\pi ^2}\frac{𝐫_{}}{r^2}\frac{(𝐫𝐫_1)_{}}{\|𝐫𝐫_1\|^2}d^2𝐫+\frac{n^2\gamma ^2}{4\pi ^2}\frac{𝐫_{}}{r^2}d^2𝐫\frac{(𝐫𝐫_1)_{}}{\|𝐫𝐫_1\|^2}d^2𝐫$$ (D18) The second term is just the product $`𝐕_0𝐕_1`$ which is equal to zero since $`𝐕_0=\mathrm{𝟎}`$. Therefore, equation (D18) returns our previous expression (D9). Accordingly, the formula (D13) for the velocity autocorrelation function can be derived very simply, independantly from the general formalism introduced in the article. Note, in contrast, that $`V_{1\|\|}`$ and $`𝐕_0𝐕_1/V_0^2`$ cannot be obtained by this method. Therefore, the general theory developed in sections VI C and VI D is necessary to compute these more complicated quantities.
warning/0004/math0004164.html	ar5iv	text	# No More than Three Favourite Sites for Simple Random Walk We prove that, with probability one, eventually there are no more than three favourite (i.e. most visited) sites of simple symmetric random walk. This partially answers a relatively long standing question of Pál Erdős and Pál Révész. Key Words: Random walks, local time, favourite sites, most visited sites 1991 AMS Subject Classification: 60J15, 60J55 1. Introduction and Main Result. Let $`S(t)`$, $`t_+`$ be a simple symmetric random walk on $``$ with initial state $`S(0)=0`$. Its upcrossings, downcrossings and (site) local time are defined for $`t`$ and $`x`$ as follows: $`U(t,x)`$ $`:=\mathrm{\#}\{0<st:S(s)=x,S(s1)=x1\},`$ $`1.1`$$`1.2`$$`1.3`$ $`D(t,x)`$ $`:=\mathrm{\#}\{0<st:S(s)=x,S(s1)=x+1\},`$ $`L(t,x)`$ $`:=\mathrm{\#}\{0<st:S(s)=x\}=U(t,x)+D(t,x).`$ The following identities are straightforward: $`U(t,x)D(t,x1)`$ $`=11_{\{0<xS(t)\}}11_{\{S(t)<x0\}}`$ $`1.4`$$`1.5`$ $`D(t,x)U(t,x+1)`$ $`=11_{\{S(t)x<0\}}11_{\{0x<S(t)\}}.`$ And from these it follows that $`L(t,x)`$ $`=D(t,x)+D(t,x1)+11_{\{0<xS(t)\}}11_{\{S(t)<x0\}}`$ $`1.6`$ $`=U(t,x)+U(t,x+1)+11_{\{S(t)x<0\}}11_{\{0x<S(t)\}}.`$ The set of favourite (or: most visited) sites of the random walk at time $`t`$, are those sites where the local time attains its maximum value: $$𝒦(t):=\{y:L(t,y)=\underset{z}{\mathrm{max}}L(t,z)\}.$$ $`1.7`$ It is clear that the number of favourite sites changes in time as follows: $$\mathrm{\#}𝒦(t+1)=\{\begin{array}{ccc}\mathrm{\#}𝒦(t)& \text{ if }& S(t+1)𝒦(t+1)\\ \mathrm{\#}𝒦(t)+1& \text{ if }& 𝒦(t+1)=𝒦(t)\{S(t+1)\}\\ 1& \text{ if }& 𝒦(t+1)=\{S(t+1)\}𝒦(t).\end{array}$$ $`1.8`$ In plain words one of the following three possibilities can occur at each step of the walk: Either the currently occupied site is not favourite and $`𝒦`$ remains unchanged. Or the currently occupied site becomes a new favourite beside the favourites of the previous stage, and thus the number of favourites increases by one. Or, finally, a favourite site is visited and so this site becomes now the only new favourite. No other possibility. Clearly $`\mathrm{\#}𝒦(t)1`$ for all $`t1`$, and it is easy to verify that for infinitely many times, $`t1`$, there are at least two favourite sites: $`\mathrm{\#}𝒦(t)2`$. Pál Erdős and Pál Révész formulated and repeatedly raised the following Question: Does it happen that $`\mathrm{\#}𝒦(t)r`$ infinitely often (i.e. almost surely for infinitely many times $`t1`$) for $`r=3,4,\mathrm{}`$? See e.g. Erdős and Révész (1984), (1987), (1991), Erdős (1994) or Révész (1990) for an extended list of related questions and problems. Questions related to the asymptotic behaviour of the favourite (or most visited) sites of a random walk were considered by many authors since the mid-eighties. We quote here a few relevant results, with no claim of exhaustiveness. $``$ Bass and Griffin (1985) prove that almost surely, the set of favourites is transient. More exactly: they prove that the distance of the set of favourite sites from the origin increases faster than $`\sqrt{n}/(\mathrm{log}n)^{11}`$ but slower than $`\sqrt{n}/(\mathrm{log}n)`$. $``$ Csáki and Shi (1998) prove that the distance between the edge of the range of the random walk and the set of favourite sites increases as fast as $`\sqrt{n}/(\mathrm{log}\mathrm{log}n)^{3/2}`$ $``$ Csáki, Révész and Shi (2000) prove that the position of a favourite site can have as large jumps as $`\sqrt{2n\mathrm{log}\mathrm{log}n}`$, i.e., comparable with the diameter of the full range of the random walk. They also extend a much earlier result of Kesten (1965), identifying the set of joint limit points of the set of favourite sites and the favourite values (i.e. max. values of local time), both rescaled by $`\sqrt{2n\mathrm{log}\mathrm{log}n}`$ $``$ There are many papers dealing with similar qestions in the context of symmetric stable processes rather than random walks (or Brownian motion). See, e.g., Eisenbaum (1997), Bass, Eisenbaum and Shi (2000) and other papers quoted there. In the present paper we answer in the negative the question of Erdős and Révész quoted above, for $`r4`$: we prove that with probability 1, there are at most finitely many times $`t1`$ when there are four or more favourite sites of the random walk $`S(t)`$. In Tóth and Werner (1997) a similar result was proved for the set of favourite edges rather than favourite sites. The present paper deals with the original question of Erdős and Révész. The starting general ideas of the present paper (see Sections 1-3) are very close to those of Tóth and Werner (1997). However: the details of the proof require more refined estimates and arguments. On the technical level (see Sections 4-6) this proof is rather different. For $`r1`$ denote by $`f(r)`$ the (possibly infinite) number of steps, when the currently occupied site is one of the $`r`$ actual favourites: $$f(r):=\mathrm{\#}\{t1:[S(t)𝒦(t)][\mathrm{\#}𝒦(t)=r]\}.$$ $`1.9`$ From (1.8) it follows that for any $`r1`$, $`f(r+1)f(r)`$. (Both sides of the inequality could be infinite.) The main result of this paper is the following ###### Theorem 1 $$𝔼\left(f(4)\right)<\mathrm{}.$$ $`1.10`$ Remarks: (1) From this theorem the negative answer to the question of Erdős and Révész clearly follows, for the cases $`r4`$. (2) The case $`r=3`$ remains open. From the proof of the above theorem it becomes clear that $`𝔼\left(f(3)\right)=\mathrm{}`$. Nevertheless we conjecture that $`f(3)<\mathrm{}`$, almost surely. The paper is organized as follows: In Section 2 we perform some straightforward manipulations (essentially: rearrangements of sums). In Section 3 we recall the Ray-Knight Theorems for the local times of simple random walks. In Section 4 first we express our relevant probabilities and expectations (found in Section 2) in terms of the Galton-Watson processes arising with the Ray-Knight representation. Then we formulate Proposition 1, stating some upper bounds on these probabilities and expectations, and using these bounds we prove Theorem 1. The proof of Proposition 1 is postponed to the end of Section 5. In Section 5 four lemmas and, as their consequence, Proposition 1 are proved. Throughout the technical parts of the proofs smaller, quite plausible statements are invoked. These are called Side-lemmas (1 to 7). Their proofs are postponed to Section 6. Throughout the paper, in various upper bounds, multiplicative constants, respectively, constants in exponential rates will be denoted generically by $`C`$, respectively, by $`\gamma `$. The values of these constants may vary even within one proof, but we hope there is no danger of confusion. 2. Preparations. In the following transformations the inverse local times, defined below for $`k`$ and $`x`$, will play an essential rôle: $`T_U(k,x)`$ $`:=inf\{t1:U(t,x)=k\},`$ $`2.1`$$`2.2`$ $`T_D(k,x)`$ $`:=inf\{t1:D(t,x)=k\}.`$ It turns out that questions related to the local time are easier to handle if the random walk is observed at the random stopping times $`T_U(k,x)`$ and $`T_D(k,x)`$ rather than at deterministic times $`t1`$. This ‘combinatorial trick’ has its origin in Knight (1963) and has been successfully applied in various contexts. See, e.g., König (1996), Tóth (1995), (1996) and references cited there. We express $`f(4)`$ with the help of some straightforward rearrangements of summations: $$f(4)=\underset{x}{}\left(u(x)+d(x)\right)$$ $`2.3`$ where $`u(x):=`$ $`{\displaystyle \underset{t=1}{\overset{\mathrm{}}{}}}11_{\{S(t)=x,S(t1)=x1,x𝒦(t),\mathrm{\#}𝒦(t)=4\}}`$ $`2.4`$ $`=`$ $`{\displaystyle \underset{t=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}11_{\{T_U(k,x)=t,x𝒦(t),\mathrm{\#}𝒦(t)=4\}}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}11_{\{x𝒦(T_U(k,x)),\mathrm{\#}𝒦(T_U(k,x))=4\}}`$ and $`d(x):=`$ $`{\displaystyle \underset{t=1}{\overset{\mathrm{}}{}}}11_{\{S(t)=x,S(t1)=x+1,x𝒦(t),\mathrm{\#}𝒦(t)=4\}}`$ $`2.5`$ $`=`$ $`{\displaystyle \underset{t=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}11_{\{T_D(k,x)=t,x𝒦(t),\mathrm{\#}𝒦(t)=4\}}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}11_{\{x𝒦(T_D(k,x)),\mathrm{\#}𝒦(T_D(k,x))=4\}}.`$ Clearly, $$u(x)\stackrel{\text{law}}{=}d(x)$$ $`2.6`$ and, consequently $$𝔼\left(f(4)\right)=2\underset{x=1}{\overset{\mathrm{}}{}}𝔼\left(u(x)\right)+2\underset{x=0}{\overset{\mathrm{}}{}}𝔼\left(d(x)\right)$$ $`2.7`$ with $`𝔼\left(u(x)\right)`$ $`={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left([x𝒦(T_U(k,x))][\mathrm{\#}𝒦(T_U(k,x))=4]\right)`$ $`2.8`$$`2.9`$ $`𝔼\left(d(x)\right)`$ $`={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left([x𝒦(T_D(k,x))][\mathrm{\#}𝒦(T_D(k,x))=4]\right)`$ We shall show in details that $$\underset{x=1}{\overset{\mathrm{}}{}}𝔼\left(u(x)\right)=\underset{x=1}{\overset{\mathrm{}}{}}\underset{k=1}{\overset{\mathrm{}}{}}\left([x𝒦(T_U(k,x))][\mathrm{\#}𝒦(T_U(k,x))=4]\right)<\mathrm{}.$$ $`2.10`$ The similar statement $$\underset{x=0}{\overset{\mathrm{}}{}}𝔼\left(d(x)\right)=\underset{x=0}{\overset{\mathrm{}}{}}\underset{k=1}{\overset{\mathrm{}}{}}\left([x𝒦(T_D(k,x))][\mathrm{\#}𝒦(T_D(k,x))=4]\right)<\mathrm{}$$ $`2.11`$ can be proved in an identical way. 3. Ray-Knight Representation. Throughout this paper we denote by $`Y_t`$ a critical branching process with geometric offspring distribution (Galton-Watson process) and by $`Z_t`$ a critical branching process with geometric offspring distribution and one intruder at each generation. $`Y_t`$ and $`Z_t`$ are Markov chains with state space $`_+`$ and transition probabilities: $`\left(Y_{t+1}=j\|Y_t=i\right)=\pi (i,j):=`$ $`\{\begin{array}{ccc}\delta _{0,j}& \text{ if }& i=0,\\ \\ 2^{ij}{\displaystyle \frac{(i+j1)!}{(i1)!j!}}& \text{ if }& i>0.\end{array}`$ $`3.1`$$`3.2`$ $`\left(Z_{t+1}=j\|Z_t=i\right)=\rho (i,j):=`$ $`\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}^{ij1}{\displaystyle \frac{(i+j)!}{i!j!}}.`$ Let $`k0`$ and $`x1`$ be fixed integers and define the following three processes: $``$$`Z_t`$, $`0tx1`$, is a Markov chain with transition probabilities $`\rho (i,j)`$ and initial state $`Z_0=k`$; $``$$`Y_t`$, $`1t<\mathrm{}`$, is a Markov chain with transition probabilities $`\pi (i,j)`$ and initial state $`Y_1=k`$; $``$Finally, $`Y_t^{}`$, $`0t<\mathrm{}`$, is another Markov chain with the same transition probabilities $`\pi (i,j)`$ and initial state $`Y_0^{}=Z_{x1}`$. The three chains are independent, except for the fact that $`Y^{}`$ starts from the terminal state of $`Z`$. Using these three chains we patch together the process $$\mathrm{\Delta }_{x,k}(y):=\{\begin{array}{cccc}Z_{xy1}& \text{ if }& 0y& x1\\ Y_{yx}& \text{ if }& x1y& \mathrm{}\\ Y_y^{}& \text{ if }& \mathrm{}y& 0.\end{array}$$ $`3.3`$ We also define $$\mathrm{\Lambda }_{x,k}(y):=\mathrm{\Delta }_{x,k}(y)+\mathrm{\Delta }_{x,k}(y1)+11_{\{0<yx\}}.$$ $`3.4`$ According to the, by now classical, Ray-Knight Theorems on the local time process of simple symmetric random walks on $``$ (cf. Knight (1963), Ray (1963)), for any integers $`x1`$ and $`k0`$: $$\left(\mathrm{\Delta }_{x,k}(y),y\right)\stackrel{\text{law}}{=}\left(D(T_U(k+1,x),y),y\right).$$ $`3.5`$ Using (1.6) and (3.4), from this we get $$\left(\mathrm{\Lambda }_{x,k}(y),y\right)\stackrel{\text{law}}{=}\left(L(T_U(k+1,x),y),y\right).$$ $`3.6`$ 4. Proof of Theorem 1. Given the Markov chains $`Y_t`$, $`Z_t`$ and $`Y_t^{}`$ we define $$\stackrel{~}{Z}_t:=Z_t+Z_{t1}+1,\stackrel{~}{Y}_t:=Y_t+Y_{t1},\stackrel{~}{Y}_t^{}:=Y_t^{}+Y_{t1}^{}$$ $`4.1`$ and for $`h`$ the following stopping times $`\sigma _h`$ $`:=inf\{t0:Y_th\},`$ $`4.2`$$`4.3`$$`4.4`$$`4.5`$$`4.6`$$`4.7`$$`4.8`$$`4.9`$ $`\sigma _h^{}`$ $`:=inf\{t0:Y_t^{}h\},`$ $`\omega `$ $`:=inf\{t0:Y_t=0\},`$ $`\omega ^{}`$ $`:=inf\{t0:Y_t^{}=0\},`$ $`\tau _h`$ $`:=inf\{t0:Z_th\},`$ $`\stackrel{~}{\sigma }_{h,0}`$ $`:=0,\stackrel{~}{\sigma }_{h,i+1}:=inf\{t>\stackrel{~}{\sigma }_{h,i}:\stackrel{~}{Y}_th\},\stackrel{~}{\sigma }_h:=\stackrel{~}{\sigma }_{h,1},`$ $`\stackrel{~}{\sigma }_{h,0}^{}`$ $`:=0,\stackrel{~}{\sigma }_{h,i+1}^{}:=inf\{t>\stackrel{~}{\sigma }_{h,i}^{}:\stackrel{~}{Y}_t^{}h\},\stackrel{~}{\sigma }_h^{}:=\stackrel{~}{\sigma }_{h,1}^{},`$ $`\stackrel{~}{\tau }_{h,0}`$ $`:=0,\stackrel{~}{\tau }_{h,i+1}:=inf\{t>\stackrel{~}{\tau }_{h,i}:\stackrel{~}{Z}_th\},\stackrel{~}{\tau }_h:=\stackrel{~}{\tau }_{h,1}.`$ For $`h1`$, $`p0`$ and $`x1`$ fixed integers we define the following events: $`A_{h,0}:=`$ $`\left\{\mathrm{max}\{\stackrel{~}{Y}_t:1t<\mathrm{}\}<h\right\}`$ $`4.10`$$`4.11`$$`4.12`$$`4.13`$ $`=`$ $`\left\{\stackrel{~}{\sigma }_h=\mathrm{}\right\},`$ $`A_{h,p}:=`$ $`\left\{[\mathrm{max}\{\stackrel{~}{Y}_t:1t<\mathrm{}\}=h][\mathrm{\#}\{1t<\mathrm{}:\stackrel{~}{Y}_t=h\}=p]\right\}`$ $`=`$ $`\left\{[\stackrel{~}{\sigma }_{h,p}<\mathrm{}=\stackrel{~}{\sigma }_{h,p+1}][\stackrel{~}{Y}_{\stackrel{~}{\sigma }_{h,i}}=h,i=1,\mathrm{},p]\right\},`$ $`A_{h,0}^{}:=`$ $`\left\{\mathrm{max}\{\stackrel{~}{Y}_t^{}:1t<\mathrm{}\}<h\right\}`$ $`=`$ $`\left\{\stackrel{~}{\sigma }_h^{}=\mathrm{}\right\},`$ $`A_{h,p}^{}:=`$ $`\left\{[\mathrm{max}\{\stackrel{~}{Y}_t^{}:1t<\mathrm{}\}=h][\mathrm{\#}\{1t<\mathrm{}:\stackrel{~}{Y}_t^{}=h\}=p]\right\}`$ $`=`$ $`\left\{[\stackrel{~}{\sigma }_{h,p}^{}<\mathrm{}=\stackrel{~}{\sigma }_{h,p+1}^{}][\stackrel{~}{Y}_{\stackrel{~}{\sigma }_{h,i}^{}}^{}=h,i=1,\mathrm{},p]\right\},`$ $`B_{x,h,0}:=`$ $`\left\{\mathrm{max}\{\stackrel{~}{Z}_t:1t<x\}<h\right\}`$ $`4.14`$$`4.15`$ $`=`$ $`\left\{\stackrel{~}{\tau }_hx\right\},`$ $`B_{x,h,p}:=`$ $`\left\{[\mathrm{max}\{\stackrel{~}{Z}_t:1t<x\}=h][\mathrm{\#}\{1t<x:\stackrel{~}{Z}_t=h\}=p]\right\}`$ $`=`$ $`\left\{[\stackrel{~}{\tau }_{h,p}<x\stackrel{~}{\tau }_{h,p+1}][\stackrel{~}{Z}_{\stackrel{~}{\tau }_{h,i}}=h,i=1,\mathrm{},p]\right\}.`$ With the help of the Ray-Knight representation and the events introduced above we get the expression: $`𝔼\left(u(x)\right)={\displaystyle \underset{p+q+r=3}{}}{\displaystyle \underset{h=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}`$ $`\left(A_{h,p}\right\|Y_0=hk1)\times `$ $`4.16`$ $`\pi (k,hk1)\times `$ $`(B_{x,h,q}[Z_{x1}=l]\|Z_0=k)\times `$ $`\left(A_{h,r}^{}\|Y_0^{}=l\right)`$ which leads directly to $`{\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}𝔼\left(u(x)\right){\displaystyle \underset{p+q+r=3}{}}{\displaystyle \underset{h=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}`$ $`\left(A_{h,p}\right\|Y_0=hk1)\times `$ $`4.17`$ $`\pi (k,hk1)\times `$ $`\left({\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}\left(B_{x,h,q}\right\|Z_0=k)\right)\times `$ $`\left(\underset{l0}{sup}\left(A_{h,r}^{}\|Y_0^{}=l\right)\right).`$ The proof of Theorem 1 will follow directly from the bounds provided by ###### Proposition 1 For any $`\epsilon >0`$ there exists a finite constant $`C<\mathrm{}`$ such that for any $`h1`$ and $`k0`$: ($`ı`$) Without any restriction on $`k`$ or $`p`$ $$\underset{x=1}{\overset{\mathrm{}}{}}\left(B_{x,h,p}\|Z_0=k\right)Ch$$ $`4.18`$ ($`ıı`$) if either $`k[(hh^{1/2+\epsilon })/2,(hh^{1/2+\epsilon })/2]`$ or $`p1`$ holds then $`\left(A_{h,p}\|Y_0=k\right)`$ $`\left(Ch^{1/2+\epsilon }\right)^{p+1}`$ $`4.19`$$`4.20`$ $`{\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}\left(B_{x,h,p}\|Z_0=k\right)`$ $`\left(Ch^{1/2+\epsilon }\right)^{p+1}h`$ Remark: For $`kh`$ the left hand side of (4.18), (4.19) and (4.20), of course, vanish. We postpone the proof of this Proposition to the end of the next section and proceed with the proof of Theorem 1. Using the bounds (4.18)-(4.20) we prove (2.10). As we already mentioned, (2.11) is proved in a completely identical way. Theorem 1 follows from (2.10), (2.11) via (2.7). In the forthcoming argument $`\mathrm{}`$ will stand as abreviation of the summand on the right hand side of (4.17). On the right hand side of (4.17) keep $`p,q,r`$ and $`h`$ fixed and decompose the sum over $`k0`$ as follows: $$\underset{k}{}\mathrm{}=\underset{k:\|h2k\|h^{1/2+\epsilon }}{}\mathrm{}+\underset{k:\|h2k\|>h^{1/2+\epsilon }}{}\mathrm{}.$$ $`4.21`$ Similar decompositions will be applied a few more times throughout the paper. ###### Side-lemma 1 For any $`\epsilon >0`$ there exist constants $`C<\mathrm{}`$ and $`\gamma >0`$ such that for any $`h1`$ $$\underset{k:\|h2k\|>h^{1/2+\epsilon }}{}\pi (k,h1k)<C\mathrm{exp}(\gamma h^{2\epsilon }).$$ $`4.22`$ Side-lemmas are proved in Section 6. Using (4.18)-(4.21) we bound the sum over $`k`$ on the right hand side of (4.17) as follows: If $`r=0`$ and $`p+q=3`$ $`{\displaystyle \underset{k}{}}\mathrm{}`$ $`\left(Ch^{1/2+\epsilon }\right)^{p+q+2}h+\left(Ch\right)\left(C\mathrm{exp}(\gamma h^{2\epsilon })\right)`$ $`4.23`$ $``$ $`C^{}h^{3/2+5\epsilon }`$ with some properly chosen $`C^{}<\mathrm{}`$. If $`r>0`$ and $`p+q+r=3`$ $`{\displaystyle \underset{k}{}}\mathrm{}`$ $`\left(Ch^{1/2+\epsilon }\right)^{p+q+r+3}h+\left(Ch\right)\left(C\mathrm{exp}(\gamma h^{2\epsilon })\right)`$ $`4.24`$ $``$ $`C^{}h^{2+6\epsilon }`$ with some properly chosen $`C^{}<\mathrm{}`$. In both cases the upper bound is summable over $`h1`$, if we choose $`\epsilon <1/10`$. Hence (2.10) and the statement of Theorem 1. ∎(Theorem 1) 5. Technical Lemmas. The present section is divided in five subsections. In subsections 5.1-5.4 we state and prove some lemmas of more technical nature, needed in the proof of Proposition 1, which is presented in subsection 5.5. Throughout this section $`\epsilon >0`$ is fixed. 5.1. The Maximal Jump. We prove that the largest jump of the Markov chains $`Y_t`$ and $`Z_t`$, before reaching level $`h`$, is less than $`h^{1/2+\epsilon }`$, with overwhelming probability. Define the maximal jumps of $`Y_t`$, respectively, $`Z_t`$ as follows: $`M_h:=`$ $`sup\{\|Y_tY_{t1}\|:1t\sigma _h\}`$ $`5.1`$$`5.2`$ $`=`$ $`sup\{\|Y_tY_{t1}\|:1t\sigma _h\omega \},`$ $`N_h:=`$ $`sup\{\|Z_tZ_{t1}\|:1t\tau _h\}.`$ By definition $`M_h=0`$ if $`Y_0h`$ and $`N_h=0`$ if $`Z_0h`$. ###### Lemma 1 There exist two constants, $`C<\mathrm{}`$ and $`\gamma >0`$, such that for any $`h1`$ and $`k0`$ the following bounds hold: $`\left(M_h>h^{1/2+\epsilon }\|Y_0=k\right)<C\mathrm{exp}(\gamma h^{2\epsilon }),`$ $`5.3`$$`5.4`$ $`\left(N_h>h^{1/2+\epsilon }\|Z_0=k\right)<C\mathrm{exp}(\gamma h^{2\epsilon }).`$ ###### Demonstration Proof We prove here (5.3) in details. The proof of (5.4) is essentially similar and it is left for the reader. For the moment let $`\gamma `$ be an arbitrary positive number. Its value will be fixed at the end of this proof. $`(M_h>h^{1/2+\epsilon }`$ $`\|Y_0=k)`$ $`5.5`$ $``$ $`\left([M_h>h^{1/2+\epsilon }][\sigma _h\omega h^2\mathrm{exp}(\gamma h^{2\epsilon })]\|Y_0=k\right)`$ $`+`$ $`\left(\sigma _h\omega >h^2\mathrm{exp}(\gamma h^{2\epsilon })\|Y_0=k\right)`$ ###### Side-lemma 2 There exists a constant $`C<\mathrm{}`$, such that for any $`h1`$ and $`k0`$ $$𝔼\left(\sigma _h\omega \|Y_0=k\right)Ch^2.$$ $`5.6`$ By using Markov’s inequality we get the following upper bound on the second term of the right hand side of (5.5): $$\left(\sigma _h\omega >h^2\mathrm{exp}(\gamma h^{2\epsilon })\|Y_0=k\right)C\mathrm{exp}(\gamma h^{2\epsilon }).$$ $`5.7`$ To bound the first term on the right hand side of (5.5) we use the following representation of the Markov chain $`Y_t`$: Let $`\left(\xi _{t,i}\right)_{t1,i1}`$ be i.i.d random variables with common geometric distribution $`\left(\xi _{t,i}=k\right)=2^{k1}`$ The process $`Y_t`$ is realized as follows: fix $`Y_0`$ and put $$Y_{t+1}=\underset{j=1}{\overset{Y_t}{}}\xi _{t+1,j}.$$ $`5.8`$ Using this representation we note that $`\left([M_h>h^{1/2+\epsilon }][\sigma _h\omega h^2\mathrm{exp}(\gamma h^{2\epsilon })]\|Y_0=k\right)`$ $`5.9`$ $`\left(\mathrm{max}\left\{\underset{1jh}{\mathrm{max}}\left\|{\displaystyle \underset{i=1}{\overset{j}{}}}\left(\xi _{t,i}1\right)\right\|:1th^2\mathrm{exp}(\gamma h^{2\epsilon })\right\}>h^{1/2+\epsilon }\right)`$ $`=1\left(1\left(\underset{1jh}{\mathrm{max}}\left\|{\displaystyle \underset{i=1}{\overset{j}{}}}\left(\xi _{1,i}1\right)\right\|>h^{1/2+\epsilon }\right)\right)^{h^2\mathrm{exp}(\gamma h^{2\epsilon })}`$ ###### Side-lemma 3 Let $`\xi _i`$ be i.i.d. random variables with the common geometric disatribution $`\left(\xi _i=k\right)=2^{k1}`$. Then there is a constant $`\theta _0>0`$ such that for any $`\lambda >0`$ and $`n`$ satisfying $`\lambda /(4n)<\theta _0`$: $$\left(\underset{1jn}{\mathrm{max}}\left\|\underset{i=1}{\overset{j}{}}\left(\xi _i1\right)\right\|>\lambda \right)2\mathrm{exp}(\lambda ^2/(8n)).$$ $`5.10`$ Using this bound we get $`\left([M_h>h^{1/2+\epsilon }][\sigma _h\omega h^2\mathrm{exp}(\gamma h^{2\epsilon })]\|Y_0=k\right)`$ $`5.11`$ $`1\left(12\mathrm{exp}(h^{2\epsilon }/8)\right)^{h^2\mathrm{exp}(\gamma h^{2\epsilon })}`$ $`2h^2\mathrm{exp}\left((\gamma 8^1)h^{2\epsilon }\right).`$ In the last inequality we use the fact that for $`0<\alpha <1<\beta `$, $`1\alpha \beta <(1\alpha )^\beta `$. We choose $`\gamma <16^1`$. From (5.5), (5.7) and (5.11) we get (5.3), with an appropriately chosen constant $`C<\mathrm{}`$. ∎(Lemma 1) 5.2. Hitting exactly $`h`$. ###### Lemma 2 There exists a constant $`C<\mathrm{}`$ such that for any $`h1`$ and $`k0`$ $`\left([\stackrel{~}{\sigma }_h<\mathrm{}][\stackrel{~}{Y}_{\stackrel{~}{\sigma }_h}=h]\|Y_0=k\right)<Ch^{1/2+\epsilon }`$ $`5.12`$$`5.13`$ $`\left(\stackrel{~}{Z}_{\stackrel{~}{\tau }_h}=h\|Z_0=k\right)<Ch^{1/2+\epsilon }`$ ###### Demonstration Proof Again, we give the details of the proof of (5.12), leaving the identical details of (5.13) for the reader. $`\left([\stackrel{~}{\sigma }_h<\mathrm{}][\stackrel{~}{Y}_{\stackrel{~}{\sigma }_h}=h]\|Y_0=k\right)=`$ $`5.14`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left([\stackrel{~}{\sigma }_h<\mathrm{}][Y_{\stackrel{~}{\sigma }_h1}=l][Y_{\stackrel{~}{\sigma }_h}=hl]\|Y_0=k\right)`$ We divide the sum in two parts, as in (4.21): $`{\displaystyle \underset{l:\|h2l\|>h^{1/2+\epsilon }}{}}`$ $`\left([\stackrel{~}{\sigma }_h<\mathrm{}][Y_{\stackrel{~}{\sigma }_h1}=l][Y_{\stackrel{~}{\sigma }_h}=hl]\|Y_0=k\right)`$ $`5.15`$ $`\left(M_h>h^{1/2+\epsilon }\|Y_0=k\right)<C\mathrm{exp}(\gamma h^{2\epsilon }),`$ by Lemma 1. On the other hand: $`{\displaystyle \underset{l:\|h2l\|h^{1/2+\epsilon }}{}}\left([\stackrel{~}{\sigma }_h<\mathrm{}][Y_{\stackrel{~}{\sigma }_h1}=l][Y_{\stackrel{~}{\sigma }_h}=hl]\|Y_0=k\right)=`$ $`5.16`$ $`{\displaystyle \underset{l:\|h2l\|h^{1/2+\epsilon }}{}}\left([\stackrel{~}{\sigma }_h<\mathrm{}][Y_{\stackrel{~}{\sigma }_h1}=l]\|Y_0=k\right){\displaystyle \frac{\pi (l,hl)}{_{mhl}\pi (l,m)}}`$ ###### Side-lemma 4 There exists a constant $`C<\mathrm{}`$, such that for any $`h1`$ and $`l[(hh^{1/2+\epsilon })/2,(h+h^{1/2+\epsilon })/2]`$ $$\frac{\pi (l,hl)}{_{mhl}\pi (l,m)}<Ch^{1/2+\epsilon }$$ $`5.17`$ From this we get $$\underset{l:\|h2l\|h^{1/2+\epsilon }}{}\left([\stackrel{~}{\sigma }_h<\mathrm{}][Y_{\stackrel{~}{\sigma }_h1}=l][Y_{\stackrel{~}{\sigma }_h}=hl]\|Y_0=k\right)Ch^{1/2+\epsilon }.$$ $`5.18`$ Finally, (5.15) and (5.18) yield (5.13). ∎(Lemma 2) 5.3. $`\stackrel{~}{Y}_t`$ does not hit level $`h`$. ###### Lemma 3 There exists a constant $`C<\mathrm{}`$ such that for any $`h1`$ and $`k[(hh^{1/2+\epsilon })/2,(h+h^{1/2+\epsilon })/2]`$ $$\left(\stackrel{~}{\sigma }_h=\mathrm{}\|Y_0=k\right)<Ch^{1/2+\epsilon }.$$ $`5.19`$ ###### Demonstration Proof $`\left(\stackrel{~}{\sigma }_h=\mathrm{}\|Y_0=k\right)`$ $`\left([\stackrel{~}{\sigma }_h=\mathrm{}][M_hh^{1/2+\epsilon }]\|Y_0=k\right)`$ $`5.20`$ $`+`$ $`\left(M_h>h^{1/2+\epsilon }\|Y_0=k\right).`$ To bound the first term on the right hand side, note that $$\left\{[\stackrel{~}{\sigma }_h=\mathrm{}][M_hh^{1/2+\epsilon }]\right\}\left\{\sigma _{(h+h^{1/2+\epsilon })/2}=\mathrm{}\right\}.$$ $`5.21`$ ###### Side-lemma 5 There exists a constant $`C<\mathrm{}`$ such that for any $`0k<h`$ $$\left(\sigma _h=\mathrm{}\|Y_0=k\right)<\frac{hk}{h}+Ch^{1/2}.$$ $`5.22`$ Thus $`\left([\stackrel{~}{\sigma }_h=\mathrm{}][M_hh^{1/2+\epsilon }]\|Y_0=k\right)`$ $`\left(\sigma _{(h+h^{1/2+\epsilon })/2}=\mathrm{}\|Y_0=k\right)`$ $`5.23`$ $``$ $`Ch^{1/2+\epsilon },`$ for $`k[(hh^{1/2+\epsilon })/2,(h+h^{1/2+\epsilon })/2]`$. This bound, together with (5.20) and (5.3) yield (5.19). ∎(Lemma 3) 5.4. Expectation of $`\stackrel{~}{\tau }_h`$. ###### Lemma 4 There exists a constant $`C<\mathrm{}`$ such that for any $`h1`$ the following bounds hold: ($`ı`$) Without any restriction on $`k`$: $$𝔼\left(\stackrel{~}{\tau }_h\|Z_0=k\right)<Ch.$$ $`5.24`$ ($`ıı`$) For $`k[(hh^{1/2+\epsilon })/2,(h+h^{1/2+\epsilon })/2]`$ $$𝔼\left(\stackrel{~}{\tau }_h\|Z_0=k\right)<Ch^{1/2+\epsilon }.$$ $`5.25`$ ###### Demonstration Proof $$𝔼\left(\stackrel{~}{\tau }_h\|Z_0=k\right)=𝔼\left(\stackrel{~}{\tau }_h11_{\{N_hh^{1/2+\epsilon }\}}\|Z_0=k\right)+𝔼\left(\stackrel{~}{\tau }_h11_{\{N_h>h^{1/2+\epsilon }\}}\|Z_0=k\right).$$ $`5.26`$ We bound the first, respectively, the second term on the right hand side, by noting $$\stackrel{~}{\tau }_h11_{\{N_hh^{1/2+\epsilon }\}}\tau _{(h+h^{1/2+\epsilon })/2},$$ $`5.27`$ respectively, $$\stackrel{~}{\tau }_h\tau _h.$$ $`5.28`$ Thus we get $`𝔼\left(\stackrel{~}{\tau }_h\|Z_0=k\right)`$ $`𝔼\left(\tau _{(h+h^{1/2+\epsilon })/2}\|Z_0=k\right)`$ $`5.29`$ $`+𝔼\left(\tau _h^2\|Z_0=k\right)^{1/2}\left(N_h>h^{1/2+\epsilon }\|Z_0=k\right)^{1/2}.`$ ###### Side-lemma 6 There exists a constant $`C<\mathrm{}`$ such that for any $`0k<h`$ $$𝔼\left(\tau _h\|Z_0=k\right)<(hk)+Ch^{1/2}.$$ $`5.30`$ ###### Side-lemma 7 There exists a constant $`C<\mathrm{}`$ such that for any $`0k<h`$ $$𝔼\left(\tau _h^2\|Z_0=k\right)<Ch^2.$$ $`5.31`$ Putting together (5.29), (5.30) and (5.31), we get (5.24) and (5.25). ∎(Lemma 4) 5.5 Proof of Proposition 1. First note that $`\left(A_{h,0}\|Y_0=k\right)`$ $`=\left(\stackrel{~}{\sigma }_h=\mathrm{}\|Y_0=k\right),`$ $`5.32`$$`5.33`$ $`{\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}\left(B_{x,h,0}\|Z_0=k\right)`$ $`={\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}\left(\stackrel{~}{\tau }_hx\|Z_0=k\right)=𝔼\left(\stackrel{~}{\tau }_h\|Z_0=k\right),`$ and for $`p1`$, using the strong Markov property of $`Y_t`$, respectively, $`Z_t`$: $`(`$ $`A_{h,p}\|Y_0=k)={\displaystyle }_{l=0}^{\mathrm{}}([\stackrel{~}{\sigma }_h<\mathrm{}][Y_{\stackrel{~}{\sigma }_h1}=hl][Y_{\stackrel{~}{\sigma }_h}=l]\|Y_0=k)\times `$ $`5.34`$$`5.35`$ $`\left(A_{h,p1}\|Y_0=l\right)`$ $`{\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}`$ $`\left(B_{x,h,p}\right\|Z_0=k)={\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}([Z_{\stackrel{~}{\tau }_h1}=hl][Z_{\stackrel{~}{\sigma }_h}=l]\|Z_0=k)\times `$ $`\left({\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}\left(B_{x,h,p1}\|Z_0=l\right)\right).`$ We prove the bounds of Proposition 1 by induction on $`p`$. According to (5.32), (5.33), for $`p=0`$, (4.18), (4.19) and (4.20) are just restatements of (5.24), (5.19) and (5.25), respectively. (See Lemma 3 and Lemma 4.) Next we consider the case $`p=1`$. Again, we divide the sum over $`l`$ in (5.34) and (5.35) in two parts, as it was done in (4.21). From (5.12) (Lemma 2) and (5.19) (Lemma 3) $`{\displaystyle \underset{l:\|h2l\|h^{1/2+\epsilon }}{}}`$ $`\left([\stackrel{~}{\sigma }_h<\mathrm{}][Y_{\stackrel{~}{\sigma }_h1}=hl][Y_{\stackrel{~}{\sigma }_h}=l]\|Y_0=k\right)\left(A_{h,0}\|Y_0=l\right)`$ $`5.36`$ $`\left(Ch^{1/2+\epsilon }\right)\left(Ch^{1/2+\epsilon }\right).`$ From (5.3) (Lemma 1) $`{\displaystyle \underset{l:\|h2l\|>h^{1/2+\epsilon }}{}}\left([\stackrel{~}{\sigma }_h<\mathrm{}][Y_{\stackrel{~}{\sigma }_h1}=hl][Y_{\stackrel{~}{\sigma }_h}=l]\|Y_0=k\right)\left(A_{h,0}\|Y_0=l\right)`$ $`5.37`$ $`\left(M_h>h^{1/2+\epsilon }\|Y_0=k\right)<C\mathrm{exp}(\gamma h^{2\epsilon })`$ From (5.36) and (5.37) we get (4.19) for $`p=1`$. Applying the same ideas to (5.35): from (5.13) (Lemma 2) and (5.25) (Lemma 4) $`{\displaystyle \underset{l:\|h2l\|h^{1/2+\epsilon }}{}}\left([Z_{\stackrel{~}{\tau }_h1}=hl][Z_{\stackrel{~}{\sigma }_h}=l]\|Z_0=k\right){\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}\left(B_{x,h,0}\|Z_0=l\right)`$ $`5.38`$ $`\left(Ch^{1/2+\epsilon }\right)\left(Ch^{1/2+\epsilon }\right)`$ From (5.4) (Lemma 1) and (5.24) (Lemma 4) $`{\displaystyle \underset{l:\|h2l\|>h^{1/2+\epsilon }}{}}\left([Z_{\stackrel{~}{\sigma }_h1}=hl][Z_{\stackrel{~}{\sigma }_h}=l]\|Z_0=k\right){\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}\left(B_{x,h,0}\|Z_0=l\right)`$ $`5.39`$ $`\left(N_h>h^{1/2+\epsilon }\|Z_0=k\right)\left(\underset{l0}{sup}{\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}\left(B_{x,h,0}\|Z_0=l\right)\right)`$ $`<\left(C\mathrm{exp}(\gamma h^{2\epsilon })\right)(Ch)`$ (5.38) and (5.39) yield (4.20) for $`p=1`$. For $`p2`$ the induction follows from the same reasonings, just one does not have to split the sum over $`l`$ as in (5.36), (5.37). After the previous arguments we may ignore these completely straightforward details. ∎(Proposition 1) 6. Proof of the Side-Lemmas. First we prove Side-lemmas 1 and 4. Then Side-lemma 3 follows, which relies on an exponential Kolmogorov inequality. These proofs are rather standard ‘classroom exercises’. Side-lemmas 2, 5, 6 and 7 follow from an estimate on the overshooting of level $`h`$ by the processes $`Y_t`$ and $`Z_t`$ stopped at $`\sigma _h\omega `$, respectively, $`\tau _h`$ and from optional stopping arguments. ###### Demonstration Proof of Side-lemma 1 Assume $`h2`$ and denote $`h2=:n`$, $`k1=:l`$. Then, using the explicit form (3.1) of $`\pi (i,j)`$, the right hand side of (4.22) becomes $$\underset{k:\|h2k\|>h^{1/2+\epsilon }}{}\pi (k,h1k)=\frac{1}{2}\left(\left\|2B_nn\right\|>(n+2)^{1/2+\epsilon }\right)$$ $`6.1`$ where $`B_n`$ is binomially distributed: $`\left(B_n=l\right)=\left(\genfrac{}{}{0pt}{}{n}{l}\right)2^n`$. Using the fact that for any $`\gamma <1/2`$ $$\underset{n}{sup}𝔼\left(\mathrm{exp}\left\{\gamma \left(2B_nn\right)^2/n\right\}\right)=C_\gamma <\mathrm{},$$ $`6.2`$ by Markov’s inequality we get (4.22). ∎(Side-lemma 1) ###### Demonstration Proof of Side-lemma 4 Note first that for $`i1`$ and $`j0`$ $$\frac{\pi (i,j+1)}{\pi (i,j)}=\frac{i+j}{2(1+j)}$$ $`6.3`$ From this it follows that the distribution $`j\pi (i,j)`$ is unimodular and for $`i2`$ fixed $`\pi (i,j)<\pi (i,j+1),\text{ for }0ji3,`$ $`6.4`$ $`\pi (i,i2)=\pi (i,i1),`$ $`\pi (i,j)>\pi (i,j+1)\text{ for }i1j<\mathrm{}.`$ We treat separately the cases $`l[h/2,(h+h^{1/2+\epsilon })/2]`$ and $`l[(hh^{1/2+\epsilon })/2,h/2]`$: For $`l[h/2,(h+h^{1/2+\epsilon })/2]`$ the following two facts imply (5.17) (1) By (6.4), $`\pi (l,hl)`$ $`\pi (l,l1)={\displaystyle \frac{1}{2}}\left({\displaystyle \genfrac{}{}{0pt}{}{2(l1)}{l1}}\right)2^{2(l1)}`$ $`6.5`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \genfrac{}{}{0pt}{}{2(h/21)}{h/21}}\right)2^{2(h/21)}Ch^{1/2}`$ (2) By the central limit theorem: $`lim_l\mathrm{}_{ml}\pi (l,m)=\frac{1}{2}`$ and thus there exists a constant $`c>0`$ such that for any $`0<h/2l`$ $$\underset{mhl}{}\pi (l,m)\underset{ml}{}\pi (l,m)c$$ $`6.6`$ Let now $`l[(h+h^{1/2\epsilon })/2,h/2]`$ and $`k:=hl+h^{1/2\epsilon }`$ Then: $`{\displaystyle \frac{\pi (l,hl)}{_{mhl}\pi (l,m)}}`$ $`(kh+l+1)^1{\displaystyle \frac{\pi (l,hl)}{\pi (l,k)}}`$ $`(kh+l+1)^1\left({\displaystyle \frac{\pi (l,k1)}{\pi (l,k)}}\right)^{kh+l}=\mathrm{}`$ $`\mathrm{}`$ $`=(kh+l+1)^1\left({\displaystyle \frac{2k}{l+k1}}\right)^{kh+l}`$ $`6.7`$ $`h^{1/2+\epsilon }\left({\displaystyle \frac{2(hl+h^{1/2\epsilon })}{h+h^{1/2\epsilon }1}}\right)^{h^{1/2\epsilon }}`$ $`h^{1/2+\epsilon }\left({\displaystyle \frac{h+h^{1/2+\epsilon }+2h^{1/2\epsilon })}{h+h^{1/2\epsilon }1}}\right)^{h^{1/2\epsilon }}`$ $`h^{1/2+\epsilon }\left(1+3h^{(1/2\epsilon )}\right)^{h^{1/2\epsilon }}`$ $`e^3h^{1/2+\epsilon }`$ In the first inequality we used (6.1). In the second one we exploited the fact that, according to (6.2), for any $`i1`$ fixed $`\pi (i,j)/\pi (i,j+1)`$ is an increasing function of $`j0`$. In the next equality, (6.2) was explicitly used. In the third inequality we inserted the value of $`k=hl+h^{1/2\epsilon }`$. In the fourth inequality $`l=(hh^{1/2\epsilon })/2`$ was inserted to maximize the expression. In the last but one inequality we used $`1h^{1/2\epsilon }h^{1/2+\epsilon }`$. Finally, in the last inequality we use the fact that $`sup_{\alpha 1}\left(1+3\alpha ^1\right)^\alpha e^3`$ ∎(Side-lemma 4) ###### Demonstration Proof of Side-lemma 3 The Exponential Kolmogorov Inequality follows directly from Doob’s maximal inequality. For its proof see e.g. page 139 of Williams (1991). ###### Exponential Kolmogorov Inequality Let $`\stackrel{~}{\xi }_j`$, $`j1`$, be i.i.d. random variables with $`𝔼\left(\mathrm{exp}\left\{\theta \left\|\stackrel{~}{\xi }_j\right\|\right\}\right)<\mathrm{}`$ for some $`\theta >0`$ and $`𝔼\left(\stackrel{~}{\xi }_j\right)=0`$. Then for any $`\lambda (0,\mathrm{})`$ and $`n`$ $$\left(\underset{1jn}{\mathrm{max}}\left\|\underset{i=1}{\overset{j}{}}\stackrel{~}{\xi }_i\right\|>\lambda \right)e^{\lambda \theta }\left\{𝔼\left(e^{\theta \stackrel{~}{\xi }_i}\right)^n+𝔼\left(e^{\theta \stackrel{~}{\xi }_i}\right)^n\right\}$$ $`6.8`$ We apply the Exponential Kolmogorov Inequality to $`\stackrel{~}{\xi }_i=\xi _i1`$, with $`\left(\xi _i=k\right)=2^{k1}`$, $`k0`$. There exists a constant $`\theta _0>0`$ such that for $`0\theta <\theta _0`$ we get: $`𝔼\left(e^{\theta (\xi _j1)}\right)=e^\theta \left(2e^\theta \right)^1=1+\theta ^2+𝒪\left(\theta ^3\right)<e^{2\theta ^2}`$ $`6.9`$$`6.10`$ $`𝔼\left(e^{\theta (\xi _j1)}\right)=e^{2\theta }\left(2e^\theta 1\right)^1=1+\theta ^2+𝒪\left(\theta ^3\right)<e^{2\theta ^2}`$ Inserting these bounds into the right hand side of (6.8) and choosing $`\theta =\lambda /(4n)`$ we obtain (5.10). ∎(Side-lemma 3) The proof of Side-lemmas 2, 5, 6 and 7 will follow from the forthcoming Overshooting Lemma and standard optional stopping considerations. The Overshooting Lemma and its Corollary are extended restatements of Lemmas 3.2 and 3.4 from Tóth and Werner (1997). ###### Overshooting Lemma For any $`0k<hu`$ the following overshoot bounds hold: $`\left(Y_{\sigma _h}u\|[Y_0=k][\sigma _h<\mathrm{}]\right)`$ $`\left(Y_1u\|[Y_0=h][Y_1h]\right)`$ $`6.11`$$`6.12`$ $`=`$ $`{\displaystyle \frac{_{v=u}^{\mathrm{}}\pi (h,v)}{_{w=h}^{\mathrm{}}\pi (h,w)}},`$ $`\left(Z_{\tau _h}u\|Z_0=k\right)`$ $`\left(Z_1u\|[Z_0=h][Z_1h]\right)`$ $`=`$ $`{\displaystyle \frac{_{v=u}^{\mathrm{}}\rho (h,v)}{_{w=h}^{\mathrm{}}\rho (h,w)}}.`$ In particular it follows that ###### Corollary There exists a constant $`C<\mathrm{}`$ such that for any $`0k<h`$: $`𝔼\left(Y_{\sigma _h}\|[Y_0=k][\sigma _h<\mathrm{}]\right)`$ $`{\displaystyle \frac{_{v=h}^{\mathrm{}}\pi (h,v)v}{_{w=h}^{\mathrm{}}\pi (h,w)}}h+Ch^{1/2}`$ $`6.13`$$`6.14`$$`6.15`$$`6.16`$ $`𝔼\left(Y_{\sigma _h}^2\|[Y_0=k][\sigma _h<\mathrm{}]\right)`$ $`{\displaystyle \frac{_{v=h}^{\mathrm{}}\pi (h,v)v^2}{_{w=h}^{\mathrm{}}\pi (h,w)}}h^2+Ch^{3/2}`$ $`𝔼\left(Z_{\tau _h}\|Z_0=k\right)`$ $`{\displaystyle \frac{_{v=h}^{\mathrm{}}\rho (h,v)v}{_{w=h}^{\mathrm{}}\rho (h,w)}}h+Ch^{1/2}`$ $`𝔼\left(Z_{\tau _h}^2\|Z_0=k\right)`$ $`{\displaystyle \frac{_{v=h}^{\mathrm{}}\rho (h,v)v^2}{_{w=h}^{\mathrm{}}\rho (h,w)}}h^2+Ch^{3/2}`$ The rightmost bounds in (6.13)-(6.16) follow from explicit computations. ###### Demonstration Proof of the Overshooting Lemma Straightforward manipulations leed to the following identities for $`1hv`$: $`\left(Y_{\sigma _h}=v\|[Y_0=k][\sigma _h<\mathrm{}]\right)=`$ $`6.17`$ $`{\displaystyle \underset{l=0}{\overset{h1}{}}}\left(Y_{\sigma _h1}=l\|[Y_0=k][\sigma _h<\mathrm{}]\right){\displaystyle \frac{\pi (l,v)}{_{w=h}^{\mathrm{}}\pi (l,w)}},`$ $$\left(Z_{\tau _h}=v\|Z_0=k\right)=\underset{l=0}{\overset{h1}{}}\left(Z_{\tau _h1}=l\|Z_0=k\right)\frac{\rho (l,v)}{_{w=h}^{\mathrm{}}\rho (l,w)}.$$ $`6.18`$ Using the explicit form (3.1), respectively, (3.2) of the transition probabilities $`\pi (i,j)`$, respectively, $`\rho (i,j)`$, it is easy to check the following inequalities, which hold for any $`0<l<hv`$, respectively, $`0l<hv`$: $`{\displaystyle \frac{\pi (l+1,v)}{\pi (l,v)}}={\displaystyle \frac{l+v}{2l}}`$ $`<{\displaystyle \frac{l+v+1}{2l}}={\displaystyle \frac{\pi (l+1,v+1)}{\pi (l,v+1)}},`$ $`6.19`$$`6.20`$ $`{\displaystyle \frac{\rho (l+1,v)}{\rho (l,v)}}={\displaystyle \frac{l+v+1}{2(l+1)}}`$ $`<{\displaystyle \frac{l+v+2}{2(l+1)}}={\displaystyle \frac{\rho (l+1,v+1)}{\rho (l,v+1)}}.`$ It follows that for any $`0l<hv<w`$ $`\pi (l+1,v)\pi (l,w)<\pi (l,v)\pi (l+1,w),`$ $`6.21`$$`6.22`$ $`\rho (l+1,v)\rho (l,w)<\rho (l,v)\rho (l+1,w).`$ Hence, for any $`0l<hu`$ $`{\displaystyle \underset{v=h}{\overset{\mathrm{}}{}}}\pi (l+1,v){\displaystyle \underset{w=u}{\overset{\mathrm{}}{}}}\pi (l,w)<{\displaystyle \underset{v=h}{\overset{\mathrm{}}{}}}\pi (l,v){\displaystyle \underset{w=u}{\overset{\mathrm{}}{}}}\pi (l+1,w),`$ $`6.23`$$`6.24`$ $`{\displaystyle \underset{v=h}{\overset{\mathrm{}}{}}}\rho (l+1,v){\displaystyle \underset{w=u}{\overset{\mathrm{}}{}}}\rho (l,w)<{\displaystyle \underset{v=h}{\overset{\mathrm{}}{}}}\rho (l,v){\displaystyle \underset{w=u}{\overset{\mathrm{}}{}}}\rho (l+1,w),`$ which directly imply (6.11), respectively, (6.12). ∎(Overshooting Lemma) ###### Demonstration Proof of Side-lemmas 2 and 5 We apply the Optional Stopping Theorem to the martingales $`Y_t`$ (for Side-lemma 5), respectively, $`Y_t^22_{s=0}^{t1}Y_s`$ (for Side-lemma 2), $`t0`$, both stopped at $`\sigma _h\omega `$. First we prove Side-lemma 5: $`k=𝔼\left(Y_{\sigma _h\omega }\|Y_0=k\right)=𝔼\left(Y_{\sigma _h}\|[Y_0=k][\sigma _h<\mathrm{}]\right)\left(\sigma _h<\mathrm{}\|Y_0=k\right)`$ $`6.25`$ $`\left(h+C\sqrt{h}\right)\left(\sigma _h<\mathrm{}\|Y_0=k\right).`$ Where, in the last inequality we applied (6.13). Hence (5.23). ∎(Side-lemma 5) Next we prove Side-lemma 2: $$k^2=𝔼\left(Y_{\sigma _h\omega }^22\underset{s=0}{\overset{\sigma _h\omega 1}{}}Y_s\|Y_0=k\right)𝔼\left(Y_{\sigma _h\omega }^2\|Y_0=k\right)2𝔼\left(\sigma _h\omega \|Y_0=k\right),$$ $`6.26`$ where in the last inequality we used the fact that $`Y_s1`$ for $`s<\omega `$. Hence $$2𝔼\left(\sigma _h\omega \|Y_0=k\right)𝔼\left(Y_{\sigma _h}^2\|[Y_0=k][\sigma _h<\mathrm{}]\right)\left(\sigma _h<\mathrm{}\|Y_0=k\right)k^2<Ch^2.$$ $`6.27`$ In the last inequality we used (6.14). ∎(Side-lemma 2) ###### Demonstration Proof of Side-lemmas 6 and 7 We apply the Optional Stopping Theorem to the martingale $`Z_tt`$ (for Side-lemma 6), respectively, to the supermartingale $`t^22tZ_t`$ (for Side-lemma 7), $`t0`$, both stopped at $`\tau _h`$. First Side-lemma 6: $$𝔼\left(\tau _h\|Z_0=k\right)=𝔼\left(Z_{\tau _h}\|Z_0=k\right)khk+C\sqrt{h}.$$ $`6.28`$ In the last inequality (6.15) was used. ∎(Side-lemma 6) Next, Side-lemma 7: $$𝔼\left(\tau _h^2\|Z_0=k\right)2𝔼\left(\tau _hZ_{\tau _h}\|Z_0=k\right)2\sqrt{𝔼\left(\tau _h^2\|Z_0=k\right)}\sqrt{𝔼\left(Z_{\tau _h}^2\|Z_0=k\right)}$$ $`6.29`$ Hence, using (6.16) we get $$𝔼\left(\tau _h^2\|Z_0=k\right)4𝔼\left(Z_{\tau _h}^2\|Z_0=k\right)Ch^2.$$ $`6.30`$ ∎(Side-lemma 7) Acknowledgement: This work was partially supported by the following grants: OTKA-T26176 (Hungarian National Foundation for Scientific Research), FKFP-0638/99 (Ministry of Education) and TKI-StochasticsTUB (Hungarian Academy of Sciences). References: Bass, R.F., Eisenbaum, N., Shi, Z. (2000): The most visited sites of symmetric stable process. Probability Theory and Related Fields (to appear) Bass, R.F., Griffin, P.S. (1985): The most visited site of Brownian motion and random walk, Z. Wahrscheinlichkeitstheorie verw. 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warning/0004/cond-mat0004253.html	ar5iv	text	# Magnetic Field Effects on Neutron Diffraction in the Antiferromagnetic Phase of 𝑈⁢𝑃⁢𝑡₃ ## Abstract We discuss possible magnetic structures in UPt<sub>3</sub> based on our analysis of elastic neutron-scattering experiments in high magnetic fields at temperatures $`T<T_N`$. The existing experimental data can be explained by a single-q antiferromagnetic structure with three independent domains. For modest in-plane spin-orbit interactions, the Zeeman coupling between the antiferromagnetic order parameter and the magnetic field induces a rotation of the magnetic moments, but not an adjustment of the propagation vector of the magnetic order. A triple-q magnetic structure is also consistent with neutron experiments, but in general leads to a non-uniform magnetization in the crystal. New experiments could decide between these structures. The coexistence of antiferromagnetic and superconducting order for five of the six heavy fermion superconductors suggests a deep connection between these two aspects of heavy fermion physics. In these materials the f-electrons are involved in the superconducting transition, just as they are in the formation of the coherent heavy fermion band, but their precise role in the development of the unconventional superconducting phase is still unclear. The magnetic field versus temperature phase diagram of $`UPt_3`$ provided compelling evidence of unconventional superconductivity in U-based heavy fermion materials. In order to explain the phase diagram of UPt<sub>3</sub> several authors proposed a multicomponent order parameter based on a multi-dimensional representation of the hexagonal point group. In these models a weak symmetry breaking field (SBF) is invoked. This SBF lifts the degeneracy of the multi-dimensional representation and leads to multiple transitions at lower temperatures and higher fields (see also the reviews in Ref. ). A natural candidate for the role of SBF is the weak antiferromagnetic order shown by neutron scattering measurements below $`T_N=6K`$. The ordered moment is unusually small, only $`0.02\mu _B`$ per U atom, and is directed in the basal plane, thus breaking the in-plane hexagonal symmetry. Evidence in support of an antiferromagnetic SBF coupled to the the superconducting order parameter is based on the correlation between changes in the magnitude of the ordered moment and the splitting of the double transition. Both the splitting and the AFM order parameter are suppressed under applied pressure of $`p_c3.5\text{kbar}`$. The effect of $`Pd`$ is the opposite; the splitting and the ordered moment increase with increasing $`Pd`$ substitution. Most thermodynamic and transport measurements have failed to detect a signature of AFM ordering near $`T_N6\text{K}`$. However, evidence of magnetic ordering is observed to onset at $`T_N`$ in the magnetoresistance. The transition has other unusual characteristics as well, including finite range correlations, $`\xi _{AFM}300500`$ Å, depending on the crystalline direction and sample. By contrast , $`(U,Th)(Pd,Pt)_3`$ alloys exhibit AFM ordering at $`T_N6\text{K}`$, but with ordered moments of conventional size, $`\mu 0.65\mu _B/U`$-ion, and resolution-limited Bragg peaks at the same positions as pure $`UPt_3`$. Based on these facts, several authors have argued that the anomaly at $`6K`$ does not indicate the onset of true long range magnetic ordering but finite-range AFM correlations, which may also be fluctuating on time scales of order $`510^{10}`$ s to $`10^7`$ s. Given the uncertainties about the nature of magnetic order in UPt<sub>3</sub>, studies of the field dependence of the magnetic order were performed in order to help clarify these issues. Two experimental groups have measured neutron scattering ratios in magnetic fields up to $`3.5\text{T}`$ and $`12T`$. Both studies concluded that applied magnetic fields have no effect on the magnetic order of $`UPt_3`$, whether it be in aligning the moments or in domain selection. Our analysis and interpretation of these experiments leads to the conclusion that there is still room for a conventional dependency on the magnetic field and that additional neutron scattering data is necessary to clarify this issue. We start from the conventional assumption of tiny antiferromagnetically ordered moments at each U site. These moments ($`\stackrel{}{m}`$) are assumed to lie on the basal plane due to a strong uniaxial anisotropy arising from spin-orbit coupling. In addition, there is an in-plane (hexagonal) anisotropy energy which favors alignment of the moments along any of the three directions perpendicular to the hexagonal lattice vectors (Fig. 1). Neutron-scattering and x-ray experiments show antiferromagnetic order with three possible propagation vectors: $`\stackrel{}{q}_1=\stackrel{}{a}_1^{}/2,\stackrel{}{q}_2=\stackrel{}{a}_2^{}/2,\stackrel{}{q}_3=(\stackrel{}{a}_1^{}\stackrel{}{a}_2^{})/2`$, where $`\stackrel{}{a}_1^{}=\frac{4\pi }{\sqrt{3}a}(1,0,0)`$, $`\stackrel{}{a}_2^{}=\frac{4\pi }{\sqrt{3}a}(1/2,\sqrt{3}/2,0)`$ and $`\stackrel{}{a}_3^{}=\frac{2\pi }{c}(0,0,1)`$ are the reciprocal vectors of the hexagonal lattice with dimensions $`a=5.74`$ Å and $`c=4.89`$ Å. The two U moments in each crystallographic unit cell have to align ferromagnetically in order to account for most of the zero-intensity Bragg points in the diffraction pattern. But, in general, the magnetic structure cannot be fully determined by standard neutron diffraction experiments, since these experiments provide information only about the Fourier components of the magnetic moment. Single- and multi-q magnetic structures display the same magnetic Bragg peaks, and cannot be distinguished unless uniaxial stress or a magnetic field is applied. The magnetic neutron scattering rate per solid angle is proportional to $$\left(\frac{d\sigma }{d\mathrm{\Omega }}\right)_\stackrel{}{Q}\underset{\stackrel{}{Q}_m}{}\|F_M(\stackrel{}{Q})\|^2\delta (\stackrel{}{Q}\stackrel{}{Q}_m)$$ (1) where $`\stackrel{}{Q}`$ is the momentum transfer, $`\stackrel{}{Q}_m`$ are the momenta of the magnetic Brag peaks and $`F_M(\stackrel{}{Q})`$ is the component of the magnetic structure factor perpendicular to the momentum transfer. We can define the magnetic structure factor as $$\stackrel{}{F}_M(\stackrel{}{Q})=\frac{1}{N}\underset{n,j}{}\stackrel{}{m}_{nj}f_{nj}(\stackrel{}{Q})e^{i\stackrel{}{Q}\stackrel{}{R}_{nj}W_j}$$ (2) where $`\stackrel{}{m}_{nj}`$ is the magnetic moment of the $`j^{\text{th}}`$ ion in the $`n^{\text{th}}`$ unit cell, $`f_{nj}`$ is its atomic form factor, $`\stackrel{}{R}_{nj}`$ is its position and $`W_j`$ is the Debye-Waller factor. The spatial distribution of magnetic moments can be Fourier expanded as $`\stackrel{}{m}_{n,j}=_\stackrel{}{q}\stackrel{}{m}_{\stackrel{}{q},j}e^{i\stackrel{}{q}\stackrel{}{R}_n}`$, where the form factor associated with this multi-q magnetic structure is $`\stackrel{}{F}_M(\stackrel{}{Q}=\stackrel{}{Q}_{nm}+\stackrel{}{q})=_j\stackrel{}{m}_{\stackrel{}{q},j}f_j(\stackrel{}{Q})e^{i\stackrel{}{Q}\stackrel{}{r}_jW_j}`$ where $`\stackrel{}{r}_j`$ are the positions of the magnetic ions in the unit cell and $`\stackrel{}{Q}_{nm}`$ label the reciprocal lattice vectors. Thus, in a material with only one type of magnetic ion the scattering rate becomes $$\left(\frac{d\sigma }{d\mathrm{\Omega }}\right)_\stackrel{}{Q}\underset{\stackrel{}{Q}_{nm},\stackrel{}{q}}{}[1(\widehat{Q}\widehat{m}_\stackrel{}{q})^2]\|f(\stackrel{}{Q})\|^2\left\|\underset{\stackrel{}{r}_i}{}e^{i\stackrel{}{Q}\stackrel{}{r}_i}\stackrel{}{m}_\stackrel{}{q}\right\|^2\delta (\stackrel{}{Q}(\stackrel{}{Q}_{nm}+\stackrel{}{q})).$$ (3) Thus, the $`UPt_3`$ diffraction pattern can either be associated with a triple-q structure where $`\stackrel{}{q}_1`$, $`\stackrel{}{q}_2`$ and $`\stackrel{}{q}_3`$ are present at each uranium site or with a single-q structure where separate regions of the crystal will order with different propagation vectors. It has been inferred from the fact that there is no intensity at the $`\stackrel{}{q}_1=[1/2,0,0]`$ position that the magnetic moment lies parallel to its propagation vector. This is the case in the U-monochalcogenides and U-monopnictides with cubic NaCl structure, which order with magnetic moments $`\mu 13\mu _B`$. A moment directed along $`\stackrel{}{q}`$ would also occur for a triple-q structure, but it is not clear that this condition must be fulfilled in the single-q structure. The intensity of $`\stackrel{}{q}_2=[0,1/2,0]`$ and $`\stackrel{}{q}_3=[1/2,1/2,0]`$ peaks has not been reported for UPt<sub>3</sub>. It is possible that the sample preparation methods make domain “1” (Fig. 1) preferable over domains “2” and “3”. However, measuring the intensity of these three peaks in the same single crystal would allow one to determine if the magnetic moments do lie parallel to the propagation vector of the domain. Below we discuss the field dependence of the magnetic neutron scattering intensity for the possible magnetic structures. We first discuss the field dependence of single-q structures, then we comment on the possibility of a triple-q magnetic structure. The magnetic unit cell of a single-q structure results from doubling the hexagonal unit cell along one in-plane direction, reducing the hexagonal symmetry to orthorhombic. Transmission electron microscope images provide direct observation of basal plane, as well as prism plane, stacking faults in pure single crystals. These defects are observed even in the crystals with the highest residual resistance ratios. We hypothesize that these defects pin AFM domain walls in the ab-plane and fix the spatial distribution of domains. In an antiferromagnet the Zeeman energy prefers the staggered magnetization to be perpendicular to the field. Thus, a sufficiently strong magnetic field applied in the hexagonal plane will give rise to domain reorientation by overcoming the in-plane anisotropy energy. The magnitude of the staggered magnetization will remain roughly the same, modulated only by a small in-plane anisotropy energy. Therefore, for a given magnetic Bragg peak, the ratio between the scattering rate at high field and at zero field is $$r=\frac{d\sigma /d\mathrm{\Omega }\|_H\mathrm{}}{d\sigma /d\mathrm{\Omega }\|_{H=0}}\frac{1(\widehat{Q}\widehat{m}_H\mathrm{})^2}{1(\widehat{Q}\widehat{m}_{H=0})^2},$$ (4) where $`\mathrm{}`$ refers to an average over domains. Let us analyze the experimental data based on Eq. (4). The staggered magnetization lies on the basal plane, $`\widehat{m}=(\mathrm{cos}\theta ,\mathrm{sin}\theta )`$. Van Dijk et al. chose a configuration with H parallel to the “a” axis ($`\theta _H=30^o`$ in Eq. 9) and a momentum transfer $`\stackrel{}{Q}=[1/2,0,1]=2\pi ((1/\sqrt{3}a),0,(1/c))`$, which gives $`\widehat{Q}=(0.441,0,0.897)`$ and $$r=\frac{1(0.441\mathrm{cos}(\theta _H+\pi /2))^2}{1(0.441\mathrm{cos}(\theta ))^2}=1.05$$ (5) for three equally populated magnetic domains. This ratio can be increased to $`r=1.18`$ by assuming that only the domain with the staggered magnetization parallel to the propagation vector is populated (domain “1” in Fig. (1)). Thus, even in the case of complete domain reorientation, the neutron scattering rate at $`\stackrel{}{Q}=[1/2,0,1]`$ in high fields can increase at most by 18% over its value at zero field. Figure 2 shows the experimental data and the theoretical curves for a model with equally populated domains and for a model with only domain “1” populated. Although the theoretical calculation associated with domain “1” is in good agreement with the data, it is not possible to conclude whether or not the U moments rotate with the field because of the small change in intensity that is expected for this Bragg peak and the large error bars that are reported for the intensity. Note that the error bars for this measurement are comparable to the maximum change in the intensity ratio. In our calculation we have assumed an anisotropy field of $`H_{\text{an}}=1.5T`$. However, much smaller values are consistent with the limited data. The precise value of the additional parameters in our model play a role only in the region of small magnetic fields. For fields $`H>2H_{an}`$ the ratio between the intensity at high fields and at zero field saturates at its upper limit, which is determined by purely geometrical arguments. Earlier analysis was based on the assumption that the staggered magnetic moment is always parallel to its propagation vector. Thus, it was expected that a sufficiently high magnetic field parallel to the “a” axis would select domain “2” with propagation vector $`\stackrel{}{q}_2`$ throughout the sample. As a consequence, the magnetic intensity at $`\stackrel{}{Q}=[1/2,0,1]=\stackrel{}{q}_1+[0,0,1]`$ was expected to drop to zero. However, as we show in Fig. 2, if we assume that the spatial distribution of domain walls is pinned, the form factors for $`\stackrel{}{Q}=[1/2,0,1]`$, which is a vector mostly out of the hexagonal plane, lead to a much smaller variation of the intensity with the field. Larger expected ratios between the low- and high-field intensities are obtained with the experimental setup used by Lussier et al. They measured the neutron scattering cross-section at three different momentum transfers, all in the basal plane: $`\stackrel{}{Q}_1=[1/2,1,0],\stackrel{}{Q}_2=[3/2,1/2,0]`$ and $`\stackrel{}{Q}_3=[1,3/2,0]`$. The magnetic field was oriented along the $`𝐛`$ axis. Lussier et al. report data for $`\stackrel{}{Q}_1`$ and $`\stackrel{}{Q}_2`$, and magnetic fields up to $`3.5T`$. We can estimate from Eq. (4) the ratio between high- and zero-field intensity for any distribution of domains in the crystal. A crystal with equally populated domains will display the following ratios for the neutron scattering rate at high fields and zero field: $$r(\stackrel{}{Q}_1)=0.86,r(\stackrel{}{Q}_2)=0.21,r(\stackrel{}{Q}_3)=1.93.$$ (6) If domain 3 is unpopulated and domains 1 and 2 are equally populated the ratios should be: $$r(\stackrel{}{Q}_1)=1.60,r(\stackrel{}{Q}_2)=0.20,r(\stackrel{}{Q}_3)=1.38,$$ (7) and if only the domain with the magnetization parallel to the propagation vector is occupied (e.g. domain “1” for $`\stackrel{}{Q}_1`$) then, $$r(\stackrel{}{Q}_1)=1,r(\stackrel{}{Q}_2)=0.25,r(\stackrel{}{Q}_3)=2.25.$$ (8) Figure 3 displays the experimental data of Ref. , and theoretical calculations for two samples, one with domains “1” and “2” equally populated at zero field, another with domains “1” and “2” unequally populated. The parameters of the model are the same ones used to fit the data at $`\stackrel{}{Q}=[1/2,0,1]`$ in Fig. 2. We conclude that the limited data for $`\stackrel{}{Q}_1`$ and $`\stackrel{}{Q}_2`$ is roughly consistent with either one or two unequally populated domains, particularly if $`H_{an}2.5\text{T}`$. Previous analysis of these results was also based on the assumption that the propagation vector of the magnetic domains follows the rotation of the magnetic moments. Thus, at high fields it was expected that the intensity of the $`\stackrel{}{Q}_2`$ and $`\stackrel{}{Q}_3`$ peaks would be suppressed to zero, while increasing the intensity of the $`\stackrel{}{Q}_1`$ peak to roughly three times its zero field value. The theoretical curves have been calculated using the free energy functional, $`\overline{F}_{AFM}=2(1\overline{T})\|\stackrel{}{m}_0\|^2+\|\stackrel{}{m}_0\|^4+\overline{U}_{an}\|\stackrel{}{m}_0\|^6(r_6\mathrm{cos}(6\theta ))+\overline{U}_{an}\overline{H}^2\|\stackrel{}{m}_0\|^2\mathrm{cos}^2(\theta \theta _H))+`$ (9) $`+r_D\overline{U}_{an}\overline{H}\|\stackrel{}{m_0}\|\|\mathrm{sin}(\theta \theta _H))\|+r_{st}\|\stackrel{}{m}_0\|^2(\left({\displaystyle \frac{(\mathrm{cos}(\theta ))}{H}}\right)^2+\left({\displaystyle \frac{(\mathrm{sin}(\theta ))}{H}}\right)^2)`$ (10) where all energies are measured in units of the exchange energy, $`U_{\text{ex}}`$, which is defined as the absolute value of the free energy at zero temperature and field in the absence of any anisotropy energy. The magnetic order parameter is restricted to the basal plane by the large uniaxial anisotropy energy (not shown in Eq. 9) and it is measured with respect to the antiferromagnetic order parameter in the exchange approximation: $`\stackrel{}{m}_0=\stackrel{}{m}/\|\stackrel{}{m}_{ex}\|=\|\stackrel{}{m}_0\|(\mathrm{cos}\theta ,\mathrm{sin}\theta ,0)`$. The renormalized temperature is defined as $`\overline{T}=T/T_N`$, with $`T_N`$ as the Néel temperature. The magnetic field $`\overline{H}`$ is measured in units of the in-plane anisotropy field, $`H_{an}`$. The first two terms of the free energy correspond to the exchange energy. For $`\overline{T}<1`$ antiferromagnetic order with magnetic moment $`\|\stackrel{}{m}_0\|=\|\stackrel{}{m}\|/\|\stackrel{}{m}_{ex}\|=\sqrt{1\overline{T}}`$ and free energy $`\overline{F}_{AFM}=F_{AFM}/U_{ex}=(1\overline{T})^2`$ is stable. The sixth-order term is the leading term in the in-plane anisotropy energy; it favors alignment along the three directions perpendicular to the hexagonal lattice vectors: $`\theta =n(\pi /3)`$, where $`n`$ is an integer. The in-plane anisotropy energy induces a hexagonal modulation of the upper critical field as a function of the orientation of the field in the basal plane. From the magnitude of this hexagonal modulation we estimate an anisotropy energy of $`\overline{U}_{an}=U_{an}/U_{ex}0.02`$. The parameter $`r_6`$ must be bigger than one in order to have a stable free energy. We use $`r_6=1.5`$ in our calculations, however, its precise value does not play any significant role in the minimization of the free energy. The fourth term is the Zeeman energy for an antiferromagnet, $`F_Z=g(\stackrel{}{m}\stackrel{}{H})^2`$, which is quadratic in $`H`$ and favors perpendicular alignment ($`g>0`$) of the staggered moment and the magnetic field. This term can be written in the form, $$F_Z=\frac{U_{an}}{U_{ex}}\left(\frac{H}{H_{an}}\right)^2\left(\frac{\stackrel{}{m}}{\|\stackrel{}{m}_{ex}\|}\right)^2\mathrm{cos}^2(\theta \theta _H),$$ (11) where $`H_{an}=(1/\|\stackrel{}{m}_{ex}\|)\sqrt{U_{an}/(gU_{ex})}`$ and $`\theta _H`$ is the angle of the magnetic field with the $`\stackrel{}{a}_1^{}`$ reciprocal vector. The fifth term in Eq. 9 is the Dzyaloshinskii-Moriya term describing the linear coupling of the sublattice magnetization to the magnetic field, $`F_D=g^{}𝐝\left(𝐇\times 𝐦_0\right)`$. This term corresponds to the Zeeman coupling of a weak ferromagnetic (FM) moment in systems which are predominantly antiferromagnetic. Its origin is the anisotropic superexchange coupling between magnetic moments, $`\stackrel{}{D}_{ij}\stackrel{}{S}_i\times \stackrel{}{S}_j`$, where $`\stackrel{}{D}_{ij}`$ are the Moriya vectors for different bonds on the lattice, and which are related to each other by lattice symmetries. In the case of $`UPt_3`$, $`\stackrel{}{D}_{ij}=0`$ when $`i`$ and $`j`$ are nearest-neighbor U sites, while $`\stackrel{}{D}_{ij}=\pm \|d\|\widehat{c}`$, independent of the direction of the staggered magnetic moment, when $`i`$ and $`j`$ refer to next-nearest-neighbor U atoms. This superexchange coupling generates the Dzyaloshinskii term in the free energy which can be expressed as $`\overline{F}_D=r_D\overline{U}_{an}\overline{H}\|\stackrel{}{m_0}\|\|\mathrm{sin}(\theta \theta _H))\|`$ shown in Eq. 9. For low temperatures the effect of the Dzyaloshinskii-Moriya term is to generate a tiny ferromagnetic moment at the price of a small reduction in the magnitude of the staggered moment. However, for temperatures close to $`T_N`$, the Dzyaloshinskii-Moriya energy is comparable to the exchange energy, and leads to a significant reduction in the magnitude of the AFM moment and, as a consequence, the intensity of the magnetic Bragg peaks. We can define a crossover temperature in terms of the parameters of the free energy, $`\overline{T}_D=1\sqrt[3]{r_D^2U_{an}^2\overline{H}^2}`$. Although the staggered moment vanishes precisely at the Néel temperature, for $`\overline{T}_D<\overline{T}<1`$ the moment decreases rapidly before the transition at $`\overline{T}=1`$. Thus, $`\overline{T}_D`$ could be mis-identified as the Néel temperature of the sample. The Dzyaloshinskii-Moriya term provides an explanation for the crossing of the intensity curves for zero and high fields as a function of temperature as shown in Fig. 4. The Dzyaloshinskii-Moriya coupling also provides an explanation for the linear term in the field dependence of the magnetoresistance, which onsets at the Neèl transition and increases for $`T<T_N`$. It has been shown that a linear term in the transverse magnetoresistance is present in antiferromagnetic structures admitting the existence of weak ferromagnetism. Indeed it follows from Onsager relations for the resistivity that a magnetoresistance which is linear in field in a AFM requires the Dzyaloshinskii-Moriya coupling. Finally, the last term in Eq. 9 describes the “stiffness” of the order parameter with respect to rotations in the ab-plane. This stiffness originates from the formation of domains in which the staggered moment points in the same direction within each domain. An inhomogeneous domain structure gives rise to domain walls separating differently oriented domains. The energy associated with the domain wall is obtained from the gradient energy, $`\kappa _{ijkl}(m_j/x_i)(m_l/x_k)`$, which must be included in the free energy functional. For an individual domain wall, the gradient energy can be written as an integral over the domain wall surface $`\mathrm{\Omega }`$, $$F_{\text{wall}}_\mathrm{\Omega }𝑑\mathrm{\Omega }_{\sigma _1}^{\sigma _2}\left[\left(\frac{\widehat{𝐦}_𝐱}{\sigma }\right)^2+\left(\frac{\widehat{𝐦}_𝐲}{\sigma }\right)^2\right]𝑑\sigma $$ (12) where $`\sigma `$ is the coordinate perpendicular at each point to the wall surface. The width of the wall is given by $`\sigma _2\sigma _1`$ and $`\widehat{𝐦}_𝐱`$, $`\widehat{𝐦}_𝐲`$ are the components of the unit vector $`\widehat{𝐦}=𝐦/\|𝐦\|`$. This unit vector satisfies the boundary conditions, $`\widehat{𝐦}(\sigma _2)=\widehat{𝐦}_{eq}(H+\mathrm{\Delta }H)`$ and $`\widehat{𝐦}(\sigma _1)=\widehat{𝐦}_{eq}(H)`$, where $`\widehat{𝐦}_{eq}(H)`$ is the equilibrium orientation of the staggered magnetic moment in the presence of a magnetic field $`\stackrel{}{H}`$. In quasiequilibrium the direction of the magnetic moment evolves smoothly through the domain wall between its values corresponding to different equilibrium field orientations, $`\widehat{𝐦}_{eq}(H+\mathrm{\Delta }H)`$ and $`\widehat{𝐦}_{eq}(H)`$. By scaling the width of the domain wall to $`\mathrm{\Delta }H`$ we obtain the stiffness energy in the form of the last term in Eq. 9. The stiffness energy is important in the region of intermediate fields, where the normalized neutron intensity increases from a value close to the one at zero field to its value at high fields. The initial drop of the neutron intensity as a function of the applied field (Fig. (2) and (3)) is a combined effect of the anisotropy and stiffness energies. This drop is due to an initial reduction of the magnitude of the magnetic moment. Small fields do not induce rotation; instead the magnitude of the staggered moment is reduced. Higher fields are able to rotate the moments by overcoming the anisotropy and stiffness energies. Consequently, the Zeeman energy is reduced to zero and the rotated moment recovers its value at zero field. So far we have discussed single-$`𝐪`$ structures or multi-domain single-$`𝐪`$ structures. Triple-q structures are also possible. By symmetry each component, $`\stackrel{}{m}_{\stackrel{}{q}_i}`$, has the same amplitude. Triple-q antiferromagnetic order occurs in the $`NaCl`$-type monopnictide $`USb`$, in the $`CsCl`$-type $`DyAg`$ and $`NdZn`$, and in the $`AuCu_3`$-type $`TmGa_3`$. These materials are cubic and the three Fourier components $`\stackrel{}{m}_{\stackrel{}{q}_i}`$ point along mutually perpendicular axes leading to the condition of a uniform magnitude of the moment. For a triple-$`𝐪`$ structure in UPt<sub>3</sub>, in order to explain the vanishing intensity at the (1/2,0,0) Bragg point we are required to have $`\stackrel{}{m}_{\stackrel{}{q}_1}`$ parallel to $`\stackrel{}{q}_1`$ and by symmetry the other two moments must also be parallel to their propagation vectors. Thus, the magnetic moment of both U ions in the n$`^{\text{th}}`$ unit cell is given by $$\stackrel{}{m}_n=\|m\|\underset{i=1}{\overset{3}{}}\widehat{𝐪}_ie^{i(\varphi _i\stackrel{}{𝐪}_i𝐑_n)}.$$ (13) It can be easily shown that it is not possible to satisfy the condition of equal magnitude of the moment at every U site. Most choices for the phases $`\varphi _1,\varphi _2,\varphi _3`$ produce a non-uniform distribution of the magnitude of the U magnetic moment. For example, Fig. 5 displays a possible spatial distribution of the moments. The three Fourier components of the triple structure have been chosen with equal phase $`\varphi _1=\varphi _2=\varphi _3`$. The magnetic unit cell is then constructed from four unit cells containing eight U ions, reducing the hexagonal symmetry to monoclinic. Note that the two U ions in the central cell have zero net moment, while the other six U ions have equal values for the magnitude of the moment. Even though a triple-q magnetic structure in UPt<sub>3</sub> is compatible with the neutron-scattering experiments the resulting non-uniform magnetization is unusual, but not unique. The triple-q magnetic structure in UPt<sub>3</sub> is similar to the magnetically frustrated structure of the uranium intermetallic $`UNi_4B`$, which also has a hexagonal crystal lattice. This material orders antiferromagnetically around $`T_N=30K`$, with approximately $`1/3`$ of the U spins remaining paramagnetic well below $`T_N`$. It has been suggested that the competition between the Kondo effect, the antiferromagnetic exchange interaction and the frustration of the crystallographic lattice is responsible for the unusual $`UNi_4B`$ magnetic structure. Such an interplay between competing interactions could also take place in UPt<sub>3</sub>. However, to our knowledge, there is no other indication of such a frustrated magnetic structure in $`UPt_3`$. Note that a triple-$`𝐪`$ structure does not preclude the coupling between the AFM and superconducting order parameters, which is considered a good candidate for the proposed SBF in the 2D order parameter models for the superconducting phases. The SBF coupling is non-vanishing for triple-$`𝐪`$ structures, except for the special case in which all three phases are identical. The coupling between the superconducting, $`\stackrel{}{\eta }=(\eta _1,\eta _2)`$, and the magnetic order parameters is $`F_{\text{AFM-SC}}A(\|\eta _1\|^2\|\eta _2\|^2)+B(\eta _1\eta _2^{}+\eta _1^{}\eta _2)`$, with $`A=_{n=1,4}(m_x^2(n)m_y^2(n))=42\mathrm{cos}^2(\varphi _2\varphi _1)2\mathrm{cos}^2(\varphi _3\varphi _1)`$, $`B=2_{n=1,4}(m_x(n)m_y(n))=2\sqrt{3}(\mathrm{cos}^2(\varphi _2\varphi _1)\mathrm{cos}^2(\varphi _3\varphi _1))`$, where the summation refers to the four unit cells contained in the magnetic unit cell shown in Fig. 5. The hexagonal triple-q shown in Fig. 5 resembles the antiferroquadrupolar order reported for $`UPd_3`$. Furthermore, $`Pt`$ and $`Pd`$ are isoelectronic, their nearest neighbor U-U distances are almost identical, and both systems have a hexagonal closed packed structures. However, the magnetic and electronic properties of $`UPt_3`$ and $`UPd_3`$ are very different. In fact $`UPd_3`$ is a localized material with well-defined crystal-field levels. Several measurements on $`UPd_3`$ show two phase transitions at $`7K`$ and $`5K`$. The transition at $`7K`$ is believed to correspond to a quadrupolar ordering of the U ions, which is accompanied by a modulated lattice distortion. The $`5K`$ transition is magnetic, with an ordered moment that is very small, as in UPt<sub>3</sub>, $`\mu 0.01\mu _B`$/U-ion. But, the moments in $`UPd_3`$ are pointing out of the basal plane. We conclude with a brief discussion of possible neutron scattering experiments which might clarify the magnetic order in UPt<sub>3</sub>. A zero-field systematic measurement of the intensity of a number of magnetic peaks in the same single crystal will determine whether the magnetic moments are indeed parallel to the propagation vector or not. Using previous experimental arrangements it would be very interesting to apply fields well above $`3T`$ and measure the intensity at three independent momentum transfers. Although polarized inelastic neutron-scattering experiments have been performed in $`UPt_3`$, the magnetic Bragg peaks have not been studied with polarized neutrons. Polarized elastic neutron-scattering would provide confirmation of the magnetic nature of the transition. This powerful method has been used successfully on $`UPd_3`$ to identify the magnetic nature of the second phase transition at $`T_2=5K`$. In summary, based on available neutron diffraction data, the magnetic field dependence of the neutron scattering intensity is consistent with antiferromagnetic order in $`UPt_3`$ based on the most conventional assumption of a single-q structure with three equivalent domains. However, a triplet-q structure is also consistent with these experiments. If realized the triple-$`𝐪`$ structure would imply a non-uniform, frustrated magnetic structure in the crystal. We thank Piers Coleman, Bill Halperin and Robert Joynt for valuable discussions on this subject. The hospitality of the Aspen Center for Physics during the 1999 Summer workshop on unconventional order in metals is gratefully acknowledged. This research was supported by NSF grants DMR 9705473, DMR 9972087 and DMR 91-20000 through the Science and Technology Center for Superconductivity.
warning/0004/gr-qc0004041.html	ar5iv	text	# Geometric Cone Surfaces and (2+1) - Gravity coupled to Particles ## 1 Introduction. Globally hyperbolic matter-free (2+1) - spacetimes with compact space of genus $`g1`$ and cosmological constant $`\mathrm{\Lambda }=0`$ have been fairly well understood. These spacetimes can be arranged into classes up to Teichmüller equivalence that is, roughly speaking, up to isometry isotopic to the identity. Both the geometric-time-free approach, which eventually identifies each spacetime by its geometric holonomy (c.f. \[W\], \[Me\]), and the cosmological approach based on fibration by constant mean curvature space-like surfaces (see \[Mo\], \[A-T-M\]), have recognized the Cotangent Bundle of the Teichmüller space $`T_g`$ (which is homeomorphic to $`B^{6g6}\times \text{}^{6g6}`$, when $`g2`$) as the parameter space. When $`g2`$, the correspondence between these two approaches is still rather implicit; nevertheless, we have shown in \[B-G 1\] that the canonical Cosmological Time Function (CTF) (that is the length of time that the events have been in existence) provides a very good cosmological resolution of the matter-free (2+1) gravity. For instance, the asymptotic states of (CTF) recover and decouple the linear part and the translation part of the geometric holonomy; the orbit in $`T_g`$ of the (CTF) has nice properties; the “initial” singularity with respect to (CTF) can be accurately described in terms of the degeneration of the geometry of the level surfaces. It turns out that all matter-free spacetimes are obtained by means of two basic constructions: (a) the Minkowskian suspension of flat or hyperbolic surfaces; (b) a distinguished kind of deformations of the hyperbolic suspensions. In this paper we shall describe the extension of constructions (a) and (b) to the case of gravitating particles in (2+1) dimensions. When gravity is coupled to particles, the picture is far from being exhaustive. We will be concerned with compact spaces with a finite number of massive particles. ’t Hooft’s approach \[’t H\] describes these spacetimes by means of the “linear” evolution of a special kind of Cauchy surfaces which are tiled by spatial planar polygons. The extrinsic curvature is null in the interior of each tile and it is singular along the edges; the evolution includes the changing of tiling combinatorics under codified transition rules. Each such a Cauchy surface is intrinsecally a flat surface with conical singularities. Among these singularities, some correspond to the intersection with particle world-lines (the spacetime has a concentrated curvature along these lines); the others are 3-dimensionally apparent singularities (but the Gauss-Bonnet constraint implies that, in general, they cannot be avoided ). Each globally hyperbolic spacetimes contains such a kind of Cauchy surface with, at least locally, such a kind of evolution; it is not clear to us if the evolution of a given surface necessarily fills all the spacetime and how the evolutions of different surfaces in the same spacetime are related each other. So, it seems hard to recover from this approach a clear identification of the parameter space. Another experimented approach (see \[Mn-S\], \[B-C-V\]) is the classical ADM formalism with the so called “instantaneaous gauge”, that requires fibration by spatial Cauchy surface with zero extrinsic curvature. This last requirement is technically very useful and allows to analytically find solutions by means of classical and very elegant mathematical tools. Unfortunately, it turns out that the only spacetimes with compact space covered by this approach are the static ones (that is, by using the terminology of the present paper, the static Minkowskian suspensions of flat surfaces with conical singularities, which we will see below). The aim of this note is to describe the spacetimes with compact space of genus $`g0`$ and $`r`$ gravitating particles ($`\mathrm{\Lambda }=0`$) that one can obtain by means of three kinds of construction: (a) the Minkowskian suspensions of flat or hyperbolic surfaces with conical singularities; (b) the distinguished deformations of hyperbolic suspensions (in strict analogy with the matter-free case); (c) the patchworking of Minkowskian suspensions (this is peculiar to gravity coupled with particles). These spacetimes have very transparent structural properties and behave somewhat similarly to the matter-free ones with respect to, for instance, the canonical (CTF), its asymptotic states, the initial singularity and so on. Moreover they form a rather wide class of spacetimes, so that we can derive from them some non trivial information about the actual parameter space. For example we will show that when the masses are big enough and the cone points are suitaby spaced (roughly speaking), then the distinguished deformations of hyperbolic suspensions determine a relevant non empty open subset of the parameter space of the form $`𝒰\times \text{}^{6g6+2r}`$, where $`𝒰`$ is an open set of the Teichmüller Space $`T_g^rB^{6g6+2r}`$. On the other hand, by patchworking of suspensions, we will produce spacetimes with the same masses of certain hyperbolic suspensions but which are not equivalent to any distinguished deformation of them. In fact we will guess that they belong to different connected components of the parameter space. So gravity coupled to particles seems to be much more flexible than pure gravity. In the last section we will state several related questions and we will develop few speculations. Several constructions concerning Minkowskian suspensions run, with minor changes, as in the matter-free case; so we will refer to the diffuse paper \[BG 1\] for more details. ## 2 Geometric Surfaces with Conical Singularities. Cone points. The local models of flat or hyperbolic surfaces at a conical singularity are respectively given, in complex coordinate, by the metrics on $`\{\|z\|<1\}`$ (set $`\alpha >0`$): $$ds_{(E,\alpha )}^2=\alpha ^2\|z\|^{2\alpha 2}\|dz\|^2,$$ $$ds_{(H,\alpha )}^2=\alpha ^2[2/(1\|z\|^{2\alpha })]^2\|z\|^{2\alpha 2}\|dz\|^2.$$ They are obtained by pull-back of the standard Euclidean or Poincaré metrics on $`\{\|w\|<1\}`$ via the map $`w=z^\alpha `$. In both cases the concentrated curvature at the conical point is $`k=2\pi (1\alpha )`$, the cone angle is $`2\pi \alpha `$. In order to have a genuine singularity, $`\alpha 1`$. Geometric cone surfaces. It is convenient to adopt the formalism of geometric $`(X,G)`$-manifolds (see, for instance, chapter B of \[B-P\] or section 3 of \[B-G 1\]). Fix a base compact oriented surface $`F_g`$ of genus $`g0`$ and fix $`p_1,\mathrm{},p_r`$ points on $`F_g`$. A marked geometric (i.e. flat or hyperbolic) surface with conical singularities, of cone angles $`2\pi \alpha _i`$, $`i=1,\mathrm{},r`$, is a homeomorphism $$\varphi :(F_g,\{p_i\})(S,\{q_i\})$$ such that $`S^{}=S\{q_i\}`$ is a $`(\text{}^2,Isom^+(\text{}^2))`$\- (resp. a $`(\text{}^2,Isom^+(\text{}^2))`$-) surface and its metric completion has a conical singularity of cone angle $`2\pi \alpha _i`$ at $`q_i`$. Gauss-Bonnet constraint. The classical Gauss-Bonnet formula leads to the following relations. Flat Case: (Gauss-Bonnet equality) $$\underset{i}{}k_i=2\pi \underset{i}{}(1\alpha _i)=2\pi (22g).$$ Hyperbolic Case: $$\underset{i}{}k_i=2\pi \underset{i}{}(1\alpha _i)=2\pi (22g)+Area(S).$$ whence: (Gauss-Bonnet inequality) $$\underset{i}{}(1\alpha _i)>22g.$$ This implies, in any case, that when $`g=0`$, necessarily $`r3`$, and we will make this assumption by default. We say that $$\delta =(\text{𝕏},g,[\alpha ]_r)=(\text{𝕏},g,(\alpha _1,\mathrm{},\alpha _r))$$ (where $`\text{𝕏}=\text{}^2`$ or $`\text{}^2`$, $`g0`$ and the $`\alpha _i`$’ satisfy the appropriate Gauss-Bonnet equality or inequality), is a virtual type of geometric surfaces with conical singularities. For a fixed type $`\delta `$ we denote by $`T_\delta `$ the Teichmüller space of marked surfaces of type $`\delta `$, (regarded up to Teichmüller equivalence \- see, for instance, section 4 of \[B-G 1\] for more details). When $`r>0`$, the fundamental group $`\pi (F_g^{}=F_g\{p_i\})`$ is a non Abelian free group with $`s=2g+r+1`$ generators. For each $`[\varphi ]T_\delta `$ it is well defined, up to conjugation, the holonomy representation $$\rho _{[\varphi ]}:\pi (F_g^{})Isom^+(\text{𝕏}).$$ The universal covering $`p:S^{}S`$ is, in a natural way, a local isometry so that $`S^{}`$ is homeomorphic to $`\text{}^2`$ and it is endowed with a geometric structure with conical singularities. We have also the developing map (well defined up to left action of $`Isom^+(\text{𝕏})`$) $$D_{[\varphi ]}:(S^{})^{}\text{𝕏}.$$ $`(S^{})^{}`$ is also homeomorphic to $`\text{}^2`$ and is endowed with a smooth geometric structure. We can choose the representatives in such a way that, for every $`\gamma \pi (S^{})`$, for every $`x(S^{})^{}`$, $$D_{[\varphi ]}(\gamma (x))=\rho _{[\varphi ]}(\gamma )(D_{[\varphi ]}(x)).$$ $`D_{[\varphi ]}`$ is a local isometry; when $`r>0`$, $`S^{}`$ and $`(S^{})^{}`$ are not (metrically) complete and, equivalently, $`D_{[\varphi ]}`$ is not a global isometry. Orbifolds. Geometric 2-dimensional compact orbifolds (with only conical singularities) make a special class of surfaces we are concerned with. Such an orbifold $`S`$ is a quotient $`\text{𝕏}/\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is a group of isometries of 𝕏 acting properly discontinuosly and such that the set of points with non trivial stabilizer is made by isolated points. For a genuine orbifold this set is nonempty. They are classified as follows (see \[T\], \[Sc\]). ###### Proposition 2.1 A geometric cone surface is a genuine Euclidean orbifold iff it is of one of the types $`(\text{}^2,0,(1/2,1/3,1/6))`$, $`(\text{}^2,0,(1/2,1/4,1/4))`$, $`(\text{}^2,0,(1/3,1/3,1/3))`$, $`(\text{}^2,0,(1/2,1/2,1/2,1/2))`$. A geometric cone surface is a genuine hyperbolic orbifold iff it is of a type $`(\text{}^2,g,[\alpha ]_r)`$ satisfying the Gauss-Bonnet inequality and such that each $`\alpha _i[\alpha ]_r`$ is of the form $`\alpha _i=1/n_i`$, $`n_i\text{}^{}`$. Moreover all these types are actually realized by orbifolds. Conformal structures. Associated to each geometric structure with conical singularities there is a natural conformal structure: a conformal atlas is simply obtained by adding a chart in complex coordinates as above at each conical singulatities, to any atlas of the geometric structure on $`S^{}`$ (use the Poincaré disk model for $`\text{}^2`$). So, for each virtual type $`\delta `$, if $`T_g^r`$ is the classical Teichmüller space of conformal structures on $`F_g`$, relatively to the marked points $`p_i`$ (which, as it is well known, is homeomorphic to an open ball $`B^{6g6+2r}`$), there is a natural continuos map (in fact in the case of flat surfaces we take off simple rescaling by normalizing the area to be equal to $`1`$) $$\psi _\delta :T_\delta T_g^r.$$ Geometric surfaces with conical singularities are classified by the following proposition. ###### Proposition 2.2 For any virtual type $`\delta `$, the natural map $`\psi _\delta `$ is a homeomorphism. The flat case is due to Troyanov (see \[Tr\]). The orbifold case is treated in \[T\]. Let us sketch the main steps of a proof in the general hyperbolic case. (1) dim ($`T_\delta )=`$ dim ($`T_g^r`$). Let us outline first a way to construct all hyperbolic cone surfaces. Fix $`(F_g,\{p_1,\mathrm{},p_r\})`$ as before. A standard spine of $`F_g^{}`$ is a 1-complex $`P`$ embedded in $`F_g^{}`$, with only 3-valent vertices, such that $`F_g^{}`$ is a regular neighbourhood of $`P`$ ($`F_g^{}`$ retracts onto $`P`$). Associated to such a $`P`$ there is a dual (topological) ideal triangulation $`\tau _P`$ of $`F_g^{}`$, that is a “relaxed” (i.e. multiple and self adjacencies between triangles are allowed) triangulation of $`F_g`$, having $`\{p_1,\mathrm{},p_r\}`$ as set of vertices. If $`v(P)=\|V(P)\|`$ denotes the number of vertices of $`P`$ (i.e. the number of triangles of $`\tau _P`$), $`e(P)=\|E(P)\|`$ the number of its edges (i.e. the number of the edges of the dual triangulation), one has $`3v(P)=2e(P)`$ so that $`e(P)=6g6+3r`$. Clearly spines exist. Fix a spine $`P`$. For any admissible map $`f:E(P)\text{}^+`$ (i.e. a map such that at each vertex $`vP`$ the values of $`f`$ on the three edges emanating from $`v`$ satisfy the triangular inequalities), we can construct a hyperbolic surface with $`r`$ conical singularities. This is obtained as a geometric realization of the dual triangulation $`\tau _P`$, by using hyperbolic triangles with edge lengths prescribed by $`f`$. Recall that each hyperbolic triangle is determined by the edge lengths as well as by the interior angles, and there are classical explicit formulas relating lengths and angles. It is not too hard to see that varying the spine and the admissible function, one can realize all the hyperbolic virtual types. On the other hand, any cone hyperbolic surface arises in this way. In fact let $`(F_g,F_g^{})(S,S^{})`$ be such a surface. Consider the subset $`Q`$ of $`S`$, such that for each $`xQ`$ there exist $`ij`$ such that $`d(x,p_i)=d(x,p_j)`$. Generically $`Q`$ is a standard spine of $`S^{}`$; the interior of an edge of $`Q`$ consists of the points with exactly two equidistant marked points $`p_i,p_j`$, the same along the given edge. The “axis” of each edge, that is the geodesic arc connecting $`p_i`$ and $`p_j`$ and passing from the point of the edge of minimal distance from them, are the edges of a geometric realization on the dual triangulation $`\tau _Q`$. In general $`Q`$ is a spine, possibly with higher valency vertices; the same procedure produces a dual ideal cellularization of $`S^{}`$ by convex hyperbolic polyhedra and we eventually obtain a geometric triangulation by subdividing without introducing new vertices. If a virtual type $`\delta `$ is realized by an admissible map $`f_0`$ on $`E(P)`$, the maps realizing the same type are obtained by imposing $`r`$ independent conditions. So one can deduce, at least, that $`T_\delta `$ is a topological manifold of the right dimension $`6g6+2r`$. (2) The map $`\psi _\delta `$ is injective. Consider $`\text{}^2`$ in the Poincaré disk model $`D=\{\|z\|<1\}`$, and let $`e^{2h}\|dz\|^2`$ be the standard Poincaré distance. Realize a given element $`\sigma `$ of $`T_g^r`$ by a smooth hyperbolic surface (with marked points) $`S=D/\mathrm{\Gamma }`$. Two hyperbolic surfaces with conical singularities of the same type, both representing $`\sigma `$, are given by two metrics $`e^{2(h+h_i)}\|dz\|^2`$, $`i=1,2`$, such that each $`h_i`$ is a $`\mathrm{\Gamma }`$-equivariant function on $`D`$, with the same kind of singularities over the marked points of $`S`$. It follows that $`h_1h_2`$ is a real analytic $`\mathrm{\Gamma }`$-equivariant function on $`D`$ satisfying the Liouville equation $$\mathrm{\Delta }(h_1h_2)=e^{2h}(e^{2h_1}e^{2h_2}).$$ As $`S`$ is compact $`h_1h_2`$ has maxima and minima. Either $`\mathrm{\Delta }(h_1h_2)>0`$ near a maximum, or $`\mathrm{\Delta }(h_1h_2)0`$ near a minimum. By the maximun principle $`h_1h_2`$ is constant near the minumum or the maximun and hence it is constant (and necessarily equal to $`0`$) everywhere. (3) Conclusion. By the invariance of domain theorem, $`\psi _\delta `$ is a homeomorphism onto a non empty open subset of $`T_g^r`$. To conclude it is enough to show that the image of $`\psi _\delta `$ is closed. This can be done by studing the convergence of the conformal factors (see the above step), or by arguing (via geometric considerations) that the image of a “diverging” sequence in $`T_\delta `$ is diverging in $`T_g^r`$. ## 3 Minkowskian Suspensions. Particle world lines. Let us give, first, the local models of the line of universe of a massive particle. They are obtained by “suspension” of the local models for geometric cone surfaces. We can take indifferently, in coordinates $`(z,t)`$, $$d\sigma _{(E,\alpha )}^2=dt^2+ds_{(E,\alpha )}^2,$$ or, assuming $`t>0`$ $$d\sigma _{(H,\alpha )}^2=dt^2+t^2ds_{(H,\alpha )}^2.$$ They are equivalent as local models, in the sense that any point $`(0,t_1)`$ in the first model and any point $`(0,t_2)`$ in the second one have isometric neighbourhoods. They are not equivalent as global spacetimes; for instance if we take the time orientation in accordance with the $`t`$ coordinate, the canonical (CTF) of the first spacetime is degenerate, constant equal to $`\mathrm{}`$, while $`t`$ is the (CTF) of the second one. We have a well defined cone angle $`2\pi \alpha `$ along such a universe line, which corresponds to a spacetime curvature concentrated along the line. In accordance with \[D-J- ‘T H\], \[‘T H\], if we normalize the gravitational constant to be $`G=1`$, the mass of the particle is related to the cone angle by $`m=(1/4)(1\alpha )`$; in $`(2+1)`$-gravity there are not physical constraints on the sign of $`Gm`$, so that an arbitrarily big $`\alpha `$ is allowed. Spacetimes with gravitating particles. A marked globally hyperbolic spacetime (coupled to massive particles) of type $$\delta =(g,[\alpha ]_r)=(g,(\alpha _1,\mathrm{},\alpha _r))$$ is an homeomorphism $$\varphi :(F_g\times \text{},\{p_i\}\times \text{})(M,L_i)$$ such that $`M^{}=M\{L_i\}`$ is an oriented and time-oriented globally hyperbolic flat Lorentzian 3-manifold (i.e. a $`(\text{𝕄}^{2+1},Isom^+(\text{𝕄}^{2+1})`$-manifold, where $`\text{𝕄}^{2+1}`$ is the standard Minkowski space) and each point of $`L_i`$ has a neighbourhood isometric to the above local models, with cone angle $`2\pi \alpha _i`$. It is convenient to restrinct to Geroch marking, that is we stipulate that the surfaces $`\varphi (F_g\times \{t\})`$ are future Cauchy surfaces. As usually we work up to Teichmüller equivalence and we denote by $`T_\delta ^{GR}`$ the corresponding Teichmüller space for a given type. To make it more meaningful it is convenient to restrict to maximal spacetimes. Identifying $`F_g`$ with $`F_g\times \{0\}`$ we have also the holonomy representation $$\rho _{[\varphi ]}:\pi (F_g^{})Isom^+(\text{𝕄}^{2+1}).$$ We also make the current assumption that the linear part of the holonomy takes values in $`SO^+(2,1)`$, the group of Lorentz transformations keeping the upper half-space invariant. Minkowskian suspensions of geometric cone surfaces. They are peculiar spacetimes which actually are $`(Y,G(Y))`$-manifolds, for suitably chosen open subsets $`Y`$ of $`\text{𝕄}^{2+1}`$, $`G(Y)`$ being the group of Minkowskian isometries keeping $`Y`$ invariant. As $`Y`$ we will take: $$Y_E=\text{𝕄}^{2+1}$$ with metric $`(dx^1)^2+(dx^2)^2(dx^3)^2`$, and thought fibred by the planes $`\{x^3=a\}`$. $$Y_H=\{x\text{𝕄}^{2+1}:(x^1)^2+(x^2)^2(x^3)^2<0,x^3>0\}$$ thought fibred by the surfaces $$\{x\text{𝕄}^{2+1}:(x^1)^2+(x^2)^2(x^3)^2=a^2,x^3>0\}$$ finally $$Y_T=\{x\text{𝕄}^{2+1}:(x^1)^2(x^3)^2<0,x^3>0\}$$ thought fibred by the surfaces $$\{x\text{𝕄}^{2+1}:(x^1)^2(x^3)^2=a^2,x^3>0\}.$$ By the change of coordinates $$x^1=\tau sh(u),x^2=y,x^3=\tau ch(u)$$ we see that $`Y_T`$ is isometric to $$P=\{(u,y,\tau )\text{}^{2+1}:\tau >0\},\mathrm{with}\mathrm{metric}\tau ^2du^2+dy^2d\tau ^2$$ and $`P`$ is fibred by the level planes of $`\tau `$. Each $`Y_{}`$ is oriented and time-oriented in the usual way. If $`S`$ is a flat cone surface of type $`(\text{}^2,g,[\alpha ]_r)`$, its Minkowskian suspension $`M(S)`$ is the obviously associated $`(Y_E,G(Y_E))`$-spacetime of type $`(g,[\alpha ]_r)`$, with holonomy equal to the holonomy of $`S`$. It is fibred by parallel copies of $`S`$. The canonical (CTF) is degenerate, costant equal to $`\mathrm{}`$. These are called static Minkowskian suspensions. If $`S`$ is a hyperbolic cone surface of type $`(\text{}^2,g,[\alpha ]_r)`$, its Minkowskian suspension $`M(S)`$ is the associated $`(Y_H,G(Y_H))`$-spacetime of type $`(g,,[\alpha ]_r)`$, with holonomy equal to the holonomy of $`S`$. It is fibred by parallel rescaled copies of $`S`$; these surfaces are the level surfaces $`S_a`$ ($`S=S_1`$) of the canonical (CTF); out of the particles they have constant mean curvature $`1/a`$ and constant intrinsic curvature equal to $`1/a^2`$. The initial singularity consists of one point. These suspensions are particularly nice when $`S`$ is an orbifold (and the matter-free spacetimes are particular cases); if the orbifold $`S=\text{𝕏}/\mathrm{\Gamma }`$, $`\mathrm{\Gamma }`$ acts isometrically also on the corresponding $`Y_{}`$, and $`M(S)=Y_{}/\mathrm{\Gamma }`$. The parameter space of $`Y_E`$ or $`Y_H`$-suspensions of a given type coincides, tautologically, with the parameter space of the suspended geometric cone surfaces (see the previous section). The $`Y_T`$-Minkowskian suspensions involve the special flat cone surfaces given by the meromorphic quadratic differentials with at most simple poles on Riemann surfaces. In fact each such a suspension is determined by a couple $`(F,q)`$, where $`F`$ is a Riemann surface and $`q`$ is a quadratic differential. That is, it is determined not only by the cone surface, but also by the horizontal and vertical measured foliations of the quadratic differential. We have already studied such spacetimes in \[B-G 2\] where we have shown how they “materialize” the classical Teichmüller flow. See also \[B-G 1\] for a description of the (CTF). In fact in \[B-G 2\] we considered only holomorphic quadratic differentials, but everything runs verbatim if one allows also simple poles. Recall that in this way one can realize all the types with $`2\pi \alpha _i=n_i\pi ,n_i1`$, satisfying the Gauss-Bonnet equality, with four exceptions (see \[M-S\]). Moreover, for any given realizable type, one knows the degrees of freedom (see \[V\]): if $`\mu (a)`$ denotes the number of cone points of cone angle $`a`$, then the degrees of freedom are $$2g+\mu (a)+(ϵ3)/2$$ where $`ϵ=1`$ iff there is at least one cone angle with odd $`n_i`$, and it is equal to $`1`$ otherwise. For example, when the type contains only $`n_i=3`$ (this corresponds to holomorphic quadratic differentials with simple zeros), the dimension of the corresponding space of $`Y_T`$-suspensions is $`6g6`$. The only orbifolds which produce such a kind of suspension are the orbifolds of type $`(\text{}^2,0,(1/2,1/2,1/2,1/2))`$. They are obtained by the natural identification of the edges of two copies of a same Euclidean rectangle. The corresponding group $`\mathrm{\Gamma }`$ is generated by two orthogonal translations and the rotation of angle $`\pi `$. Groups that determine the same $`Y_E`$-suspension (up to equivalence), do determine in general different $`Y_T`$-suspensions; in fact if we look at these groups acting on $`P`$, the horizontal and vertical foliations on each $`\tau `$-level plane induce different foliations on the (CTF)-level surfaces of the two suspensions. ## 4 Distinguished Deformations of Hyperbolic Suspensions. In this section we will refer heavily to \[B-G 1\] and to \[Me\]. All matter-free spacetimes with space of genus $`g2`$ are obtained by specific “deformations” of Minkowskian suspensions $`M(S)=Y_H/\mathrm{\Gamma }`$, and each such a deformation $`M(S,)`$ is uniquely determined by a measured geodesic lamination $``$ on $`S`$. $`M(S)`$ and $`M(S,)`$ have holonomies with the same linear part. The lifted lamination $`^{}`$ to the universal covering $`S^{}=\text{}^2`$ is “dual” to a real tree which is isometric to the initial singularity of $`M(S,)^{}`$. The initial singularity of $`M(S,)`$ must be properly interpreted in terms of a natural action of the fundamental group $`\pi (S)`$ on this real tree; the natural actions of $`\pi (S)`$ on the universal covering of the level surfaces of the (CTF) of $`M(S,)`$, asymptotically degenerate to that action on the real tree. If $`S`$ is a hyperbolic cone surface of type $`\delta =(g,[\alpha ]_r)`$ (we have omitted the “$`\text{}^2`$” in $`\delta `$), and $``$ is a measured geodesic lamination with compact support in $`S^{}`$ we can repeat those constructions (working with $`S^{}`$ which is now a cone hyperbolic surface) getting, by definition, a distinguished deformation $`M(S,)`$ of $`M(S)`$, which is again a spacetime of type $`\delta `$. The simplest deformations arise when $``$ is a multicurve, i.e. it consists of the finite union of disjoint simple geodesics endowed with positive weights. Assume, for simplicity, that there is one geodesic $`\sigma `$, with weight $`s`$ and length $`r`$. Consider the quotient $`A^{}(s,r)`$ of $`B^{}(s,r)=\{(u,y,\tau )P;0ys\}`$ by the group generated by the translation $`(u,y,\tau )(u+r,y,\tau )`$. Actually it is better to consider the isometric quotient $`A(s,r)`$ of $`B(s,r)Y_T`$, obtained via the explicit change of coordinates given in section 3. Then, to construct $`M(S,)`$, cut-open $`M(S)`$ along the suspension of $`\sigma `$ and insert $`A(s,r)`$ in the natural way. $`M(S,)`$ is, by construction, fibred by $`C^1`$-embedded space-like surfaces (made by the union of pieces of constant negative curvature and flat annuli); in fact they are the level surfaces of the canonical (CTF). The above construction is very simple nevertheless, as multicurves are dense in the space of measured geodesic laminations, by making the multicurve “complicated” enough, we can fairly well approximate the shape of any general distinguished deformation. Given a hyperbolic type $`\delta =(g,[\alpha ]_r)`$, we denote by $`D(\delta )`$ the subset of $`T_\delta ^{GR}`$ determined by the distinguished deformations of Minkowskian suspensions of hyperbolic cone surface of type $`\delta `$. Of course, a suspension is meant as the trivial deformation of itself and there is a natural projection $`p:D(\delta )T_\delta `$. The following proposition gives partial information on $`D(\delta )`$. We will use some notations introduced in section 2. The set of hyperbolic “$`(g,r)`$-types” can be identified with an open set of $`\text{}^r`$. ###### Proposition 4.1 (1) For each hyperbolic type $`\delta `$ there is an open (possibly empty) maximal subset $`𝒰_\delta `$ of $`T_\delta `$ such that $`p^1(𝒰_\delta )D(\delta )`$ is homeomorphic to $`𝒰_\delta \times \text{}^{6g6+2r}`$ (and $`p`$ becomes the natural projection onto the first factor). (2) For each $`(g,r)`$ there is a maximal non empty open subset $`𝒲_{(g,r)}`$ of the space of $`(g,r)`$-types, such that for each $`\delta 𝒲_{(g,r)}`$, $`𝒰_\delta `$ is non empty. (3) For any $`\delta `$, $$12g12+4rdimD(\delta )6g6+2r.$$ (4) If $`𝒰_\delta `$ is non empty, then $`𝒰_\delta \times \text{}^{6g6+2r}`$ is an open subset of $`T_\delta ^{GR}`$. Let us give a scketch of a proof. By using the result of section 2, the first statement is equivalent to show that the space of measured geodesic laminations with compact support on $`S^{}`$, for a given hyperbolic cone surface $`S`$ in $`𝒰`$ (for a suitable $`𝒰`$), is homeomorphic to $`\text{}^{6g6+2r}`$. This fact is known in the “limit” case when each $`\alpha _i=0`$, that is when $`S^{}`$ is a complete finite area hyperbolic surface with $`r`$ cusps (see \[Pe\]). Let us denote by $`HT_g^r`$ (which is homeomorphic to $`T_g^r`$ ) the Teichmüller space of such hyperbolic surfaces with $`r`$ cusps and fix one surface $`F`$. It is known that each geodesic lamination with compact support on $`F`$ has support contained in $`F^{\prime \prime }`$ obtained by removing from $`F`$ all the horocycles of length $`<1`$ around all the cusp points (see \[Pe\] pag. 72). It turns out that any hyperbolic cone surface $`S`$ which is “geometrically” close to $`F`$ has, up to homeomorphism, the same space of measured geodesic laminations with compact support on $`S^{}`$. The crucial fact is that if $`S`$ is close enough to a cusped $`F`$, each isotopy class of essential (i.e. non contractible nor contractible to one cone point) simple curves on $`S^{}`$ has a simple geodesic (in $`S^{}`$) shortest length representative. For some notions on the “geometric topology” see, for instance, chapter E of \[B-P\]. In our situation “$`S`$ geometrically close to $`F`$” roughly means that, removing suitable “round” disks with centres at the cone points of $`S`$, we find $`S^{\prime \prime }`$ which is quasi-isometric to $`F^{\prime \prime }`$, by a quasi-isometry close to an isometry. It follows that for any fixed compact subset $`K`$ of $`HT_g^r`$ there is an open subset $`U_K`$ (possibly empty) of $`T_\delta `$, which satisfies the first statement of the proposition. To prove the second statement, it is enough to show that, for any fixed $`F`$ as before, there are cone surface $`S`$ close to $`F`$ in the above sense. Fix a geodesic ideal triangulation $`𝒯`$ of $`F`$ (i.e. a “relaxed” triangulation of $`F`$ by ideal hyperbolic triangles). For each $`0<a<1`$ consider the horocycles of length $`a`$ around the cusps of $`F`$. Associate to each edge of the triangulation the length of the subarc determined by the horocycles. Consider the cone surface $`S`$ obtained accordingly with the construction after proposition 2.2, by using the same $`𝒯`$ as topological ideal triangulation of $`F_g^{}`$ and those lengths as edge-lengths. If $`a`$ is small enough, $`S`$ is close to $`F`$. $`S`$ is not close enough to a cusped $`F`$ when, at a very qualitative level, the masses are not big enough or the particles are too close each other on a given level surface of the (CTF) of the corresponding Minkowskian suspension. In such a case the basic trouble consists in the fact that the shortest length representative (if any) of an essential isotopy class of simple curves on $`S^{}`$ might be a broken geodesic passing through some cone points or not even a simple curve. The third statement is clear from the above discussion. To achieve the last statement it is enough to show that $`T_\delta ^{GR}`$ is of dimension $`12g12+4r`$; we are going to argue it without any assumption on the spacetime type $`\delta =(g,[\alpha ]_r)`$. The degrees of freedom of $`T_\delta ^{GR}`$. Fix a marked spacetime $`M`$ of type $`\delta `$ and a relatively compact globally hyperbolic open neigbourhood $`U`$ of the Cauchy surface image of $`F_g\times \{0\}`$. Let $`\rho :\pi (F_g^{})ISO^+(2,1)`$ be its holonomy. As $`\pi (F_g^{})`$ is a free group, a deformation of $`\rho `$ is simply obtained by modifying $`\rho `$ on a set of $`2g2+r+1`$ free generators. If a deformation $`\rho ^{}`$ is small enough, then, by the stability property of holonomies, $`\rho ^{}`$ is still the holonomy of a spacetime structure on the interior of $`U`$, with $`r`$ gravitating particles. So, as the holonomy is defined only up to conjugation, the dimension of the set of all these spacetimes “close” to $`M`$ is $`12g12+6r`$. In order to impose that the spacetimes have the specific cone angles prescribed by $`\delta `$, we have to impose $`2r`$ (that is $`(6d)r`$, where $`d`$ is the dimension of the conjugation orbit of a “rotation”) more independent conditions, and we finally get the required number of degrees of freedom $`12g12+4r`$. ## 5 Patchworking of Minkowskian Suspensions. A simple variation of the costruction of the distinguished modification of hyperbolic suspensions, based on multicurves, that we have described in the previous section, will produce interesting new examples of spacetimes. Let $`M(S,)`$ be as in the previous section. Assume that we have a finite union of simple closed geodesic on $`S^{}`$ disjoint from $``$. For simplicity, assume that there is a single geodesic $`\sigma `$ of length $`a`$. Let $`(F,q)`$ be a Riemann surface with a meromorphic quadratic differential $`q`$, with at most simple poles. Let $`M(F,q)`$ be the corresponding $`Y_T`$\- Minkowskian suspension (see section 3). Assume that the q-horizontal foliation on $`F`$ contains a simple closed leaf $`c`$ of length $`a`$. Then we can construct new spacetimes as follows: cut-open $`M(S,)`$ along the suspension of $`\sigma `$ and $`M(F,q)`$ along the suspension of $`c`$; then glue, pairwise, pieces of $`M(S,)`$ with pieces of $`M(F,q)`$ along isometric boundary components in the natural way. Note that there is, in general, a finite number of possible combinations, and the resulting Lorentz manifolds may be not connected, so we can take each connected component as a new spacetime. Call $`M([S,,\sigma ],[F,q,c])`$ any spacetime obtained in this way. By construction, it is fibred by space-like surfaces (made by rescaled pieces of $`S`$ and by “stretched ” pieces of $`F`$) which actually are the level surfaces of the canonical (CTF) of $`M([S,,\sigma ],[F,q,c])`$. Note also that the construction can be iterated, starting from suitable $`M([S,,\sigma ],[F,q,c])`$; so one can produce a wide class of new examples. This patchworking is peculiar of spacetimes with gravitating particles; in fact if we formally apply it to matter-free spacetimes we get nothing else than distinguished deformations of hyperbolic suspensions. In particular, let us use as $`(F,q)`$ the orbifolds of type $$(\text{}^2,0,4[1/2])=(\text{}^2,0,(1/2,1/2,1/2,1/2))$$ with the horizontal and vertical foliations of $`q`$ (with $`4`$ simple poles) parallel to the edges of the “fundamental” rectangle. It is not hard to construct by the patchworking procedure hyperbolic types of the form $`\delta =(g,[\alpha ]_r)=(g,[\alpha ]_r^{}2h[1/2])`$, $`2h+r^{}=r`$. On the other hand, these new spacetimes do not belong to $`D(\delta )`$ because, for instance, the level surfaces af the canonical (CTF) are not isometric (there are cone points of cone angle $`\pi `$ with no isometric neighbourhoods). Other differences manifest themselves by studying the past asymptotic states of the respective (CTF). By small perturbation of the holonomy of these examples one could produce examples out of $`D(\delta ^{})`$ for any $`[\alpha ]_r^{}`$ close to $`[\alpha ]_r`$. ## 6 Final Questions and Considerations. We are going to conclude with some questions, problems and, sometimes, with a guess about them. (1) Is $`T_\delta ^{GR}`$ connected ? The answer could depend on the type. We guess that the above examples not belonging to $`D(\delta )`$, actually do not even belong to the same connected component of any element of $`D(\delta )`$. (2) Does any spacetime satisfy the Gauss-Bonnet constraint $`(1\alpha _i)22g`$ ? We guess that by suitable small perturbations of the holonomy of static Minkowskian suspensions (which satisfy the Gauss-Bonnet equality) one could obtain spacetimes with $`_i(1\alpha _i)<22g`$ . We note that all the examples of spacetime that we have produced starting from non static Minkowskian suspensions have the following property : each particle line of universe has a neighboourhood isometric to the set of points of spatial distance $`<bt`$, for some positive $`b`$, from the $`t`$-axis in the model $`\{(z,t),\|z\|<1,t>0\}`$ with metric $`d\sigma _{(H,\alpha )}^2`$ (see section 3). (3) Does the same property hold for any spacetime with tame \- see \[B-G 1\]\- canonical (CTF) with values onto $`(0,\mathrm{})`$? It would be interesting to find, if any, examples where the linear function $`bt`$ must be replaced by some positive function $`f(t)`$ going faster to $`0`$ when $`t0`$. (4) Find an intrinsic characterization of hyperbolic cone surfaces belonging to $`𝒰_\delta `$. One expects that it could be expressed in terms of inequalities involving the cone angle, the genus and the distances between the cone points. (5) Describe $`𝒲_{(g,r)}`$. In particular, does $`m(g,r)`$ exist, with $`1>m(g,r)>0`$, such that for any $`\delta 𝒲_{(g,r)}`$ and for any mass $`m_i`$ associated to $`\delta `$, one has $`m_i>m(g,r)`$ ? For example, beside the “rigid” case $`(g,r)=(0,3)`$, the very peculiar case $`(g,r)=(0,4)`$ has $`𝒲_{(0,4)}`$ which coincides with the whole space of $`(0,4)`$-types; moreover for each type $`\delta 𝒲_{(0,4)}`$, $`𝒰_\delta `$ coincides with $`T_\delta `$, so that $`D(\delta )`$ = $`T_\delta \times \text{}^2`$. On the other hand, we guess, for example, that for each $`(0,r)`$, $`r>4`$, the last question has negative answer. (6) Is $`D(\delta )`$ always of dimension $`>6g6+2r`$ ? In other words, one is asking if there are always non trivial distinguished deformations. We guess that when $`g2`$ and the masses are all positive, then $`D(\delta )`$ contains at least $`T_\delta \times \text{}^{6g6}`$; in other words one expects that there is at least the same “amount” of distinguished deformations of the matter-free case of the same genus. (7) Let $`C`$ be any closed subset of $`𝒰_\delta `$. Is $`p^1(C)𝒰_\delta \times \text{}^{6g6+2r}`$ closed in $`T_\delta ^{GR}`$ ? Finally we note that in several instances of the present paper we have seen how very natural perturbations of a given spacetime do not preserve the type (see for instance the constructions of section 2 or the argument at the end of section 4). It would suggest that the study of (2+1)-gravity (coupled to particles) “type by type”, or even “space-genus by space-genus”, could be misleading. Spacetimes would be considered “all toghether” and it becomes quite demanding to figure out the structure of the corresponding (infinite dimensional) parameter space. 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Harer Combinatorics of Train Tracks, Annals of Math. Studies 125, Princeton 1992. \[Sc\] P. Scott, Bull. London Math. Soc. 15 (1983),no 5, 401-487. \[T\] W.P. Thurston, Geometry and Topology of 3-manifolds, “The” Notes, Princeton, 1982. \[Tr\] M. Troyanov, L’Enseignement Math. t. 32 (1986), 79-94. \[V\] W. A. Veech, Annals of Math. 124 (1986), 441-530. \[W\] E. Witten, Nucl. Phys. B311 (1988) 46. Dipartimento di Matematica Università di Pisa Via F. Buonarroti, 2 I-56127 PISA (Italia) benedett@dm.unipi.it Dipartimento di Fisica Università di Pisa Via F. Buonarroti, 2 I-56127 PISA (Italia) guadagni@difi.unipi.it
warning/0004/cond-mat0004042.html	ar5iv	text	# Crossover from the chiral to the standard universality classes in the conductance of a quantum wire with random hopping only ## I Introduction Since the introduction of the scaling approach to the problem of Anderson localization, it is known that transport characteristics of a disordered metal are universal, provided the disorder is sufficiently weak, the temperature sufficiently low so that quantum coherence is maintained over large distances, and the interaction between electrons can be neglected. An example is the phenomenon of weak localization, a small deviation from Ohm’s law for the conductance of a weakly disordered metal, which is suppressed by the application of a time-reversal symmetry breaking magnetic field. Though small, the weak-localization correction is universal in the sense that it does not depend on the shape of the sample, nor on any other microscopic or macroscopic property other than its dimensionality and the presence or absence of time-reversal symmetry and spin-rotation invariance. Another example is the phenomenon of universal conductance fluctuations: The sample-to-sample fluctuations of the conductance of a disordered metal or semiconductor are of order $`e^2/h`$ with a prefactor that only depends on dimensionality and symmetry. Both the weak-localization correction and the universal conductance fluctuations are precursors of the true Anderson localization, where as a result of destructive interference of multiple scattered quantum mechanical waves the dirty metal turns into an insulator for sufficiently strong disorder, or, in one or two dimensions, for a sufficiently large sample size. The original paper by Anderson, and most of the effort devoted to the problem of Anderson localization since then, considers the case of a particle on a lattice with a random on-site potential (diagonal disorder) and non-random hopping amplitudes. In that case, one distinguishes three universality classes, corresponding to the presence or absence of time-reversal and of spin-rotation symmetry. These three classes, are called orthogonal, unitary, and symplectic, respectively. Here, we will refer to these as the three “standard” universality classes. The electronic localization problem was soon generalized to lattice models with randomness in the hopping amplitudes (off-diagonal disorder). (This type of randomness was previously known from the description of phonons and narrow-gap semiconductors.) The localization problem with off-diagonal disorder has received comparatively much less attention, although it has been known since the work of Dyson that random systems with off-diagonal disorder, but without diagonal disorder, can behave in a way dramatically different from that of systems with diagonal disorder only, or with both types of disorder. For instance, the average density of states (DoS) for a one-dimensional chain with random nearest-neighbor hopping was found to be singular at the center of the band, $`\epsilon =0`$. According to the Thouless formula, such a singular DoS implies that at $`\epsilon =0`$ the conductance distribution be anomalous as well. Wegner and Gade in Ref. (see also Refs. ) found a two-dimensional counterpart to the singular behavior of the average DoS within their analysis of a non-linear-$`\sigma `$ model with a sublattice symmetry. Interest in the effect of off-diagonal disorder has revived in the 90’s on two fronts. Motivated by quenched approximations to interacting theories such as the quantum Hall effect at half-filling or gauge approaches to high $`T_c`$ superconductivity, the random flux problem (a special case of off-diagonal disorder in which hopping amplitudes have a random phase only) has been extensively studied, although very little consensus on its localization properties has emerged. A second thrust of activity has been motivated by the close resemblance between the anomalies at zero energy induced by pure off-diagonal disorder in two dimensions and the nature of the plateau transitions in the integer quantum Hall effect (IQHE): Both models might share the property that all eigenstates are localized except at one special energy. The reason why the localization properties of the random hopping problem can depart from those of the standard problem of Anderson localization is the existence of an additional sublattice symmetry in systems with off-diagonal but without diagonal disorder: In that case, the lattice can be divided into two sublattices, such that the Hamiltonian changes sign under a transformation where the wavefunction changes sign on one sublattice, but not on the other. As a result, the spectrum is symmetric with respect to a reflection about $`\epsilon =0`$ (i.e., eigenvalues appear in pairs $`\pm \epsilon `$). The fact that the band center $`\epsilon =0`$ is a very special energy in the presence of the sublattice symmetry explains why anomalies in the DoS and the localization properties occur at precisely this value of the energy. When the energy moves away from zero, the effects of the sublattice symmetry on the spectrum and the wavefunctions decreases and a crossover to the standard behavior takes place. The sublattice symmetry is broken by the presence of on-site disorder, long-range hopping, or (in some cases) by periodic boundary conditions. Counterparts to this sublattice symmetry in other disordered systems or in quenched approximation to interacting problems are numerous. They occur in, e.g., the QCD Hamiltonian, random $`XY`$ spin chains, diffusion in random environments, supersymmetric quantum mechanics, non-Hermitean quantum mechanics, and two-dimensional disordered models in the continuum such as Dirac fermions with random vector potentials. Following previous works in this field, which adopted the nomenclature of QCD, we will refer to the sublattice symmetry as chiral symmetry and will restrict our attention to random hopping problems with this symmetry. One-dimensional disordered systems with chiral symmetry have been well-studied with all kinds of approaches and in various contexts (for references, see the previous paragraph), and despite a continuing confusion about semantics, their localization properties can be considered well-understood. For two-dimensional systems the situation is different (see Refs. and references therein). Reliable analytical and numerical results are notoriously hard to obtain, and no consensus has been reached to date, not even on some most elementary issues. In view of this controversy, it is particularly instructive to study the natural intermediate between one and two dimensions, the thick (or “quasi-one-dimensional”) disordered wire. On the one hand, it shares the existence of both a localized and a diffusive regime of quantum transport with two dimensional disordered systems, while on the other hand, it allows for a controllable analytic treatment, just like the truly one-dimensional system. Moreover, quasi-one-dimensional systems appear as a logical intermediate step in the finite-size scaling approach for numerical simulations in two and three dimensions. Localization properties at the band center of a quasi-one-dimensional quantum wire with off-diagonal disorder were investigated in several previous publications by the authors, together with Simons and Altland. In those works we derived a chiral counterpart to the so-called DMPK equation, a Fokker-Planck equation that governs the distribution of the transmission eigenvalues of a quantum wire without chiral symmetry. Solution of the chiral DMPK equation for lengths beyond the localization length of the standard DMPK equation showed that there is no exponential localization if the number $`N`$ of propagating channels is odd (including the one-dimensional case), while the conductance decays exponentially with length if $`N`$ is even. This parity effect is strikingly similar to the sensitivity of the low-energy sector of a single antiferromagnetic spin-$`N/2`$ chain to the parity of $`N`$, on the one hand, or to the sensitivity of the low-energy sector of $`N`$ coupled antiferromagnetic spin-1/2 chains to the parity of $`N`$, on the other hand. In the special case of the chiral Fokker-Planck equation without time reversal invariance (random phase quantum wire), it was possible to calculate exactly the crossover from the diffusive to the localized regime for all moments of the conductance and to verify the validity of the assumption of universality against a numerical simulation of the random flux problem. The numerical simulations also confirmed that sufficiently far away from the center of the band, transport is governed by the standard universality classes. A limitation of the approach relying on the Fokker-Planck equations for the transmission eigenvalues is that it cannot describe how the conductance distribution crosses over from the chiral to the standard universality class as $`\epsilon `$ is tuned away from zero. In the renormalization group language, each Fokker-Planck equation describes a fixed point corresponding to a case of pure symmetry and the fixed points by themselves cannot be used to infer how the scaling flows take place between them. One possibility to obtain information about the crossover energy and length scales below (above) which the physics is that of the chiral (standard) universality classes, is to study the DoS of a chiral quantum wire. However, unlike in the case of a one-dimensional wire, where the Thouless formula connects conductance and DoS, for a quasi-one-dimensional wire it is not possible to infer transport properties from the DoS. In this paper, we use an alternative approach, developed by one of us for the study of transmission through a random waveguide with absorption. Focusing on weak-localization corrections and universal conductance fluctuations, we compute how, in the diffusive regime, the conductance distribution of a quantum wire with random hopping crosses over from the chiral to the standard universality classes as the energy is tuned away from zero. We are not able to compute the crossover in the localized regime. Instead, for the localized regime, we consider the conductance distributions in the pure symmetry classes and compare to numerical simulations to establish the crossover scale and to verify the validity of our predictions. The paper is organized as follows. In Sec. II, we define our microscopic model and derive the symmetries of the scattering matrix in the presence of the chiral symmetry. We then explain the scaling approach in Sec. III. The localized regime is studied in Sec. IV. Our main results are presented in Sec. V, where we consider the crossover from the chiral to the standard universality classes in the diffusive regime. In Sec. VI we compare our theoretical predictions to a numerical simulation of a random hopping model on a square lattice. We conclude in Sec. VII. ## II Microscopic model and scattering matrix A convenient microscopic model that describes a single particle hopping randomly between two sublattices is defined. The symmetries obeyed by the scattering matrix associated to this Hamiltonian are derived. Assuming weak disorder, the microscopic model is approximated by a model defined in the continuum, for which the scattering matrix is explicitly constructed. ### A Microscopic lattice model with chiral symmetry In a general form, the Schrödinger equation for an $`N`$-chain system with random hopping between two sublattices and without on-site randomness reads $$\epsilon \mathrm{\Psi }^{}(m)=T_m^{}\mathrm{\Psi }^{}(m+1)+T_{m1}^{}\mathrm{\Psi }^{}(m1).$$ (1) For a spinless particle, $`\mathrm{\Psi }(m)`$ is the $`N`$-component wavefunction where the index $`m`$ labels the position along the chain. For a particle with spin-1/2, $`\mathrm{\Psi }(m)`$ is the $`N`$-component wavefunction made of spinors. In that case, the $`N\times N`$ hopping matrix $`T_m`$ consists of quaternions. The system, and the allowed hopping matrix elements are depicted in Fig. 1(a). Note that the case of a square lattice with nearest-neighbor random hopping is included in the general formula (1), see Fig. 1(b). In order to model transport, we consider a disordered region of finite length $`L=Ma`$, $`a`$ being the lattice constant, and attach ideal leads with hopping matrix $`T_m=𝟙_{}`$ on both ends, see Fig. 1(a). Following Refs. , we draw the hopping matrices $`T_m`$ with $`m`$ inside the disordered region from a distribution centered around the $`N\times N`$ unit matrix, $$T_m=\mathrm{exp}(\delta T_m).$$ (2) We distinguish three symmetry classes depending on the presence or absence of time-reversal and spin-rotation symmetry. For a spinless particle (or for a spin-1/2 particle in the presence of spin-rotation symmetry), the hopping matrix $`\delta T_m`$ is real (complex) if time-reversal symmetry is present (absent). These two cases are commonly referred to as the orthogonal and unitary symmetry class and are labeled by the symmetry index $`\beta =1`$ and $`2`$, respectively. The case of broken spin-rotation symmetry with time-reversal symmetry is denoted $`\beta =4`$ and is referred to as the symplectic class. When $`\beta =4`$, the elements of the $`N\times N`$ matrix $`\delta T_m`$ are real quaternions. The situation when both time-reversal symmetry and spin-rotation symmetry are broken reduces to the unitary class ($`\beta =2`$) and will not be considered separately in this paper. We further assume that $`\delta T_m`$ has a Gaussian distribution, with zero mean and with variance given by $`(\delta T_m)_{kl}^{}[(\delta T_m^{})_{k^{}l^{}}^{}]^{}`$ $`=`$ $`{\displaystyle \frac{2\beta a}{\gamma \mathrm{}}}\delta _{mm^{}}^{}`$ (4) $`\times \left(\delta _{kk^{}}^{}\delta _{ll^{}}^{}{\displaystyle \frac{1\eta }{N}}\delta _{kl}^{}\delta _{k^{}l^{}}^{}\right),`$ $`(\delta T_m)_{kl}^{}(\delta T_m^{})_{k^{}l^{}}^{}`$ $`=`$ $`{\displaystyle \frac{2\beta }{\beta }}(\delta T_m)_{kl}^{}[(\delta T_m^{})_{k^{}l^{}}^{}]^{},`$ (5) (6) where $$\gamma =\beta N+2\beta \frac{2(1\eta )}{N},$$ (7) and $`\mathrm{}`$ is the mean free path. (Why $`\mathrm{}`$ can be identified as the mean free path is explained below.) Here, the symbol $``$ denotes the operation of complex conjugation for $`\beta =1,2`$, whereas it denotes the operation of Hermitean conjugation for quaternions for $`\beta =4`$. We assume weak disorder, $`\mathrm{}a`$. The parameter $`\eta `$ governs the relative randomness of the determinant of $`T_m`$. (See Ref. for the reason for its introduction.) We have chosen the statistical distribution (6) for technical convenience; it allows for an exact solution of the transport problem. As a justification for this choice, we recall that the transport properties do not depend on details of the microscopic model as long as disorder is weak, $`\mathrm{}a`$, and the length $`L`$ of the system is much larger than the mean free path. All properties of the microscopic model are summarized in the two parameters $`\mathrm{}`$ and $`\eta `$. \[The proper value of the parameter $`\eta `$ depends on the details of the microscopic model under consideration. For instance, for the random flux model (which is a special case of a random hopping model), $`\eta =0`$, while $`\eta >0`$ in generic random hopping models.\] To emphasize this universality, we compare our final results to numerical simulations for nearest-neighbor random hopping on a square lattice, cf. Fig. 1(b). In the leads on the left (L) and right (R), the Schrödinger equation (1) at energy $`\epsilon `$ is solved by a sum of plane waves moving towards the disordered region (denoted by a subscript $`\mathrm{i}`$) and away from the sample (denoted by a subscript $`\mathrm{o}`$) (see Fig. 2), $`\mathrm{\Psi }_\epsilon ^\mathrm{L}(m)`$ $`=`$ $`\psi _\epsilon ^{\mathrm{iL}}e^{ikma}+\psi _\epsilon ^{\mathrm{oL}}e^{ikma},`$ (8) $`\mathrm{\Psi }_\epsilon ^\mathrm{R}(m)`$ $`=`$ $`\psi _\epsilon ^{\mathrm{iR}}e^{ikma}+\psi _\epsilon ^{\mathrm{oR}}e^{ikma}.`$ (9) Here $`0k\pi /a`$, $`\epsilon =2\mathrm{cos}ka`$, and $`\psi _\epsilon ^{\mathrm{iL}}`$ and $`\psi _\epsilon ^{\mathrm{iR}}`$ ($`\psi _\epsilon ^{\mathrm{oL}}`$ and $`\psi _\epsilon ^{\mathrm{oR}}`$) are $`N`$-components vectors containing the amplitudes of the incoming (outgoing) plane waves in the left and right leads, respectively. The amplitudes of the ingoing and outgoing waves are connected through the Schrödinger equation (1) in the disordered region. This relation is formulated in terms of the $`2N\times 2N`$ scattering matrix $`S_\epsilon `$, $$\left(\begin{array}{c}\psi _\epsilon ^{\mathrm{oL}}\\ \psi _\epsilon ^{\mathrm{oR}}\end{array}\right)=S_\epsilon ^{}\left(\begin{array}{c}\psi _\epsilon ^{\mathrm{iL}}\\ \psi _\epsilon ^{\mathrm{iR}}\end{array}\right).$$ (10) Current conservation implies $`S^{}S^{}=𝟙_\mathrm{𝟚}.`$ (11) (Here and below we suppress the index $`\epsilon `$ if only scattering matrices at the same energy are involved.) For the cases $`\beta =1,4`$, i.e., if time-reversal symmetry is present, the complex conjugate of any eigenfunction is itself an eigenfunction with the same energy. (For $`\beta =4`$, complex conjugation is meant in the quaternion sense.) Since outgoing and incoming plane waves are interchanged under complex conjugation, we infer that time-reversal invariance is represented by the additional constraint $$S^{}S=𝟙_\mathrm{𝟚}.$$ (12) The Schrödinger equation (1) has an additional symmetry: the Hamiltonian changes sign under the transformation $`\mathrm{\Psi }(m)()^m\mathrm{\Psi }(m)`$. Correspondingly, for any realization of the disorder, the spectrum of energy eigenvalues is symmetric about the band center $`\epsilon =0`$. This symmetry, which originates from the fact that the disorder preserves the bipartite structure of the lattice, is referred to as chiral symmetry. The chiral symmetry is a special attribute of random hopping between different sublattices; it is broken by e.g. on-site randomness or next-nearest-neighbor hopping. It is the chiral symmetry that is responsible for the anomalous transport properties at the special energy $`\epsilon =0`$ of a quantum wire with random hopping. To find the effect of the chiral symmetry on the scattering matrix, we note that the transformation $`\mathrm{\Psi }(m)()^m\mathrm{\Psi }(m)`$ interchanges incoming waves at energy $`\epsilon `$ into outgoing waves at energy $`\epsilon `$, and vice versa. Applied to Eq. (10), this gives $`\left(\begin{array}{c}\psi _\epsilon ^{\mathrm{iL}}\\ \psi _\epsilon ^{\mathrm{iR}}\end{array}\right)`$ $`=S_\epsilon ^{}\left(\begin{array}{c}\psi _\epsilon ^{\mathrm{oL}}\\ \psi _\epsilon ^{\mathrm{oR}}\end{array}\right)=\left(S_\epsilon ^{}\right)^1\left(\begin{array}{c}\psi _\epsilon ^{\mathrm{oL}}\\ \psi _\epsilon ^{\mathrm{oR}}\end{array}\right).`$ (13) (The second equality follows from Eq. (10) at energy $`\epsilon `$.) Taken together with flux conservation (11), we thus find that the presence of the chiral symmetry results in the constraint $$S_\epsilon ^{}=\left(S_\epsilon ^{}\right)^{}$$ (14) for the scattering matrix $`S`$. Unlike the constraints of flux conservation and time-reversal symmetry, Eq. (14) involves scattering matrices at different energies. The exception is the band center $`\epsilon =0`$, where we find that $`S`$ is Hermitian, $$S_0^{}=S_0^{}.$$ (15) The scattering matrix is decomposed into four $`N\times N`$ subblocks $`r`$, $`r^{}`$ and $`t`$, $`t^{}`$, the reflection and transmission matrices, $$S=\left(\begin{array}{cc}r& t^{}\\ t& r^{}\end{array}\right).$$ (16) The transmission and reflection matrices determine the transport properties of the quantum wire. They are related to the conductance of the wire through the Landauer formula, $$G=\frac{2e^2}{h}\text{tr}t^{}t\frac{2e^2}{h}g,$$ (17) and to the shot noise power $$P=\frac{4e^3V}{h}\text{tr}\left[t^{}t(1t^{}t)\right]\frac{4e^3V}{h}p,$$ (18) $`V`$ being the applied voltage. (See Ref. for more applications to quantum transport.) A further decomposition of $`S`$ follows from the polar decomposition of the matrices $`r,r^{}`$ and $`t,t^{}`$, $$S=\left(\begin{array}{cc}𝒱^{}& 0\\ 0& 𝒰\end{array}\right)\left(\begin{array}{cc}\mathrm{tanh}X& (\mathrm{cosh}X)^1\\ (\mathrm{cosh}X)^1& \mathrm{tanh}X\end{array}\right)\left(\begin{array}{cc}𝒱& 0\\ 0& 𝒰^{}\end{array}\right),$$ (19) where $`𝒰`$, $`𝒰^{}`$, $`𝒱`$, and $`𝒱^{}`$ are $`N\times N`$ unitary matrices and $`X`$ is an $`N\times N`$ diagonal matrix with real numbers $`x_j`$ ($`j=1,\mathrm{},N`$) on the diagonal. In the presence of time-reversal symmetry, one has $`𝒰^{}𝒰^{}=𝒱^{}𝒱^{}=𝟙_{}^{}.`$ (20) Chiral symmetry implies a relationship between the unitary matrices $`𝒰`$, $`𝒰^{}`$, $`𝒱`$, and $`𝒱^{}`$ at opposite energies, $`𝒰_\epsilon ^{}=𝒰_\epsilon ^{},𝒱_\epsilon ^{}=𝒱_\epsilon ^{},X_\epsilon ^{}=X_\epsilon ^{}.`$ (21) In terms of the eigenvalues $`x_j`$, the equations (17) and (18) for the conductance and the shot noise power read $$g=\underset{j=1}{\overset{N}{}}\frac{1}{\mathrm{cosh}^2x_j},p=\underset{j=1}{\overset{N}{}}\frac{\mathrm{tanh}^2x_j}{\mathrm{cosh}^2x_j}.$$ (22) ### B Continuum model with chiral symmetry For weak disorder (mean free path $`\mathrm{}`$ much larger than the lattice spacing $`a`$), we may replace the lattice model (1) by a continuum model. We linearize the spectrum of the kinetic energy of Schrödinger equation (1) in the close vicinity of the band center $`\epsilon =0`$. Choosing a representation with left and right movers, we arrive at the continuum Schrödinger equation $$\epsilon \psi (y)=\left[i\sigma _3𝟙_{}_𝕪+\sigma _\mathrm{𝟛}𝕧(𝕪)+\sigma _\mathrm{𝟚}𝕨(𝕪)\right]\psi (y).$$ (23) Here $`\psi `$ is a $`2N`$ component vector (elements of $`\psi `$ occur in pairs that correspond to left and right movers), $`v`$ and $`w`$ are $`N\times N`$ Hermitean matrices, and the $`\sigma _\mu `$ $`(\mu =1,2,3)`$ are the Pauli matrices. In the presence of time-reversal symmetry $`w`$ ($`v`$) is (anti)symmetric. The continuum limit has been taken along the chains only; discreteness is maintained in the transverse direction through the $`N`$ components of $`\psi `$. The Fermi velocity has been set to one. The randomness in the hopping amplitudes has been translated to the matrices $`v`$ and $`w`$, by means of the identifications $`i\left(\delta T_m^{}\delta T_{m+1}^{}\right)+\text{ h.c.}`$ $``$ $`v(y),`$ (24) $`\left(\delta T_m^{}\delta T_{m+1}^{}\right)+\text{ h.c.}`$ $``$ $`w(y).`$ (25) With the choice (6), the disorder in $`v`$ is statistically independent from the disorder in $`w`$. Both $`v`$ and $`w`$ are Gaussian distributed with zero mean and with variance $`v_{ij}^{}(y)v_{kl}^{}(y^{})^{}`$ $`={\displaystyle \frac{\beta \delta (yy^{})}{\gamma \mathrm{}}}[\delta _{ik}^{}\delta _{jl}^{}{\displaystyle \frac{2\beta }{\beta }}\delta _{il}^{}\delta _{jk}^{}`$ (28) $`{\displaystyle \frac{2(\beta 1)(1\eta )}{\beta N}}\delta _{ij}^{}\delta _{kl}^{}],`$ $`w_{ij}^{}(y)w_{kl}^{}(y^{})^{}`$ $`={\displaystyle \frac{\beta \delta (yy^{})}{\gamma \mathrm{}}}[\delta _{ik}^{}\delta _{jl}^{}+{\displaystyle \frac{2\beta }{\beta }}\delta _{il}^{}\delta _{jk}^{}`$ (30) $`{\displaystyle \frac{2(1\eta )}{\beta N}}\delta _{ij}^{}\delta _{kl}^{}].`$ The symmetries (flux conservation, time-reversal, and chiral symmetry) of the scattering matrix in the continuum model are the same as for the lattice model. (Note that in the continuum model, the chiral transformation is represented by $`\psi \sigma _1\psi `$. The chiral symmetry then follows from the fact that $`\sigma _1`$ anticommutes with the Hamiltonian.) ## III Scaling approach The idea behind the scaling approach to the theory of localization in a quantum wire is to calculate how the scattering matrix $`S`$ of the quantum wire changes if a thin slice is added to the disordered region \[see Fig. 3(a)\]. Here we are mostly interested in the eigenvalues of the matrix product $`t^{}t=1r^{}r`$, i.e., in the parameters $`x_j`$ of the decomposition (19). Hence, it is sufficient to consider the reflection matrix $`r`$, and calculate how it is changed upon the addition of a thin slice. This change follows from the composition law $$r^{}=r_1^{}+t_1^{}\left(1r_2^{}r_1^{}\right)^1r_2^{}t_1^{},$$ (31) that gives the reflection matrix of two scatterers $`1`$ and $`2`$ in series, in terms of the reflection matrix of the right scatterer (2) and all reflection and transmission matrices of the left scatterer (1), see Fig. 3(b). If applied to a quantum wire, the only input in this approach is the statistical distribution of the transmission and reflection matrices $`t_1^{}`$, $`t_1^{}`$, $`r_1^{}`$, and $`r_1^{}`$ of the thin slice. The width $`\delta L`$ of the slice is taken much smaller than the mean free path $`\mathrm{}`$, so that the change of $`r`$ is small as well, although $`\delta L`$ must remain large compared to the lattice spacing $`a`$ for the continuum limit to be a good approximation. Then, the scattering matrix $`S_1`$ of the thin slice can be calculated in the second-order Born approximation from the Schrödinger Equation (23). The result is $`r_1^{}`$ $`=W+{\displaystyle \frac{i}{2}}[V,W],`$ (33) $`t_1^{}`$ $`=1+iV{\displaystyle \frac{1}{2}}V^2{\displaystyle \frac{1}{2}}W^2+i\epsilon \delta L,`$ (34) $`r_1^{}`$ $`=W+{\displaystyle \frac{i}{2}}[V,W],`$ (35) $`t_1^{}`$ $`=1iV{\displaystyle \frac{1}{2}}V^2{\displaystyle \frac{1}{2}}W^2+i\epsilon \delta L,`$ (36) where $$V=_0^{\delta L}𝑑yv(y),W=_0^{\delta L}𝑑yw(y).$$ Here we neglected terms that are of order $`(\delta L)^2`$. \[We also ignored the $`y`$-ordering of the integrals in Eq. (3) as it does not affect the statistical distribution of $`S_1`$ in view of the delta-function correlation of the random potentials $`v`$ and $`w`$.\] Using Eq. (II B) for the distribution of the random potentials $`v`$ and $`w`$, we find that the matrices $`V`$ and $`W`$ are Gaussian distributed with zero average and with variance proportional to the width $`\delta L`$ of the thin slice, $`V_{ij}^{}(V_{kl}^{})^{}`$ $`={\displaystyle \frac{\beta \delta L}{\gamma \mathrm{}}}[\delta _{ik}^{}\delta _{jl}^{}{\displaystyle \frac{2\beta }{\beta }}\delta _{il}^{}\delta _{jk}^{}`$ (39) $`{\displaystyle \frac{2(\beta 1)(1\eta )}{\beta N}}\delta _{ij}^{}\delta _{kl}^{}],`$ $`W_{ij}^{}(W_{kl}^{})^{}`$ $`={\displaystyle \frac{\beta \delta L}{\gamma \mathrm{}}}[\delta _{ik}^{}\delta _{jl}^{}+{\displaystyle \frac{2\beta }{\beta }}\delta _{il}^{}\delta _{jk}^{}`$ (41) $`{\displaystyle \frac{2\left(1\eta \right)}{\beta N}}\delta _{ij}^{}\delta _{kl}^{}].`$ Equations (31)–(III) define the scaling approach. They are exact for the continuum model (23) with the statistical distribution (II B) of the random potentials, which in turn was derived from the random hopping lattice model (1,6) in the limit of weak disorder. A different choice for the distribution of the hopping matrices in Eq. (6) would have led to different statistical properties of the scattering matrix for a thin slice. However, as we will verify in Sec. VI by numerical simulations, such differences are irrelevant in the sense of the renormalization group, i.e., they disappear for sufficiently long wires (longer than the mean free path $`\mathrm{}`$). Note that the reflection probability $`N^1\text{tr}r_1^{}r_1^{}`$ of a thin slice has average $$N^1\text{tr}r_1^{}r_1^{}=\delta L/\mathrm{},$$ (42) which justifies our choice that $`\mathrm{}`$ is the mean free path. In terms of the matrices $`V`$ and $`W`$, upon addition of the thin slice, the reflection matrix $`r`$ changes according to $$rr+\delta r,$$ (44) with $`\delta r`$ $`=`$ $`2i\epsilon \delta LrW+rWri(VrrV)+rWrWr`$ (46) $`{\displaystyle \frac{1}{2}}(W^2r+rW^2+V^2r+rV^2)+VrV.`$ We have not included terms of order $`VW`$ as their contributions vanish upon disorder averaging. Several observations can be made already on the level of the evolution equation (3), in combination with the Gaussian distribution (III) of the matrices $`V`$ and $`W`$. First, the distribution of $`r`$ is symmetric under a change of sign, $`rr`$. This implies that the average of any odd function of $`r`$ must be zero, for all values of the energy $`\epsilon `$. Second, at the band center $`\epsilon =0`$, the chiral symmetry implies that $`r`$ is Hermitian, cf. Eq. (15). The Hermiticity is broken by the first term in Eq. (46), which is proportional to the energy. Third, the distribution of $`r`$ is invariant under transformations $`rUrU^{}`$, where $`U`$ is an orthogonal (unitary) $`N\times N`$ matrix for $`\beta =1`$ ($`2`$). For zero energy, where $`r`$ is Hermitian, this implies that the distribution of $`r`$ depends on its eigenvalues $`\mathrm{tanh}x_j`$ only, cf. Eq. (19). As was shown in Refs. , in this case, the scaling flow can be represented in terms of a Fokker-Planck equation for the distribution $`P(x_1,\mathrm{},x_N;L)`$, $`{\displaystyle \frac{P}{L}}={\displaystyle \frac{1}{\gamma \mathrm{}}}{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \frac{}{x_i}}\left(\delta _{ij}^{}{\displaystyle \frac{1\eta }{N}}\right)J{\displaystyle \frac{}{x_j}}J^1P,`$ (47) $`J={\displaystyle \underset{k<l}{}}\|\mathrm{sinh}(x_lx_k)\|^\beta .`$ (48) Away from the center of the band, $`r`$ is no longer Hermitian, and its distribution depends on both eigenvalues and eigenvectors. However, for $`\epsilon `$ sufficiently far away from $`0`$ (this notion will be made precise below), the chiral symmetry has no effect on the scattering matrix, and $`P(x_1,\mathrm{},x_N;L)`$ obeys the Fokker-Planck equation for the standard orthogonal, symplectic, or unitary symmetry classes, the so-called Dorokhov-Mello-Pereyra-Kumar (DMPK) equation, $`{\displaystyle \frac{P}{L}}={\displaystyle \frac{1}{2(\beta N+2\beta )\mathrm{}}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{}{x_j}}J{\displaystyle \frac{}{x_j}}J^1P,`$ (49) $`J={\displaystyle \underset{k}{}}\|\mathrm{sinh}2x_k\|{\displaystyle \underset{k<l}{}}\|\mathrm{sinh}^2x_l\mathrm{sinh}^2x_k\|^\beta .`$ (50) There is no parameter $`\eta `$ in the DMPK equation; the presence of the parameter $`\eta `$ is special for the case of chiral symmetry at the band center $`\epsilon =0`$. In the language of the Fokker-Planck equation (48), $`\eta `$ controls the relative strength of the diffusion of the center of mass $`\overline{x}=(x_1+\mathrm{}+x_N)/N`$ compared to that of the relative coordinates $`x_j\overline{x}`$. The most important difference between the Fokker-Planck equations (48) and (50) are the symmetries of the Jacobians $`J`$. In Eq. (48), i.e., at the band center $`\epsilon =0`$, $`J`$ is invariant under a simultaneous translation $`x_jx_j+\delta x`$ and under a simultaneous reflection $`x_jx_j`$ for all $`j`$. \[The translation invariance decouples the motion of the “center of mass” $`\overline{x}=(x_1+\mathrm{}+x_N)/N`$ from the relative coordinates $`x_j\overline{x}`$, and hence calls for the presence of the parameter $`\eta `$ in Eq. (48).\] In the standard DMPK equation (50), i.e., for energies $`\epsilon `$ far away from the band center, $`J`$ is invariant under a reflection $`x_jx_j`$ for each $`j`$ separately; there is no longer translation invariance. It is the absence of this “local” reflection symmetry at $`\epsilon =0`$ that is responsible for anomalies in transport properties at $`\epsilon =0`$. In the remainder of this paper, we describe these in more detail, focusing on the distribution of the conductance in the localized regime $`LN\mathrm{}`$ and on the quantum interference corrections to the conductance in the diffusive regime $`\mathrm{}LN\mathrm{}`$. For the localized regime, we use the Fokker-Planck equations (48) and (50) to compare the transport properties for $`\epsilon =0`$ and $`\epsilon `$ far away from $`0`$. (A comparison for the case of broken time-reversal symmetry only has already been given in Ref. .) In the diffusive regime we start from the evolution equation (3) directly, in order to include the $`\epsilon `$-dependence of the transport properties. Knowledge of the crossover as a function of $`\epsilon `$ will allow us to specify what is meant by “$`\epsilon `$ sufficiently far away from $`0`$”, and hence when the standard DMPK equation (50) replaces the special Fokker-Planck equation (48) in the random hopping problem. ## IV localized regime Differences between the conductance distribution at the band center $`\epsilon =0`$ and away from $`\epsilon =0`$ are most pronounced in the localized regime $`LN\mathrm{}`$. Away from the band center, the conductance decreases exponentially with length, as is the case in the standard orthogonal, symplectic, and unitary classes. At the band center, however, the exponential decrease of the conductance is only observed if the number of channels is even, while for an odd number of channels the conductance decreases only algebraically. Exact calculations for the moments of the conductance in the standard symmetry classes have been obtained for all $`\beta `$, while for the chiral symmetry classes governed by the Fokker-Planck equation (48) only exact results for $`\beta =2`$ and $`\eta =1`$ are known. While we do not know of a way to extend our exact analysis of Ref. to the cases of orthogonal and symplectic symmetries, it is still possible to extract the conductance distribution deep inside the localized regime $`LN\mathrm{}`$ using the approximation scheme of Refs. . This is done here. We are thus able to compare the average and variance of the conductance, and the average and variance of its logarithm at and away from the band center $`\epsilon =0`$ for the orthogonal, symplectic, and unitary symmetry classes for all values of $`\eta `$. Our starting point is the Fokker-Planck equation (48), which we rewrite in the form $$\frac{P}{L}=\frac{1}{\gamma \mathrm{}}\underset{i,j=1}{\overset{N}{}}\left(\delta _{ij}\frac{1\eta }{N}\right)\frac{}{x_i}\left(\frac{P}{x_j}+\beta P\frac{\mathrm{\Omega }}{x_j}\right)$$ (51) with the “potential” $$\mathrm{\Omega }=\underset{k=1}{\overset{N}{}}\underset{l=k+1}{\overset{N}{}}\mathrm{ln}\|\mathrm{sinh}(x_kx_l)\|.$$ (52) Equation (51) has the interpretation that as $`L`$ increases, fictitious particles with the coordinates $`x_j`$ perform a Brownian motion subject to the repulsive two-body potential $`\mathrm{\Omega }`$. Since $`\mathrm{\Omega }`$ has a hard core, we may assume that $`x_1<x_2<\mathrm{}<x_N`$ for all $`L`$. In fact, as a result of their repulsive interaction, the distances between the $`x_j`$’s will grow with increasing length, until eventually for sufficiently large $`L`$ $$x_1x_2\mathrm{}x_N.$$ (53) Then we may approximate $$\frac{\mathrm{\Omega }}{x_j}N+12j,$$ (54) and find that Eq. (51) is solved by a Gaussian distribution for the $`x_j`$, $`P(x_1,\mathrm{},x_N;L)`$ $`\mathrm{exp}\{{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \frac{\gamma \mathrm{}}{4L}}(x_i{\displaystyle \frac{L}{\xi _i}})`$ (56) $`\times \left[(𝟙_{}{\displaystyle \frac{\mathrm{𝟙}\eta }{}}𝔼_{})^1\right]_{ij}(x_j{\displaystyle \frac{L}{\xi _j}})\}.`$ Here, the $`N\times N`$ matrix $`E_N`$ has the entries $`(E_N)_{ij}=1`$, and the channel-dependent “localization length” $`\|\xi _j\|`$ reads $$\xi _j=\frac{\gamma \mathrm{}}{\beta (N+12j)}.$$ (57) For comparison, in the standard orthogonal and unitary symmetry classes, the probability distribution $`P(x_1,\mathrm{},x_N;L)`$ in the localized regime is also given by a Gaussian of the type (56), but with $`\eta =1`$, $`\gamma =2(\beta N+2\beta )`$, and $`\xi _j=(\beta N+2\beta )\mathrm{}/(1+\beta j\beta )`$. In the localized regime $`LN\mathrm{}`$ only the $`x_j`$ that are closest to $`0`$ contribute to the conductance, cf. Eq. (22). For even $`N`$, they are $`x_{N/2}`$ and $`x_{(N/2)+1}`$, both of which are an average distance $$x_{N/2}=x_{(N/2)+1}=\frac{L}{\xi },\xi =\frac{\gamma \mathrm{}}{\beta },$$ (58) away from zero. The length scale $`\xi `$ serves as the localization length for even $`N`$. For odd $`N`$, the conductance is determined by only one eigenvalue, $`x_{(N+1)/2}`$, which has zero average, $$x_{(N+1)/2}=0.$$ (59) The presence of the eigenvalue $`x_{(N+1)/2}`$ with zero average is responsible for the absence of exponential localization in this case. The average and variance of the conductance and the average and variance of its logarithm follow from the probability distribution (56). For even $`N`$ the results are, with an accuracy $`𝒪(L^0/\xi ^0)`$ for the logarithms displayed, $`\mathrm{ln}g`$ $`=`$ $`{\displaystyle \frac{\beta }{4}}\left(1{\displaystyle \frac{1\eta }{N}}\right)^{\frac{1}{2}}{\displaystyle \frac{L}{\xi }}{\displaystyle \frac{1}{2}}\mathrm{ln}\left({\displaystyle \frac{L}{\xi }}\right),`$ (61) $`\mathrm{ln}\text{var}g`$ $`=`$ $`\mathrm{ln}g,`$ (62) and $`\mathrm{ln}g`$ $`=`$ $`{\displaystyle \frac{2L}{\xi }}+2\sqrt{{\displaystyle \frac{2}{\beta \pi }}\left(12{\displaystyle \frac{1\eta }{N}}\right){\displaystyle \frac{L}{\xi }}},`$ (63) $`\text{var}\mathrm{ln}g`$ $`=`$ $`{\displaystyle \frac{4}{\beta }}\left[1+\left(1{\displaystyle \frac{2}{\pi }}\right)\left(12{\displaystyle \frac{1\eta }{N}}\right)\right]{\displaystyle \frac{L}{\xi }}.`$ (64) The latter result shows that, in the localized regime, the conductance distribution is well approximated by a log-normal distribution; unlike the average conductance $`g`$ itself, which has fluctuations that are much bigger than the average, its logarithm $`\mathrm{ln}g`$ provides a good characteristic of the ensemble. For odd $`N`$, there is no exponential localization. The conductance has a broad distribution, which is neither characterized by the (average of the) conductance nor its logarithm, $$P(g)\frac{\mathrm{exp}\left[\frac{\gamma \mathrm{}}{4L}\left(1\frac{1\eta }{N}\right)^1\text{arccosh}^2g^{\frac{1}{2}}\right]}{g\sqrt{1g}}.$$ (65) With this distribution and up to corrections of order $`L^0/\xi ^0`$, the average conductance decays algebraically, $`g`$ $`=\left({\displaystyle \frac{\beta }{\pi }}\right)^{\frac{1}{2}}\left(1{\displaystyle \frac{1\eta }{N}}\right)^{\frac{1}{2}}\left({\displaystyle \frac{\xi }{L}}\right)^{\frac{1}{2}},`$ (67) $`g^2`$ $`={\displaystyle \frac{2}{3}}g,`$ (68) while the average of its logarithm grows proportional to $`L^{1/2}`$ rather than $`L`$, $`\mathrm{ln}g`$ $`=4\sqrt{{\displaystyle \frac{1}{\beta \pi }}\left(1{\displaystyle \frac{1\eta }{N}}\right){\displaystyle \frac{L}{\xi }}},`$ (69) $`\mathrm{var}\mathrm{ln}g`$ $`={\displaystyle \frac{8}{\beta }}\left(1{\displaystyle \frac{2}{\pi }}\right)\left(1{\displaystyle \frac{1\eta }{N}}\right){\displaystyle \frac{L}{\xi }}.`$ (70) Away from the band center $`\epsilon =0`$, the conductance distribution follows from the standard DMPK equation (50). It is close to log-normal, with $`\mathrm{ln}g`$ $`=`$ $`{\displaystyle \frac{L}{2\xi _{\mathrm{st}}}}{\displaystyle \frac{3}{2}}\mathrm{ln}\left({\displaystyle \frac{L}{\xi _{\mathrm{st}}}}\right),`$ (72) $`\mathrm{ln}\text{var}g`$ $`=`$ $`\mathrm{ln}g,`$ (73) $`\mathrm{ln}g`$ $`=`$ $`{\displaystyle \frac{2L}{\xi _{\mathrm{st}}}},`$ (74) $`\text{var}\mathrm{ln}g`$ $`=`$ $`{\displaystyle \frac{4L}{\xi _{\mathrm{st}}}},`$ (75) up to an accuracy of $`𝒪(L^0/\xi ^0)`$. Here the localization length for the standard symmetry classes is given by $$\xi _{\mathrm{st}}=(\beta N+2\beta )\mathrm{}.$$ (76) The most striking difference in the conductance distribution appears for odd $`N`$, where the absence of exponential localization at $`\epsilon =0`$ is contrasted with the exponential decay of the conductance for $`\epsilon 0`$. However, also for an even number of channels, there is an important difference. At $`\epsilon =0`$, the localization length $`\xi N\mathrm{}`$ is $`\beta `$-independent for large $`N`$, cf. Eqs. (7) and (58), while the localization length $`\xi _{\mathrm{st}}\beta N\mathrm{}`$ away from the band center is proportional to $`\beta `$ for $`\epsilon 0`$, cf. Eq. (76). Hence, upon moving away from the band center, the localization length increases by a factor $`\beta `$. \[The mean free path does not depend on $`\epsilon `$, see Eq. (42).\] The absence of a $`\beta `$-dependence for the localization length at the band center may be related to the anomaly in the DoS for random hopping models at that energy. In Ref. , it was shown that in the absence of time-reversal symmetry the DoS $`\rho (\epsilon )`$ near zero energy has a pseudogap, $`\rho (\epsilon )\epsilon \|\mathrm{ln}\epsilon \|`$, while in the presence of time-reversal symmetry $`\rho `$ has a logarithmic divergence, $`\rho (\epsilon )\|\mathrm{ln}\epsilon \|`$. We conclude that, upon breaking time-reversal symmetry, the decrease in the DoS available for transport, cancels the suppression of destructive interference responsible for the increase of the localization length in the standard case. The average and variance of the conductance in the localized regime are dominated by rare events, where the smallest $`x_j`$ is close to zero (corresponding to a transmission coefficient close to unity). For wires without chiral symmetry, approximation of $`P(x_1,\mathrm{},x_N;L)`$ by a Gaussian similar to Eq. (56) fails for $`x_j`$ close to zero because it does not account for the repulsion between $`x_j`$ and its mirror image $`x_j`$, cf. Eq. (50). While it does not affect the leading $`𝒪(L)`$ behavior of $`\mathrm{ln}g`$ and $`\mathrm{ln}\text{var}g`$, this failure shows up in the subleading logarithmic terms in Eqs. (72,73) which are different from what one would have obtained from a Gaussian distribution for the $`x_j`$. \[The results quoted in Eqs. (72,73) above follow from an exact solution of the DMPK equation.\] In the presence of the chiral symmetry, there is no repulsion between $`x_j`$ and $`x_j`$, so that the approximation (56) remains valid for $`x_j`$ close to zero. In this respect, we remark that the logarithmic terms in Eq. (IV), which were obtained with the help of Eq. (56) indeed agree with the exact solution of Ref. for the case $`\beta =2`$. ## V Diffusive regime In the diffusive regime $`\mathrm{}LN\mathrm{}`$, the effects of quantum interference do not take such a dramatic form as in the localized regime. The typical conductance of any sample is given by the classical Ohm’s law, $`g=N\mathrm{}/L`$, and does not know of quantum mechanical phase coherence, the presence or absence of time-reversal symmetry, or, as we shall see below, the presence or absence of chiral symmetry. The role of quantum mechanics, and hence the role of the symmetries of the microscopic Hamiltonian in this regime is confined to small corrections to the average conductance and to its sample-to-sample fluctuations. In spite of their smallness, these corrections are of prime importance, as they are a universal signature of quantum phase coherence, their size being determined by the fundamental symmetries of the system only. They do not depend on microscopic properties of the quantum wire, nor on its macroscopic characteristics, such as mean free path, width, or length. The two corrections are referred to as “weak-localization” and “universal conductance fluctuations”. The former is a small correction $`\delta g`$ ($`\delta p`$) to the ensemble averaged (dimensionless) conductance $`g`$ (shot-noise $`p`$) that is suppressed if time-reversal symmetry is broken by a magnetic field. For a standard quantum wire, it reads $$\delta g=\frac{\beta 2}{3\beta }\left(\delta p=\frac{\beta 2}{45\beta }\right).$$ (77) Since it signals the first departure from Ohm’s law, the weak localization correction to the conductance is precursor to the exponential suppression of the conductance in the localized regime. The universal conductance fluctuations refer to the sample-to-sample fluctuations of the conductance, which have variance, $$\text{var}g=\frac{2}{15\beta }.$$ (78) The breaking of time-reversal (spin-rotation) symmetry reduces the conductance fluctuations by a universal factor of $`\sqrt{2}`$ ($`2`$). In this section we calculate those quantum corrections for the case of a quantum wire with random hopping. Our calculations are inspired by the approach of Mello and Stone, who have derived and solved scaling equations for the moments of the conductance in the standard universality classes from the DMPK equation in the limit of large $`N`$. We consider both the quantum corrections for the pure symmetry classes, corresponding to the Fokker-Planck equations (48) and (50) at $`\epsilon =0`$ and $`\epsilon `$ far away from $`0`$, respectively, and for the intermediate regime, where the crossover between the two symmetry classes takes place. Since in the latter case no Fokker-Planck equation for the transmission eigenvalues $`x_j`$ is available, a modification of the approach of Ref. is needed, which is based on the more fundamental scaling equation for the reflection matrix $`r`$, Eq. (3), rather than on a Fokker-Planck equation for the transmission eigenvalues $`x_j`$. Such a method was proposed by one of the authors in the context of the transmission through a random waveguide with absorption. Below, we adapt this method to the present case (Sec. V A), and present solutions for the chiral symmetry classes at the band center $`\epsilon =0`$ (Sec. V B) and for the crossover from the chiral symmetry classes to the standard universality classes as $`\epsilon `$ moves away from the band center $`\epsilon =0`$ (Sec. V C). ### A Scaling equations Although we are primarily interested in the statistics of the transmission matrix $`t`$, and in particular in the (dimensionless) conductance $`g=\text{tr}t^{}t`$ and shot noise power $`p=\text{tr}\left[t^{}t(1t^{}t)\right]`$, we find it more convenient to formulate our scaling equations in terms of the reflection matrix $`r`$. Once we know $`r`$, unitary of the scattering matrix allows us to find the transmission properties without much effort. Before we write down the most general scaling equation for a trace of an arbitrary product of the reflection matrix $`r`$ and its Hermitian conjugate, we would like to focus on the scaling equation for $`\text{tr}r^{}r`$ in order to demonstrate the method and the approximations involved. Addition of a thin disordered slice to a disordered wire causes a small change $`rr+\delta r`$ to the reflection matrix $`r`$, see Eq. (46). Hence, upon addition of this slice, the trace $`\text{tr}r^{}r`$ changes to $`\text{tr}(r+\delta r)^{}(r+\delta r)`$. Using Eq. (46) for $`\delta r`$, we thus find, up to $`𝒪(\delta L)`$ \[Recall that the variance of $`W`$ is of order $`\delta L`$, so keeping terms up to $`𝒪(\delta L)`$ means up to $`𝒪(W^2)`$\], $`\delta \mathrm{tr}r^{}r`$ $`=`$ $`\text{tr}[r(1r^{}r)W+W(1r^{}r)r^{}`$ (80) $`(1rr^{})W(1r^{}r)W+r(1r^{}r)WrW+Wr^{}W(1r^{}r)r^{}].`$ All terms that involve the disorder potential $`V`$ in Eq. (46) canceled due to the cyclicity of the trace. Next we perform a disorder average over $`W`$ and over the reflection matrix $`r`$ of the wire of length $`L`$. We thus find $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}{\displaystyle \frac{\delta \text{tr}r^{}r}{\delta L}}`$ $`=`$ $`[\text{tr}(1r^{}r)]^2\text{tr}r\text{tr}r(1r^{}r)\text{tr}r^{}\text{tr}r^{}(1rr^{})`$ (82) $`+{\displaystyle \frac{2\beta }{\beta }}\text{tr}(1r^{}r)(1r^{}rr^2r^2){\displaystyle \frac{2\left(1\eta \right)}{\beta N}}\text{tr}(1rr^{})(1r^{}rr^2r^2).`$ Finally, we take the limit $`\delta L\mathrm{}`$, and replace the finite differences on the l.h.s. of Eq. (82) by differentials. It is apparent that the scaling equation obeyed by $`\mathrm{tr}r^{}r`$ is not closed: On the r.h.s. traces and products of traces of up to four reflection matrices appear. Closure requires an infinite family of scaling equations, and cannot be achieved on the level of scaling equations for the moments, but only with the help of the Fokker-Planck equation for the transmission eigenvalues $`x_j`$ in the cases of pure symmetry. However, for lengths $`LN\mathrm{}`$ it is possible to decouple this infinite set, and to find a solution order by order in $`L/(N\mathrm{})`$. Formally, this decoupling scheme proceeds along the lines of a large-$`N`$ expansion: In addition to the explicit factors $`N`$ in Eq. (82), each trace contributes a factor $`N`$. Further, we assume that, to leading order in $`N`$, the average of a product of traces equals the product of the averages. As we will see below, corrections correspond to a (co)variance of traces, and are of order $`N^0`$. Similarly, if we have a product of $`n`$ traces, we can expand in cumulants, where an $`n`$th cumulant will turn out to be of relative size $`N^{2n}`$. Such a decoupling scheme is known to work for the case of the standard DMPK equation, and its consistency can be verified from the scaling equations for traces and products of traces that we derive in this section. Let us now see how the scaling equation for $`\text{tr}r^{}r`$ decouples in this large-$`N`$ decoupling scheme. Recalling that $`\gamma `$ is of order $`N`$, cf. Eq. (7), we thus find that the r.h.s. of Eq. (82) is of order $`N^2`$, i.e., $$\mathrm{}_L\mathrm{tr}r^{}r=N2\mathrm{tr}r^{}r+\frac{1}{N}\mathrm{tr}r^{}r^2+𝒪(N^0).$$ (83) Here we have used the fact that the average of the trace of an odd product of $`r`$’s and $`r^{}`$’s is zero, see our discussion below Eq. (3). The initial condition at $`L=0`$ corresponds to perfect transmission, i.e., $`\mathrm{tr}r^{}r=0`$. The solution is easily found, $$\mathrm{tr}r^{}r=\frac{Ns}{s+1}+𝒪(N^0),$$ (84) where $`s=L/\mathrm{}`$. This solution corresponds to Ohm’s law for the conductance $`g=N\text{tr}r^{}r`$, $$g=\frac{N}{s+1}+𝒪(N^0).$$ (85) To this order in $`N`$, the result is entirely classical. The average $`\mathrm{tr}r^{}r`$ (and hence $`g`$) does not depend on the energy $`\epsilon `$ nor on the presence or absence of time-reversal symmetry. The dependence on time-reversal symmetry shows up through the term proportional to $`\frac{2\beta }{\beta }`$ on the r.h.s. of Eq. (82), which is of order $`N`$. It is this term in the scaling equation that gives rise to the weak-localization correction to the conductance. The scaling equation for $`\mathrm{tr}r^{}r`$ does not contain an explicit energy dependence. Instead, the energy-dependence shows up through the appearance of the traces like $`\text{tr}r^2`$ or $`\text{tr}r^{}r^3`$ in Eq. (82) in the weak-localization correction. Such traces that contain different numbers of $`r`$’s and $`r^{}`$’s strongly depend on energy, as can be seen from the scaling equation of, e.g., $`\text{tr}r^2`$, $$\frac{\gamma \mathrm{}}{\beta }_L\mathrm{tr}r^2=[\text{tr}(1r^2)]^22\text{tr}r\text{tr}r(1r^2)+\frac{4i\epsilon \gamma \mathrm{}}{\beta }\mathrm{tr}r^2+\left[\frac{2\beta }{\beta }\frac{2(1\eta )}{\beta N}\right]\mathrm{tr}(1r^2)(13r^2).$$ (86) With the same decoupling scheme as before, we find a closed scaling equation for $`\mathrm{tr}r^2`$ up to $`𝒪(N)`$, $$\mathrm{}_L\mathrm{tr}r^2=N2(12i\epsilon \mathrm{})\mathrm{tr}r^2+\frac{1}{N}\mathrm{tr}r^2^2+𝒪(N^0),$$ (87) which has the solution $`\mathrm{tr}r^2`$ $`=`$ $`N\left\{12i\epsilon \mathrm{}+2\sqrt{\epsilon \mathrm{}(i+\epsilon \mathrm{})}\mathrm{cot}\left[2\sqrt{\epsilon \mathrm{}(i+\epsilon \mathrm{})}s\right]\right\}^1+𝒪(N^0).`$ (88) One verifies that for $`\epsilon 0`$, the average $`\text{tr}r^2`$ equals the average $`\text{tr}r^{}r`$ that we computed above, since for zero energy one has $`r=r^{}`$. One also verifies that for $`\epsilon \mathrm{}1`$ the average $`\text{tr}r^2`$ approaches zero, as is the case in the standard symmetry classes. We are now ready to discuss the scaling equations for the trace of the product of an arbitrary number of reflection matrices and for the product of such traces. Hereto we write $`r_0=r`$ and $`r_1=r^{}`$, and define $$R_{j_1\mathrm{}j_n}\mathrm{tr}r_{j_1}\mathrm{}r_{j_n},$$ (89) where the indices $`j_k`$ can take the values $`0`$ or $`1`$. We define the symbol $`R`$ without indices as $`R=N`$. We also define products of traces through the symbols $$Q_{𝐧_1\mathrm{}𝐧_m}^{}=R_{i_1^{(1)}\mathrm{}i_{n_1}^{(1)}}^{}\mathrm{}R_{i_1^{(m)}\mathrm{}i_{n_m}^{(m)}}^{},$$ (90) where $`𝐧_j`$ denotes the $`n`$-tuple $`i_1^{(j)},\mathrm{},i_{n_j}^{(j)}`$. Proceeding along the same lines as above, we then find that the scaling equation for a single trace is given by (see Fig. 4) $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_L`$ $`R_{j_1\mathrm{}j_n}={\displaystyle \frac{2i\epsilon \gamma \mathrm{}[\underset{k=1}{\overset{n}{}}(1)^{j_k}]n\gamma }{\beta }}R_{j_1\mathrm{}j_n}`$ (95) $`+{\displaystyle \underset{1kln}{}}R_{j_k\mathrm{}j_l}R_{j_l\mathrm{}j_nj_1\mathrm{}j_k}+{\displaystyle \frac{2\beta }{\beta }}R_{j_k\mathrm{}j_lj_k\mathrm{}j_1j_n\mathrm{}j_l}{\displaystyle \frac{2(1\eta )}{\beta N}}R_{j_k\mathrm{}j_lj_l\mathrm{}j_nj_1\mathrm{}j_k}`$ $`+{\displaystyle \underset{1k<ln}{}}R_{j_{k+1}\mathrm{}j_{l1}}R_{j_{l+1}\mathrm{}j_nj_1\mathrm{}j_{k1}}+{\displaystyle \frac{2\beta }{\beta }}R_{j_{k+1}\mathrm{}j_{l1}j_{k1}\mathrm{}j_1j_n\mathrm{}j_{l+1}}{\displaystyle \frac{2(1\eta )}{\beta N}}R_{j_{k+1}\mathrm{}j_{l1}j_{l+1}\mathrm{}j_nj_1\mathrm{}j_{k1}}`$ $`{\displaystyle \underset{1k<ln}{}}R_{j_k\mathrm{}j_{l1}}R_{j_{l+1}\mathrm{}j_nj_1\mathrm{}j_k}+{\displaystyle \frac{2\beta }{\beta }}R_{j_k\mathrm{}j_{l1}j_k\mathrm{}j_1j_n\mathrm{}j_{l+1}}{\displaystyle \frac{2(1\eta )}{\beta N}}R_{j_k\mathrm{}j_{l1}j_{l+1}\mathrm{}j_nj_1\mathrm{}j_k}`$ $`{\displaystyle \underset{1k<ln}{}}R_{j_{k+1}\mathrm{}j_l}R_{j_l\mathrm{}j_nj_1\mathrm{}j_{k1}}+{\displaystyle \frac{2\beta }{\beta }}R_{j_{k+1}\mathrm{}j_lj_{k1}\mathrm{}j_1j_n\mathrm{}j_l}{\displaystyle \frac{2(1\eta )}{\beta N}}R_{j_{k+1}\mathrm{}j_lj_l\mathrm{}j_nj_1\mathrm{}j_{k1}}.`$ Here, it is understood that $`j_{n+1}j_1`$, $`j_0j_n`$. Moreover for $`nl=k+1>1`$, $`R_{j_{k+1}\mathrm{}j_{l1}}\mathrm{tr}𝟙_{}=`$, $`R_{j_{k+1}\mathrm{}j_{l1}j_{k1}\mathrm{}j_1j_n\mathrm{}j_{l+1}}R_{j_{k1}\mathrm{}j_1j_n\mathrm{}j_{k+2}}`$, and $`R_{j_{k+1}\mathrm{}j_{l1}j_{l+1}\mathrm{}j_nj_1\mathrm{}j_{k1}}R_{j_{k+2}\mathrm{}j_nj_1\mathrm{}j_{k1}},`$ respectively, whereas when $`k=1`$ and $`l=n`$ $`R_{j_{l+1}\mathrm{}j_nj_1\mathrm{}j_{k1}}\mathrm{tr}𝟙_{}=,`$ $`R_{j_{k+1}\mathrm{}j_{l1}j_{k1}\mathrm{}j_1j_n\mathrm{}j_{l+1}}R_{j_{k+1}\mathrm{}j_{l1}j_{l+1}\mathrm{}j_nj_1\mathrm{}j_{k1}}R_{j_2\mathrm{}j_{n1}},`$ respectively. Note that there is a one-to-one correspondence between contributions involving a product of two traces, say, $$R_{j_k\mathrm{}j_l}R_{j_l\mathrm{}j_nj_1\mathrm{}j_k}\mathrm{tr}\left(r_{j_k}\mathrm{}r_{j_l}\right)\mathrm{tr}\left(r_{j_l}\mathrm{}r_{j_n}r_{j_1}\mathrm{}r_{j_k}\right),$$ and contributions arising in the presence of time reversal symmetry $$\frac{2\beta }{\beta }R_{j_k\mathrm{}j_lj_k\mathrm{}j_1j_n\mathrm{}j_l}\frac{2\beta }{\beta }\mathrm{tr}\left[\left(r_{j_k}\mathrm{}r_{j_l}\right)\left(r_{j_k}\mathrm{}r_{j_1}r_{j_n}\mathrm{}r_{j_l}\right)\right]=\frac{2\beta }{\beta }\mathrm{tr}\left[\left(r_{j_k}\mathrm{}r_{j_l}\right)\left(r_{j_l}\mathrm{}r_{j_n}r_{j_1}\mathrm{}r_{j_k}\right)^\mathrm{t}\right],$$ or due to the randomness in the determinant of the hopping matrices $$\frac{2(1\eta )}{\beta N}R_{j_k\mathrm{}j_lj_l\mathrm{}j_nj_1\mathrm{}j_k}\mathrm{tr}\left[\left(r_{j_k}\mathrm{}r_{j_l}\right)\left(r_{j_l}\mathrm{}r_{j_n}r_{j_1}\mathrm{}r_{j_k}\right)\right].$$ For products of traces, we find $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_LQ_{𝐧_1\mathrm{}𝐧_m}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \frac{\gamma \mathrm{}}{\beta }}Q_{𝐧_1\mathrm{}𝐧_{j1}𝐧_{j+1}\mathrm{}𝐧_m}_LR_{𝐧_j}+{\displaystyle \underset{1k<lm}{}}Q_{𝐧_1\mathrm{}𝐧_{k1}𝐧_{k+1}\mathrm{}𝐧_{l1}𝐧_{l+1}\mathrm{}𝐧_m}^{}F_{𝐧_k𝐧_l}^{},`$ (97) where $`(\gamma \mathrm{}/\beta )_LR_{𝐧_j}`$ is given by the r.h.s. of Eq. (95) with omission of the angular brackets for the disorder averaging, and where $`F_{\mathrm{𝐦𝐧}}^{}=_{k=1}^m_{l=1}^nf_{k,l}`$ with $`f_{k,l}`$ $`=R_{i_k\mathrm{}i_mi_1\mathrm{}i_kj_l\mathrm{}j_nj_1\mathrm{}j_l}^{}+{\displaystyle \frac{2\beta }{\beta }}R_{i_k\mathrm{}i_mi_1\mathrm{}i_kj_l\mathrm{}j_1j_n\mathrm{}j_l}^{}{\displaystyle \frac{2(1\eta )}{\beta N}}R_{i_k\mathrm{}i_mi_1\mathrm{}i_k}^{}R_{j_l\mathrm{}j_nj_1\mathrm{}j_l}^{}`$ (101) $`+R_{i_{k+1}\mathrm{}i_mi_1\mathrm{}i_{k1}j_{l+1}\mathrm{}j_nj_1\mathrm{}j_{l1}}^{}+{\displaystyle \frac{2\beta }{\beta }}R_{i_{k+1}\mathrm{}i_mi_1\mathrm{}i_{k1}j_{l1}\mathrm{}j_1j_n\mathrm{}j_{l+1}}^{}{\displaystyle \frac{2(1\eta )}{\beta N}}R_{i_{k+1}\mathrm{}i_mi_1\mathrm{}i_{k1}}^{}R_{j_{l+1}\mathrm{}j_nj_1\mathrm{}j_{l1}}^{}`$ $`R_{i_k\mathrm{}i_mi_1\mathrm{}i_kj_{l+1}\mathrm{}j_nj_1\mathrm{}j_{l1}}^{}{\displaystyle \frac{2\beta }{\beta }}R_{i_k\mathrm{}i_mi_1\mathrm{}i_kj_{l1}\mathrm{}j_1j_n\mathrm{}j_{l+1}}^{}+{\displaystyle \frac{2(1\eta )}{\beta N}}R_{i_k\mathrm{}i_mi_1\mathrm{}i_k}^{}R_{j_{l+1}\mathrm{}j_nj_1\mathrm{}j_{l1}}^{}`$ $`R_{i_{k+1}\mathrm{}i_mi_1\mathrm{}i_{k1}j_l\mathrm{}j_nj_1\mathrm{}j_l}^{}{\displaystyle \frac{2\beta }{\beta }}R_{i_{k+1}\mathrm{}i_mi_1\mathrm{}i_{k1}j_l\mathrm{}j_1j_n\mathrm{}j_l}^{}+{\displaystyle \frac{2(1\eta )}{\beta N}}R_{i_{k+1}\mathrm{}i_mi_1\mathrm{}i_{k1}}^{}R_{j_l\mathrm{}j_nj_1\mathrm{}j_l}^{}.`$ Here we denoted $`𝐦=i_1,\mathrm{},i_m`$ and $`𝐧=j_1,\mathrm{},j_n`$. Below, we are interested in averages and (co)variances of traces of an even number of reflection matrices up to order $`N^0`$. In both cases, the terms proportional to $`(1\eta )`$ do not play a role. For the average of a single trace, this is immediately clear from Eq. (95). To see this for the (co)variance of two traces, some further inspection of Eq. (V A) is needed. First, $`\eta `$ appears explicitly in the quantity $`F_{\mathrm{𝐦𝐧}}`$, multiplying a product of two traces, see Eq. (LABEL:eq:Fmn). A priori, the leading contribution, which is obtained by replacement of those traces by their averages, is of the same order \[$`𝒪(N)`$\] as the other terms in Eq. (LABEL:eq:Fmn). However, as $`m`$ and $`n`$ are even, each of the two traces multiplying $`(1\eta )`$ contains an odd number of reflection matrices, so that their averages vanish. Hence, to leading order in $`N`$, the contribution from the term proportional to $`(1\eta )`$ vanishes. Second, $`\eta `$ appears implicitly through the derivative $`(\gamma \mathrm{}/\beta )_LR_{𝐧_j}`$ in Eq. (97). Again, to leading order in $`N`$, its contribution vanishes, and one is left with a term of relative size $`N^2`$. It should be mentioned that Eqs. (95) and (V A) can be extended to the case in which a weak staggering of the hopping amplitude is present in the microscopic model, cf. Eqs. (17). (How to generalize the Fokker-Planck equation (48) to include dimerization was shown in Ref. , see also Ref. .) Weak staggering of the hopping amplitude is implemented by requiring that the disorder potential $`W`$ has the Gaussian distribution with variance (41) and mean $`W_{jk}=\frac{\beta \delta L}{\gamma \mathrm{}}\mathrm{\Delta }\delta _{jk}`$. Here $`\mathrm{\Delta }`$ measures the strength of the dimerization along the chain direction. With weak dimerization, Eq. (95), say, is modified by the addition on the r.h.s. of the contribution $`\mathrm{\Delta }_{k=1}^nR_{j_k\mathrm{}j_n\mathrm{}j_k}R_{j_{k+1}\mathrm{}j_n\mathrm{}j_{k1}}`$. We see that the scaling equations now couple traces over an even and odd number of reflection matrices as is expected since the probability distribution of $`W`$ is not anymore symmetric about $`W=0`$, cf. Eq. (46). Equations (95) and (V A) are the central results of this section. These equations are more general than the Fokker-Planck equations (48) and (50) in the sense that they are valid both at the center of the band $`\epsilon =0`$ and in its proximity. Their limitation is that they can only be solved in the diffusive regime $`LN\mathrm{}`$. In particular, they cannot be used to probe the localized regime (in contrast to their counterparts in the problem of a wave guide with absorption, see Ref. ). The next two subsections are devoted to a solution in the diffusive regime. The case of pure chiral symmetry ($`\epsilon =0`$) is considered in Sec. V B; the energy dependence of the solution is discussed in Sec. V C. ### B Diffusive regime in the chiral limit The general scaling equations (95) and (V A) simplify considerably at the band center $`\epsilon =0`$. At the band center, the scattering matrix is Hermitian, and hence $`r=r^{}`$. Restricting our attention to single traces and products of two traces, we find the scaling equations $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_L\text{tr}r^n`$ $`=`$ $`{\displaystyle \frac{n\gamma }{\beta }}\text{tr}r^n+{\displaystyle \frac{n}{2}}{\displaystyle \underset{k=0}{\overset{n}{}}}\text{tr}r^{nk+1}\text{tr}r^{k+1}+{\displaystyle \frac{n}{2}}{\displaystyle \underset{k=1}{\overset{n1}{}}}\left(\text{tr}r^{nk1}\text{tr}r^{k1}2\text{tr}r^{nk}\text{tr}r^k\right)`$ (104) $`+{\displaystyle \frac{n}{2}}\left[{\displaystyle \frac{2\beta }{\beta }}{\displaystyle \frac{2(1\eta )}{\beta N}}\right]\left[(n+1)\text{tr}r^{n+2}+(n1)\text{tr}r^{n2}2(n1)\text{tr}r^n\right],`$ $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_L\text{tr}r^m\text{tr}r^n`$ $`=`$ $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}\left[\text{tr}(_L\text{tr}r^m)r^n+\text{tr}r^m(_L\text{tr}r^n)\right]+{\displaystyle \frac{2mn}{\beta }}\text{tr}r^{m+n2}(1r^2)^2`$ (106) $`{\displaystyle \frac{2mn(1\eta )}{\beta N}}\text{tr}r^{m1}(1r^2)\text{tr}r^{n1}(1r^2).`$ Here $`(\gamma \mathrm{}/\beta )_L\text{tr}r^n`$ is the r.h.s. of Eq. (104) with the omission of the disorder averaging brackets. (If $`n=1`$, the last term in Eq. (104) should be omitted.) \[Alternatively, one could have used the Fokker-Planck equation (48) to derive these scaling equations. Both methods agree as we have verified explicitly.\] The average and variance of the conductance $`g=N\mathrm{tr}r^2`$ can be computed by straightforward solution of Eqs. (104) and (106) using the decoupling scheme of Sec. V A. The result is, up to corrections of order $`N^1`$, $`g`$ $`=`$ $`{\displaystyle \frac{N}{s+1}}+{\displaystyle \frac{2\beta }{\beta }}{\displaystyle \frac{s^2}{(1+s)^3}},`$ (108) $`\mathrm{var}g`$ $`=`$ $`{\displaystyle \frac{4}{15\beta }}\left[1{\displaystyle \frac{6s+1}{(s+1)^6}}\right],`$ (109) where, as before $`s=L/\mathrm{}`$. For the derivation of these results we needed the following intermediate results, $`\text{tr}r^4`$ $`=`$ $`{\displaystyle \frac{Ns^2(3s^2+8s+6)}{3(1+s)^4}},`$ $`\text{tr}r^6`$ $`=`$ $`{\displaystyle \frac{Ns^3(15s^4+82s^3+177s^2+180s+75)}{15(1+s)^7}},`$ up to corrections of order $`N^0`$ and $`\mathrm{tr}r\mathrm{tr}r`$ $`=`$ $`{\displaystyle \frac{2}{3\beta }}\left[1{\displaystyle \frac{1}{(s+1)^3}}\right],`$ $`\text{tr}r\text{tr}r^3`$ $`=`$ $`{\displaystyle \frac{2}{15\beta }}\left[4{\displaystyle \frac{5s^3+15s^2+24s+4}{(s+1)^6}}\right],`$ up to corrections of order $`N^1`$. In the diffusive regime $`\mathrm{}LN\mathrm{}`$ we observe that the variance of the conductance at $`\epsilon =0`$ is twice the value taken in the standard case, for $`\epsilon `$ far away from $`0`$, cf. Eq. (78). The result that the presence of the extra chiral symmetry leads to a doubling of conductance fluctuations was found previously for the random flux model (corresponding to our case $`\beta =2`$) from numerical simulations and from an exact solution of the Fokker-Planck equation (48). The factor two decrease in the fluctuations as the chiral symmetry is broken, is reminiscent of the factor two decrease of the conductance fluctuations upon breaking time-reversal symmetry or upon breaking a spatial symmetry. According to Eq. (108), application of a magnetic field has an effect on the average conductance, but this effect vanishes in the diffusive limit $`\mathrm{}LN\mathrm{}`$, i.e., $`s1`$. In other words, there is no weak-localization correction to the conductance in the diffusive regime. It is instructive to note a coincidence between the $`\beta `$-dependence of the average conductance $`g`$ and the $`\beta `$-dependence of the localization length $`\xi `$ which was considered in the previous section. In the case of the random hopping model at zero energy, there is no weak-localization correction, and the localization length $`\xi `$ does not depend on the presence or absence of time-reversal symmetry. On the other hand, without chiral symmetry (for large energies), the negative correction to the average conductance for $`\beta =1`$ foreshadows the localization transition, which occurs on a length scale $`\xi _{\mathrm{st}}`$ that is proportional to $`\beta `$, i.e., localization takes place twice as fast without than with a time-reversal symmetry breaking magnetic field. Absence of weak-localization correction to the conductance had been pointed out by Gade and Wegner in their study of a (two-dimensional) non-linear-$`\sigma `$ model implementing the chiral symmetry. (See also Ref. .) Finally, notice that there is no precursor in Eqs. (V B) of the even-odd effect seen in the localized regime. This agrees with the exact solution for $`\beta =2`$, where it was found that the even-odd effect is non-perturbative in the expansion parameter $`L/N\mathrm{}`$. To order $`N`$, the average conductance is the same as in the case of a wire without chiral symmetry. Differences show up only to order $`N^0`$, where we find that there are no weak-localization corrections for the chiral case. This is not a coincidence that is limited to the average of the conductance $`g=\text{tr}t^{}t`$. It extends to the averages of traces of arbitrary powers of $`r`$ or $`t`$. To see this and in order to allow for a more detailed comparison to the case where chiral symmetry is absent, we rephrase the scaling equation (104) in terms of the transmission matrix $`t`$. In the limit of large $`N`$, one thus obtains $`\left(N+{\displaystyle \frac{2\beta }{\beta }}\right)\mathrm{}_L\mathrm{tr}(t^{}t)^n`$ $`=`$ $`n{\displaystyle \underset{m=1}{\overset{n}{}}}\mathrm{tr}(t^{}t)^{n+1m}\mathrm{tr}(t^{}t)^m+n{\displaystyle \underset{m=1}{\overset{n1}{}}}\mathrm{tr}(t^{}t)^{nm}\mathrm{tr}(t^{}t)^m`$ (111) $`n(2n+1){\displaystyle \frac{2\beta }{\beta }}\mathrm{tr}(t^{}t)^{n+1}+2n^2{\displaystyle \frac{2\beta }{\beta }}\mathrm{tr}(t^{}t)^n+𝒪(N^0).`$ For a quantum wire without the chiral symmetry, the leading $`𝒪(N)`$ contribution to $`\mathrm{tr}(t^{}t)^n`$ is precisely the same as the first line of Eq. (111). The weak-localization correction proportional to $`\frac{2\beta }{\beta }`$ differs in the standard case from the second line of Eq. (111) as it reads $`n^2\frac{2\beta }{\beta }\mathrm{tr}(t^{}t)^n+n(n1)\frac{2\beta }{\beta }\mathrm{tr}(t^{}t)^{n+1}`$. Hence, whereas the solution of Eq. (111) is the same as for an ordinary quantum wire to leading order in $`N`$, $$\mathrm{tr}(t^{}t)^n=\frac{N}{2s}B(n,1/2),$$ (112) where $`B(x,y)=\mathrm{\Gamma }(x)\mathrm{\Gamma }(y)/\mathrm{\Gamma }(x+y)`$ is the beta function, the combination $`(2n+1)\mathrm{tr}(t^{}t)^{n+1}2n\mathrm{tr}(t^{}t)^n`$ on the second line of Eq. (111) conspires with the coefficient $`B(n,1/2)`$ to insure the disappearance of a weak-localization correction for all averages $`\mathrm{tr}(t^{}t)^n`$ in the presence of the chiral symmetry. As a corollary, we find that the average density of the transmission eigenvalues $`x_j`$ $$\rho (x)=\underset{j=1}{\overset{N}{}}\delta (xx_j)$$ (113) has no weak-localization correction as well. With these results, it is little work to compute the average shot noise power $`p=\text{tr}t^{}t(1t^{}t)`$ and its weak-localization correction, $`p`$ $`=`$ $`{\displaystyle \frac{N}{3}}\left[{\displaystyle \frac{1}{s+1}}{\displaystyle \frac{1}{(s+1)^4}}\right]`$ (115) $`+{\displaystyle \frac{2\beta }{\beta }}\left[{\displaystyle \frac{s^2}{3(1+s)^3}}{\displaystyle \frac{7s^2}{3(1+s)^6}}\right].`$ Just like in the case of the conductance there is no weak-localization correction in the diffusive regime $`\mathrm{}LN\mathrm{}`$. ### C Crossover between the chiral and standard universality classes For any nonzero energy $`\epsilon `$, the chiral symmetry of Eq. (15) is broken. Hence one expects that for a sufficiently long length $`L`$ of the quantum wire, its transmission properties will flow to those of the standard symmetry class. This flow is governed by a crossover length scale $`\mathrm{}_\epsilon `$ so that for $`L\mathrm{}_\epsilon `$, the transmission properties are still alike those in the chiral symmetry class, while for $`L\mathrm{}_\epsilon `$ they resemble those of the standard symmetry class. We distinguish three possible regimes where this crossover can take place: * The crossover takes place in the ballistic regime, $`\mathrm{}_\epsilon \mathrm{}`$, * The crossover takes place in the diffusive regime, $`\mathrm{}\mathrm{}_\epsilon N\mathrm{}`$, or * The crossover takes place in the localized regime, $`\mathrm{}_\epsilon N\mathrm{}`$. This regime cannot be treated with the methods used in the paper. For the case $`N=1`$ of a single-channel quantum wire, this regime has been studied in Refs. . The full set of scaling equations (95) and (V A) can be used to describe the first two regimes (and the intermediate region between them). Although the solution of the scaling equations is straightforward — within the large-$`N`$ decoupling scheme, the scaling equations are linear ordinary differential equations that can be solved one-by-one (see appendix A) — it is a quite cumbersome task, and many expressions get quite lengthy. \[The expression (88) for $`\text{tr}r^2`$ is the only example whose solution can be represented by a one-line equation.\] To simplify our presentation and to save the reader from those lengthy expressions, we focus on the regime $`\mathrm{}\mathrm{}_\epsilon N\mathrm{}`$, where the crossover takes place inside the regime of diffusive dynamics. The length scale for the crossover can be identified from Eq. (88) as $$\mathrm{}_\epsilon =\sqrt{\frac{\mathrm{}}{2\epsilon }}.$$ (116) (Here we have neglected $`\epsilon \mathrm{}`$ with respect to $`i`$ in $`\sqrt{i+\epsilon \mathrm{}}`$. This is consistent with our focus on the regime $`\mathrm{}_\epsilon \mathrm{}`$.) We note that $`\epsilon `$ is none but the Thouless energy for a diffusive process with diffusion constant $`v_f\mathrm{}`$ (having momentarily reinstated the Fermi velocity $`v_f`$) in a system of linear size $`\mathrm{}_\epsilon `$. Using the hierarchy of length scales $`\mathrm{}\mathrm{}_\epsilon `$ we then find that the solution of the scaling equations takes a relatively simple form. For the average and variance of the conductance $`g`$ we find up to $`𝒪(N^0)`$ $`g`$ $`=`$ $`{\displaystyle \frac{N\mathrm{}}{\sigma \mathrm{}_\epsilon }}{\displaystyle \frac{2\beta }{\beta }}\left[{\displaystyle \frac{1}{3}}{\displaystyle \frac{z\mathrm{coth}(z^{}\sigma )+z^{}\mathrm{coth}(z\sigma )}{4\sigma }}\right].`$ (117) $`\text{var}g`$ $`=`$ $`{\displaystyle \frac{2}{15\beta }}+{\displaystyle \frac{2}{\beta }}[{\displaystyle \frac{3z\sigma \mathrm{coth}(z\sigma )2}{16\sigma ^4}}+{\displaystyle \frac{i}{8\sigma ^2\mathrm{sinh}^2(z^{}\sigma )}}+\mathrm{c}.\mathrm{c}.].`$ (118) Here we defined $`z=1+i`$ and $`\sigma =L/\mathrm{}_\epsilon `$. For the average of the shot-noise power we find up to $`𝒪(N^0)`$ $`p={\displaystyle \frac{N\mathrm{}}{3\sigma \mathrm{}_\epsilon }}{\displaystyle \frac{2\beta }{\beta }}\{{\displaystyle \frac{1}{45}}+[{\displaystyle \frac{(3z2z^{}\sigma ^2)\mathrm{coth}(z\sigma )}{24\sigma ^3}}+{\displaystyle \frac{i}{4\sigma ^2\mathrm{sinh}^2(z^{}\sigma )}}+\mathrm{c}.\mathrm{c}.]\}.`$ (119) For the derivation of these results, we needed the following intermediate results, all up to corrections of order $`N^0`$, $`\mathrm{tr}(r^{}r)^3`$ $`=`$ $`N{\displaystyle \frac{N\mathrm{}}{\mathrm{}_\epsilon }}{\displaystyle \frac{23}{15\sigma }},`$ (120) $`\mathrm{tr}r^2`$ $`=`$ $`N{\displaystyle \frac{N\mathrm{}}{\mathrm{}_\epsilon }}z^{}\mathrm{coth}(z^{}\sigma ),`$ (121) $`\mathrm{tr}r^{}r^3`$ $`=`$ $`N{\displaystyle \frac{N\mathrm{}}{\mathrm{}_\epsilon }}\left[{\displaystyle \frac{(4\sigma ^2+i)z^{}\mathrm{coth}(z^{}\sigma )}{4\sigma ^2}}{\displaystyle \frac{1}{2\sigma \mathrm{sinh}^2(z^{}\sigma )}}\right],`$ (122) $`\mathrm{tr}r^2r^2`$ $`=`$ $`N{\displaystyle \frac{N\mathrm{}}{\mathrm{}_\epsilon }}\left[{\displaystyle \frac{4\sigma z^{}\mathrm{coth}(z\sigma )z\mathrm{coth}(z^{}\sigma )}{2\sigma ^2}}\right],`$ (123) $`\mathrm{tr}(r^{}r^2)^2`$ $`=`$ $`N{\displaystyle \frac{N\mathrm{}}{\mathrm{}_\epsilon }}\left[{\displaystyle \frac{(16\sigma ^4z^{}+8\sigma ^2z+z^{})\mathrm{coth}(z^{}\sigma )}{16\sigma ^4}}{\displaystyle \frac{z\mathrm{coth}(z\sigma )}{4\sigma ^4}}{\displaystyle \frac{8\sigma ^22z\sigma \mathrm{coth}(z^{}\sigma )+5i}{8\sigma ^3\mathrm{sinh}^2(z^{}\sigma )}}\right],`$ (124) $`\mathrm{tr}(r^{}r)^2r^2`$ $`=`$ $`N{\displaystyle \frac{N\mathrm{}}{\mathrm{}_\epsilon }}\left[{\displaystyle \frac{(96\sigma ^4z^{}+40z\sigma ^2+3z^{})\mathrm{coth}(z^{}\sigma )}{96\sigma ^4}}{\displaystyle \frac{40\sigma ^26z\sigma \mathrm{coth}(z^{}\sigma )+3i}{48\sigma ^3\mathrm{sinh}^2(z^{}\sigma )}}\right].`$ (125) In the limit $`\sigma 0`$, corresponding to $`L\mathrm{}_\epsilon `$, the weak-localization corrections to the conductance and the shot noise power and the conductance fluctuations approach their values for the chiral symmetry class cf. Eqs. (V B) and (115). For $`\sigma 1`$, corresponding to $`L\mathrm{}_\epsilon `$, one verifies that the values corresponding to the standard symmetry class are recovered. In this subsection, we have described the effect of a finite energy by a crossover length scale $`\mathrm{}_\epsilon `$. For the conductance and the shot noise power the limits of large and small $`\epsilon `$ correspond to the limits of $`L`$ large or small compared to $`\mathrm{}_\epsilon `$. However, upon inspection of Eq. (88) or (125) one observes that traces like $`\text{tr}r^2`$ that contain different numbers of $`r`$’s and $`r^{}`$’s, do not approach their large-energy limits $`\text{tr}r^2=0`$ as $`L\mathrm{}`$. The origin of this difference is that the reflection matrix is dominated by (interference of) paths that only enter a distance of the order of a mean free path $`\mathrm{}`$ into the quantum wire, while the conductance and the shot noise power depend on quantum interference throughout the entire wire. Hence, as long as $`\mathrm{}_\epsilon \mathrm{}`$, the finite energy cannot alter the interference of most paths that contribute to $`r`$. Hence, to judge whether the finite energy is relevant for the traces of reflection matrices, one has to compare $`\mathrm{}_\epsilon `$ to $`\mathrm{}`$ instead of $`L`$. We are now ready to define what is meant by “$`\epsilon `$ sufficiently large” in the crossover from the chiral symmetry class to the standard symmetry class. As far as quantum interference corrections to the transmission properties are concerned, the results of this subsection show that “$`\epsilon `$ sufficiently large” corresponds to the inequality of length scales $`L\mathrm{}_\epsilon `$, or equivalently, $`\epsilon \mathrm{}/L^2`$. However for reflection traces like $`\text{tr}r^2`$, a much more strict criteria is needed, $`\mathrm{}_\epsilon \mathrm{}`$, or $`\epsilon \mathrm{}`$. In the next section, these criteria, as well as the functional forms (117), (118) for the crossover will be compared to numerical simulations. ## VI Numerical simulations In this section we report on numerical simulations of the conductance of a quantum wire with random hopping only, and compare them with the theory of sections II-V. The simulations are for the random hopping model on a square lattice, described by the Schrödinger equation, $`\epsilon \psi _{m,j}`$ $`=`$ $`t_{m,j1;}\psi _{m,j1}t_{m,j;}^{}\psi _{m,j+1}`$ (127) $`t_{m1,j;}\psi _{m1,j}t_{m,j;}^{}\psi _{m+1,j},`$ where $`\psi _{m,j}`$ is the wave function at the lattice site $`(m,j)`$. A site is labeled by the chain index $`j=1,\mathrm{},N`$ and by the column index $`m`$. We impose open boundary conditions in the transverse direction, $`t_{m,0;}=t_{m,N;}=0`$. The system consists of a disordered region ($`0<m<L`$), coupled to the left and right to perfect leads ($`m<1`$ and $`m>L`$). In the leads, the longitudinal and transverse hopping amplitudes are $`t_{m,j;}=1`$ and $`t_{m,j;}=t`$, where $`0<t1`$. With this choice, there is a window of energies $`1+t<\epsilon <1t`$ around the band center, where the number of transmission channels does not depend on energy (and equals the number of chains $`N`$). In the disordered region, the hopping amplitudes are taken from a distribution centered around the values $`t_{m,j;}=1`$ and $`t_{m,j;}=t`$ for the leads. We consider two types of randomness, that we refer to as the real random hopping and random flux models. * In the real random hopping (RRH) model, the hopping amplitudes $`t_{m,j;}`$ and $`t_{m,j;}`$ are chosen uniformly and independently in the intervals $`t(1\delta )<t_{m,j;}<t(1+\delta )`$ and $`1\delta <t_{m,j;}<1+\delta `$, respectively, where $`\delta `$ measures the disorder strength. A uniform magnetic field with a flux $`\varphi _{\mathrm{pl}}`$ through each plaquette is modeled by multiplication of $`t_{m,j;}`$ with a Peierls phase $`e^{2\pi i\varphi _{\mathrm{pl}}(j1)}`$. * In the random flux (RF) model, the longitudinal hopping amplitudes $`t_{m,j;}=1`$, while the transverse hopping amplitudes $`t_{m,j;}`$ are complex numbers $`t_{m,j;}=te^{i\theta _{m,j}}`$. Here the $`\theta _{m,j}`$ are chosen such that the fluxes $`\varphi _{m,j}=\theta _{m,j}\theta _{m1,j}`$ are independently and uniformly distributed in the interval $`\pi p<\varphi _{m,j}<\pi p`$, where $`p`$ is a measure for the strength of the disorder. In the random flux model, the parameter $`\eta =0`$, see Ref. ; in the real random hopping model the precise value of $`\eta `$ is not known. However, nonzero $`\eta `$ (of order $`N^0`$ by assumption) will only give rise to corrections of relative order $`1/N`$, ($`1/N^2`$ for the average conductance), which can be neglected for large $`N`$. Since the statistics of the conductance in the RF model in a quasi-one-dimensional geometry has been studied extensively in Ref. at and away from the band center $`\epsilon =0`$, we restrict our attention here to the crossover as a function of energy. The wavefunctions that solve the Schrödinger equation (127) at energy $`\epsilon `$ can be written as $`\psi _{m,j}`$ $`=`$ $`{\displaystyle \underset{\nu =1}{\overset{N}{}}}{\displaystyle \frac{1}{\mathrm{sin}k_\nu }}[e^{ik_\nu m}\mathrm{sin}(q_\nu j)\psi _\epsilon ^{\mathrm{iL}}(\nu )`$ (129) $`+e^{ik_\nu m}\mathrm{sin}(q_{N+1\nu }j)\psi _\epsilon ^{\mathrm{oL}}(\nu )]`$ in the left lead and as $`\psi _{m,j}`$ $`=`$ $`{\displaystyle \underset{\nu =1}{\overset{N}{}}}{\displaystyle \frac{1}{\mathrm{sin}k_\nu }}[e^{ik_\nu m}\mathrm{sin}(q_{N+1\nu }j)\psi _\epsilon ^{\mathrm{iR}}(\nu )`$ (131) $`+e^{ik_\nu m}\mathrm{sin}(q_\nu j)\psi _\epsilon ^{\mathrm{oR}}(\nu )]`$ in the right lead, where the wave number $`k_\nu >0`$ is determined from $`\epsilon =2\mathrm{cos}k_\nu 2t\mathrm{cos}q_\nu `$ with $`q_\nu =\pi \nu /(N+1)`$. With this parameterization, the definition of the scattering matrix $`S_\epsilon `$ and its symmetries are the same as in Sec. II. For each realization of the disorder, the dimensionless conductance $`g`$ is computed from the Landauer formula (17). The recursive Green’s function method is used to calculate $`S_\epsilon `$. (Application of the method to the random hopping or random flux models is discussed in Ref. .) Our numerical simulations use the parameters $`t=0.6`$, $`\delta =0.2`$, and $`p=0.3`$, for the RRH and RF models. ### A Localized regime in the RRH model In the localized regime, the even-odd effect manifests itself most dramatically. Taking an average over $`2\times 10^4`$ realizations of the disorder, we have computed the mean and variance of $`g`$ at the band center $`\epsilon =0`$ for the RRH model with $`N=20`$ and $`21`$, and with and without a time-reversal breaking magnetic field, see Fig. 5. The magnetic field corresponds to a flux $`\varphi _{\mathrm{pl}}=8\times 10^4`$ per plaquette, or $`1`$ flux quantum per $`50`$ lattice spacings along the chain, so that time-reversal symmetry is broken for all but the shortest wire lengths shown in Fig. 5. For odd $`N(=21)`$ both $`g`$ and $`\mathrm{var}g`$ decrease algebraically whereas they decay exponentially for even $`N(=20)`$. We observe that, for odd (even) $`N`$ and fixed $`L`$, $`g`$ and $`\text{var}g`$ are larger (smaller) in the presence of a magnetic field, $`\beta =2`$, than without, $`\beta =1`$, in agreement with Eqs. (IV) and (65). Note that for small $`L`$, $`\text{var}g`$ is $`L`$-independent for the chiral unitary class, while $`\text{var}g`$ decreases linearly with $`L`$ for small $`L`$ in the chiral orthogonal class. Similar $`L`$-dependencies for small $`L`$ have been obtained for the standard symmetry classes, see Ref. . Results for the crossover from the chiral universality classes to the standard ones as a function of energy are shown in Figs. 6 and 7. Figure 6 shows the energy dependence of the localization length $`\xi =2lim_L\mathrm{}L/\mathrm{ln}g`$ \[cf. Eq. (63)\] for $`N=20`$; Fig. 7 shows numerical data for the ratio $`C=lim_L\mathrm{}\mathrm{ln}g/\mathrm{var}\mathrm{ln}g`$. Here, the averages were taken over $`500`$$`1000`$ realizations of the disorder and magnetic fields corresponding to fluxes $`\varphi _{\mathrm{pl}}=2,4,6\times 10^4`$ per plaquette, respectively, have been used. In the absence of a magnetic field, $`\xi (\epsilon )`$ shows non-monotonic behavior with a maximum around $`\epsilon 5\times 10^6`$, while, within $`10\%`$, the localization length $`\xi `$ is the same in the chiral orthogonal class ($`\epsilon =0`$) and in the standard orthogonal class ($`\epsilon 10^4`$ for the choice of parameters in the simulations), in agreement with Sec. IV. As we discussed in Sec. IV, the fact that $`\xi (\epsilon =0)=\xi (\epsilon 0)`$ in the absence of a magnetic field, could be interpreted as the result of a cancellation of two effects: The presence of an extra symmetry at the band center (the chiral symmetry), which tends to make $`\xi `$ shorter than away from the band center, and the enhancement of the DoS at the band center, which tends to make $`\xi `$ larger. Apparently, these two competing effects are not balanced in the crossover region, thereby giving rise to the non-monotonic energy dependence to $`\xi `$ displayed in Fig. 6. Such a non-monotonicity of $`\xi (\epsilon )`$ is reminiscent of the non-monotonic voltage dependence of differential conductance found in a normal-metal/superconductor microbridge. In the presence of a magnetic field, $`\xi `$ increases by a factor $`1.8`$ in the crossover from the chiral-unitary to the standard unitary symmetry class, which is slightly less than the factor $`2`$ predicted in Eq. (76). Note that the increase in localization length is most rapid around the same energy scale $`\epsilon _\mathrm{c}5\times 10^6`$ for which $`\xi `$ reaches its maximum, $`\xi _\mathrm{c}`$, in the absence of a magnetic field. Moreover, $`\epsilon _\mathrm{c}`$ is related to $`\xi _\mathrm{c}`$ by Thouless relation $`\xi _\mathrm{c}\sqrt{\mathrm{}/\epsilon _\mathrm{c}}`$ where the mean free path $`\mathrm{}`$ is obtained by dividing the localization length $`\xi (\epsilon =10^{10})`$ by $`N=20`$ in Fig. 6. A plot of the ratio $`C=lim_L\mathrm{}\mathrm{ln}g/\mathrm{var}\mathrm{ln}g`$ is shown in Fig. 7. In the standard symmetry classes, $`C`$ takes the universal value $`C=1/2`$, while in the chiral classes one has $`C`$ $`=`$ $`{\displaystyle \frac{\beta N/2}{N+\left(1\frac{2}{\pi }\right)(N2+2\eta )}}`$ (132) $`=`$ $`{\displaystyle \frac{\beta }{44/\pi }}+𝒪(1/N).`$ (133) The data shown in Fig. 7 confirm that $`C=0.5`$ in the standard symmetry classes (corresponding to $`\epsilon 10^4`$ for the parameters of our simulation). However, for the chiral symmetry classes, a $`20\%`$ discrepancy with Eq. (133) is found. While the numerical simulations for the localized regime qualitatively confirm the theory of Sec. IV, quantitative agreement is only up to $`20\%`$. As a possible source of this discrepancy, we point to the fact that the simulations are done for an appreciable disorder strength $`\delta =0.2`$, while the theory is derived for weak disorder, corresponding to $`\delta 0`$. Hence, the system cannot be considered truly (quasi) one-dimensional, and corrections from two-dimensional dynamics on shorter length scales need to be taken into account. Another cause of the observed discrepancies could be the uncertainty of the precise value of $`\eta `$ for the RRH model. While we believe that $`\eta `$ should not affect the conductance distribution significantly for large $`N`$, it remains difficult to make a quantitative assessment of finite-$`N`$ corrections as long as $`\eta `$ is unknown. At this moment, we are not aware of a direct way to obtain $`\eta `$ from the numerical simulations. ### B Diffusive regime in the RRH and RF models We next consider the crossover from the chiral symmetry classes to the standard symmetry classes in the diffusive regime $`\mathrm{}L\xi `$. In Figs. 8 and 9 we show numerical simulations of the average and variance of the conductance as a function of the energy $`\epsilon `$ and the magnetic flux $`\varphi =L(N1)\varphi _{\mathrm{pl}}`$ through the disordered part of the wire. The simulations are performed with $`N=45`$ and $`L=800`$, in order to ensure that the conditions $`\mathrm{}L\xi `$ for diffusive transport and $`NL`$ for quasi-one-dimensionality are both met. The ensemble average is taken over $`10^4`$ samples. Numerical results for the variance of the conductance in the RRH model versus $`\varphi `$ and $`\epsilon `$ are shown in Fig. 8. The numerical data of $`\text{var}g`$ versus $`\varphi `$ (Fig. 8a) agree within $`10\%`$ with the theoretical predictions $`\text{var}g=4/15`$ ($`2/15`$) for $`\varphi =0`$ and $`\text{var}g=2/15`$ ($`1/15`$) for $`\varphi 1`$ for $`\epsilon =0`$ ($`\epsilon 0`$). The crossover between the orthogonal and unitary classes happens for $`\varphi 1`$, both with chiral symmetry ($`\epsilon =0`$), and without ($`\epsilon =0.005`$). The $`\epsilon `$-dependence of $`\text{var}g`$ is shown in Fig. 8b, together with the theoretical result (118), where we fitted the crossover energy scale that enters into the definition of $`\sigma =L/\mathrm{}_\epsilon `$, cf. Eq. (116). Again we find quantitative agreement well within $`10\%`$. (The fact that the numerical data for $`\varphi =0`$ are below the theoretical curve can probably be attributed to the suppression of $`\text{var}g`$ as $`L`$ approaches the localization length $`\xi `$, see the remark in the discussion of Fig. 5.) Numerical results for the average conductance are shown in Fig. 9. All data shown are for the same length $`L=800`$ and for the same number of channels $`N=45`$. For weak disorder one can ignore the $`\varphi `$ and $`\epsilon `$ dependence of the mean free path $`\mathrm{}`$, and hence of the Drude term in the conductance. The only effect of a variation of $`\epsilon `$ or $`\varphi `$ is thus to change the symmetry of the quantum wire, which affects the weak-localization correction to the conductance. According to Eq. (117), we expect a nonzero weak-localization correction $`\delta g`$ in the standard orthogonal symmetry class, i.e., for $`\varphi =0`$ and $`\epsilon 0`$, while $`\delta g=0`$ if time-reversal symmetry is broken ($`\varphi 1`$) or if the chiral symmetry is present ($`\epsilon =0`$). This behavior is confirmed in Fig. 9. However, quantitatively, the numerical results differ $`30\%`$ from Eq. (117). In addition, Fig. 9 shows a small $`\epsilon `$ dependence of $`g`$ at large magnetic fields that cannot be accounted for within the theory of Sec. V. In particular, note the cusp-like structure at small $`\epsilon `$ in the $`\varphi =8`$ data in Fig. 9b. This effect seems to be too large to be explained by a spurious $`\epsilon `$-dependence of the mean free path $`\mathrm{}`$. Since the feature at small $`\epsilon `$ is suppressed at larger lengths $`L`$, a possible cause might be a contact resistance effect. (Contact resistance is known to play a role for disordered normal-metal–superconductor junctions, when the particle-hole degeneracy is destroyed by a finite voltage or by a magnetic field.) As we discussed in the previous subsection, other causes for the discrepancy between theory and numerical simulations may be the fact that the disorder is not small, or that the parameter $`\eta `$ is not known. ## VII Conclusions In the vicinity of the band center, the physics of localization in a quantum wire with chiral symmetry exhibits differences with respect to the case of a quantum wire in one of the standard symmetry classes. The most prominent differences are observed in the localized regime. For wires with chiral symmetry (as is the case with off-diagonal disorder), at the band center, the statistics of the conductance depends sensitively on the parity in the number of transmission channels $`N`$. For odd $`N`$, the band center represents a critical point that is characterized by the absence of exponential localization. The logarithm of the conductance is not self-averaging and the mean conductance or its variance decay algebraically with the length $`L`$ of the wire. For even $`N`$, exponential localization takes place with a self-averaging localization length $`\xi `$ that does not depend on the presence or absence of time-reversal and spin rotation symmetry. As the energy is tuned away from the band center, the system crosses over to the standard universality classes: The parity effect disappears and the localization length acquires a dependence on the presence or absence of time-reversal and spin rotation symmetry. In the presence of time-reversal symmetry, the localization length $`\xi `$ for even $`N`$ is the same with and without chiral symmetry. Our numerical simulations indicate that the crossover is non-monotonous: in the crossover between the chiral and standard symmetry classes, $`\xi `$ differs from the values in the pure symmetry classes. A complete theoretical description of this crossover is still lacking. In the diffusive regime, the differences between the chiral and standard universality classes are less pronounced. They show up in quantum interference corrections to the classical (Drude) conductance, which is the same in both cases. We have found that in the chiral classes, weak-localization corrections to the mean conductance $`g`$ and, more generally, to the density of transmission eigenvalues vanish at the band center. This is the quasi-one-dimensional counterpart of a similar observation made by Gade and Wegner in their study of two-dimensional disordered systems with chiral symmetry. The conductance fluctuations are twice as large at the band center relative to the standard universality classes, compatible with the presence of an extra symmetry in the system. We have calculated these quantum interference corrections as a function of energy, for the entire crossover from the chiral universality class at the band center $`\epsilon =0`$ to the standard unitary classes for $`\epsilon `$ far away from $`0`$. The theoretical predictions for this crossover agree qualitatively with numerical simulations though there remains sizable deviations between theory and numerics of the order of 10-30$`\%`$. While the chiral symmetry classes have received an enormous amount of theoretical attention (see the introduction of this paper for a brief summary), there are several hurdles to take before a chiral quantum wire can be realized in practice. Besides the effect of electron-electron interactions, which is not taken into account here, the main obstacle is the fact that the chiral symmetry is very fragile, since it is easily broken by, e.g., on-site random energies, next-nearest-neighbor hopping, or a small shift of the chemical potential, which will drive the system away from the chiral symmetric band center. Our calculation of the quantum interference corrections in the crossover from the chiral symmetry classes to the standard ones can be seen as a first and necessary step to tackle the latter obstacle. ###### Acknowledgements. C. M. acknowledges support from the Swiss National Science Foundation. P. W. B. acknowledges support by the NSF under grant nos. DMR 94-16910, DMR 96-30064, and DMR 97-14725 for the work done at Harvard. The work of A. F. was supported by a Grant-in-Aid from Japan Society for the Promotion of Science (No. 11740199). The numerical computations were performed at the Yukawa Institute Computer Facility. ## A Scaling equations In this appendix, we present the scaling equations needed to calculate the crossover for the weak localization corrections of the conductance, shot noise, and the universal conductance fluctuations of the conductance as done in Sec. V C. The notation was defined in Eq. (89) and we use the short-hands $$RR_{10},R_2R_{1010},R_3R_{101010}.$$ Equations needed up to corrections of order $`N^0`$ are $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_LR`$ $`=NN2R+R^2+{\displaystyle \frac{2\beta }{\beta }}NR_{00}^{}2RR_{00}^{}+R_{1000}^{}+R_{1000}^{}+R_2^{},`$ (A1) $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_LR_2^{}`$ $`=4RR_2^{}\left(NR\right)`$ (A4) $`+{\displaystyle \frac{2\beta }{\beta }}R_{00}^{}+4R+R_{00}^{}8R_2^{}4R_{1000}^{}4R_{1000}^{}+4R_3^{}+2R_{101000}^{}+2R_{101000}^{}+R_{100100}^{}+R_{100100}^{},`$ $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_L\mathrm{var}R`$ $`=4\left(RN\right)\mathrm{var}R+R_{00}^{}+R_{00}^{}+2R2R_{1000}^{}2R_{1000}^{}4R_2^{}+R_{100100}^{}+R_{100100}^{}+2R_3^{}`$ (A7) $`+{\displaystyle \frac{2\beta }{\beta }}R_{00}^{}+R_{00}^{}+2R2R_{1000}^{}2R_{1000}^{}4R_2^{}+R_{100100}^{}+R_{100100}^{}+2R_3^{}.`$ Equations needed up to corrections of order $`N`$ are $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_LR_{00}^{}`$ $`={\displaystyle \frac{4i\epsilon \gamma \mathrm{}}{\beta }}R_{00}^{}+NN2R_{00}^{}+R_{00}^{}^2,`$ (A8) $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_LR_{1000}^{}`$ $`={\displaystyle \frac{4i\epsilon \gamma \mathrm{}}{\beta }}R_{1000}^{}+N2R_{00}^{}+2R4R_{1000}^{}R_{00}^{}+2R2R_{1000}^{}R_{00}^{}R2R_{1000}^{}R,`$ (A10) $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_LR_{1100}^{}`$ $`=NR_{00}^{}+2R+R_{00}^{}4R_{1100}^{}2RR_{1100}^{}R_{00}^{}2R_{00}^{}2R_{1100}^{}R+R_{1100}^{}R_{00}^{},`$ (A12) $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_LR_3^{}`$ $`=3\left(2NR_2^{}R_3^{}+R4R_2^{}+2R_3^{}R+R_2^{}^2\right),`$ (A14) $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_LR_{101000}^{}`$ $`={\displaystyle \frac{4i\epsilon \gamma \mathrm{}}{\beta }}R_{101000}^{}+2\left(N2R_{1000}^{}+R_2^{}3R_{101000}^{}+RR_{1000}^{}R_2^{}+R_{101000}^{}R_{00}^{}\right)`$ (A17) $`+R6R_{1000}^{}2R_2^{}+4R_{101000}^{}R+R_{1000}^{}+2R_2^{}R_{1000}^{},`$ $`{\displaystyle \frac{\gamma \mathrm{}}{\beta }}_LR_{100100}^{}`$ $`={\displaystyle \frac{4i\epsilon \gamma \mathrm{}}{\beta }}R_{100100}^{}+2N2R_{1000}^{}+R_{1100}^{}3R_{100100}^{}+R_{00}^{}4R_{1000}^{}+2R_{100100}^{}R_{00}^{}`$ (A20) $`+2R2R_{1000}^{}2R_{1100}^{}+2R_{100100}^{}R+2R_{1000}^{}^2+R_{1100}^{}^2.`$
warning/0004/cond-mat0004012.html	ar5iv	text	# Contents ## 1 Introduction In recent years, the Calogero-Sutherland type $`N`$-body problems in one dimension have attracted considerable attention not only because they are exactly solvable but also due to their relationship with (1+1)-dimensional conformal field theory, random matrix theory etc. In particular, the connections between exactly solvable models and random matrix theory have been very fruitful. For example, by mapping these models to random matrices from an orthogonal, unitary or symplectic Gaussian ensemble, Sutherland was able to obtain all static correlation functions of the corresponding many body theory. The key point of this model is the pairwise long-range interaction among the $`N`$ particles. One may add here that the family consisting of exactly solvable models, related to fully integrable systems, is quite small and their importance lies in the fact that their small perturbations describe wide range of physically interesting situations. Further, recent developments relating equilibrium statistical mechanics to random matrix theory owing to non-integrability of dynamical systems has made the pursuit of unifying seemingly disparate ideas a very important theme. The results presented in this paper belong to the emerging intersection of several frontiers like quantum chaos, random matrix theory, many-body theory and equilibrium statistical mechanics . The universality in level correlations in linear (Gaussian) random matrix ensembles agrees very well with those in chaotic quantum systems as also in many-body systems like nuclei . On the other hand, random matrix theory was connected to the world of exactly solvable models when the Brownian motion model was presented by Dyson , and later on, by the works on level dynamics . However, there are dynamical systems which are neither chaotic nor integrable - the so-called pseudo-integrable systems . It is known that the spectral statistics of such systems are “non-universal with a universal trend” . In particular, for Aharonov-Bohm billiards, the level spacing distribution is linear for small spacing and it falls off exponentially for large spacing . Similar features are numerically observed for the Anderson model in three dimensions at the metal-insulator transition point . To understand these statistical features, and in the context of random banded matrices, a new random matrix model (which has been called as the short-range Dyson model in ) was introduced wherein the energy levels are treated as in the Coulomb gas model with the difference that only nearest neighbours interact. This new model explains features of intermediate statistics in some polygonal billiards. In view of all this it is worth enquiring if one can construct an $`N`$-body problem which is exactly solvable and which is connected to the short-range Dyson model (SRDM)? If possible, then using this correspondence one can hope to calculate the correlation functions of the corresponding many-body theory and see if the system exhibits long-range order and/or off-diagonal long-range order. The purpose of this paper is to present two such models in one dimension, one on a line and the other on a circle. We obtain the exact ground state and a part of the excitation spectrum on a line and the exact ground state on a circle in case the $`N`$ particles are interacting via nearest and next-to-nearest neighbour interactions . Further, in both the cases we show how the norm of the ground state wave function is related to the joint probability density function of the eigenvalues of short-range Dyson models. Using this mapping, we obtain one- and two-point functions of a related many-body theory in the thermodynamic limit and prove the absence of long-range order in the system. However, quite remarkably, we prove the existence of an off-diagonal long-range order in the symmetrized version of the corresponding many-body theory . We also extend this work in several different directions. For example, we consider an $`N`$-body problem with nearest and next-to-nearest neighbour interaction in an arbitrary number of dimensions $`D`$ and show that the ground state and a part of the excitation spectrum can still be obtained analytically. We also obtain a part of the bound state spectrum in one dimension (both on a full line and on a circle) by replacing the root system $`A_{N1}`$ by $`BC_N,D_N`$ etc. Besides, we also consider a model in two dimensions for which novel correlations are present in the ground as well as the excited states. The plan of the paper is the following. In Sec.II we consider an $`N`$-body problem on a line characterized by the Hamiltonian (throughout this paper we shall use $`\mathrm{}=m=1`$) $`H={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{^2}{x_i^2}}+g{\displaystyle \underset{i=1}{\overset{N1}{}}}{\displaystyle \frac{1}{(x_ix_{i+1})^2}}`$ $``$ $`G{\displaystyle \underset{i=2}{\overset{N1}{}}}{\displaystyle \frac{1}{(x_{i1}x_i)(x_ix_{i+1})}}`$ (1) $`+`$ $`V\left({\displaystyle \underset{i=1}{\overset{N}{}}}x_i^2\right)`$ with $`G0`$ while $`g>1/4`$ to prevent the collapse that a more attractive inversely quadratic potential would cause. We show that the ground state and at least a part of the excitation spectrum can be obtained if $$g=\beta (\beta 1),G=\beta ^2,V=\frac{\omega ^2}{2}\underset{i=1}{\overset{N}{}}x_i^2.$$ (2) Note that with the above restriction on $`G`$ and $`g`$, $`\beta 1/2`$. Further we also point out the connection between the norm of the ground state wave function and the joint probability distribution function for eigenvalues in SRDM. In Sec. III we consider another $`N`$-body problem, but this time on a circle characterized by the Hamiltonian $`H={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{^2}{x_i^2}}+g{\displaystyle \frac{\pi ^2}{L^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{1}{\mathrm{sin}^2[\frac{\pi }{L}(x_ix_{i+1})]}}`$ $`G{\displaystyle \frac{\pi ^2}{L^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{cot}\left[(x_{i1}x_i){\displaystyle \frac{\pi }{L}}\right]\mathrm{cot}\left[(x_ix_{i+1}){\displaystyle \frac{\pi }{L}}\right],(x_{N+1}=x_1),`$ (3) (where again $`G0`$ while $`g>1/4`$) and obtain the exact ground state in case $`g`$ and $`G`$ are again as related by eq. (2). Further, we also point out the connection between the norm of the ground state wave function and the joint probability distribution function for eigenvalues of short-range circular Dyson model (SRCDM). Using this connection, in Secs. IV and V we obtain several exact results for the corresponding many-body theory in the thermodynamic limit. In particular, in Sec. IV we calculate the two-particle correlation functions of a related many-body theory in the thermodynamic limit and prove the absence of long-range order in the system. In Sec. V we consider the symmetrized version of the model considered in Sec. III and show the existence of an off-diagonal long-range order in the bosonic system in the thermodynamic limit. In Sec. VI we consider the $`BC_N`$ generalization of the model (1) and obtain the exact ground state of the system. In Sec. VII we consider the $`BC_N`$ generalization of the model (3) and obtain the exact ground state of the system. In Sec. VIII we consider a generalization of the model (1) to higher dimensions and obtain the ground state and a part of the excitation spectrum. In Sec. IX we consider a variant of the model (1) in two dimensions and obtain the ground state as well a class of excited states all of which have a novel correlation built into them. Finally, in Sec. X we summarize the results obtained and point out several open questions. ## 2 N-body problem in one dimension on a line Let us start from the Hamiltonian (1) and restrict our attention to the sector of configuration space corresponding to a definite ordering of the particles, say $$x_ix_{i+1},i=1,2,\mathrm{},N1.$$ (4) On using the ansatz $$\psi =\varphi \underset{i=1}{\overset{N1}{}}(x_ix_{i+1})^\beta ,$$ (5) in the corresponding Schrödinger equation $`H\psi =E\psi `$, it is easily shown that, provided $`g`$ and $`G`$ are related to $`\beta `$ by eq. (2), $`\varphi `$ satisfies the equation $$\frac{1}{2}\underset{i=1}{\overset{N}{}}\frac{^2\varphi }{x_i^2}\beta \underset{i=1}{\overset{N1}{}}\frac{1}{(x_ix_{i+1})}\left(\frac{\varphi }{x_i}\frac{\varphi }{x_{i+1}}\right)+(VE)\varphi =0.$$ (6) Following Calogero we start from $`\varphi `$ as given by eq. (6) and assume that $$\varphi =P_k(x)\mathrm{\Phi }(r).$$ (7) where $`r^2=_{i=1}^Nx_i^2`$. The function, $`\mathrm{\Phi }`$ satisfies the equation $$\mathrm{\Phi }^{\prime \prime }(r)+[N+2k1+2(N1)\beta ]\frac{1}{r}\mathrm{\Phi }^{}(r)+2[EV(r)]\mathrm{\Phi }(r)=0,$$ (8) provided $`P_k(x)`$ is a homogeneous polynomial of degree $`k`$ ($`k=0,1,2,\mathrm{}`$) in the particle-coordinates and satisfies generalized Laplace equation $$\left[\underset{i=1}{\overset{N}{}}\frac{^2}{x_i^2}+2\beta \underset{i=1}{\overset{N1}{}}\frac{1}{(x_ix_{i+1})}\left(\frac{}{x_i}\frac{}{x_{i+1}}\right)\right]P_k(x)=0.$$ (9) We shall discuss few solutions of the Laplace equation (9) below. Let us now specialize to the case of the oscillator potential i.e. $`V(r)=\frac{\omega ^2}{2}r^2`$. In this case, eq. (8) is the well known radial equation for the oscillator problem in more than one dimension and its solution is $$\mathrm{\Phi }(r)=\mathrm{exp}(\omega r^2/2)L_n^a(\omega r^2),n=0,1,2,\mathrm{}.$$ (10) where $`L_n^a(x)`$ is the associated Laguerre polynomial while the energy eigenvalues are given by $$E_n=\left[2n+k+\frac{N}{2}+(N1)\beta \right]\omega =E_0+(2n+k)\omega ,$$ (11) with $`a=\frac{E}{\omega }2n1`$. Few comments are in order at this stage. 1. For large $`N`$, the energy $`E`$ is proportional to $`N`$ so that $$\underset{N\mathrm{}}{lim}\frac{E}{N}=\left(\beta +\frac{1}{2}\right)\omega ,$$ (12) i.e., the system has a good thermodynamic limit. In contrast, notice that the long-ranged Calogero model does not have a good thermodynamic limit since in that case for large $`N`$, $`E/N`$ goes like $`N`$. 2. The spectrum can be interpreted as due to noninteracting bosons (or fermions) plus $`(n,k)`$\- independent (but $`N`$-dependent) shift. The ground state eigenvalue and eigenfunction of the model is thus given by ($`n=k=0`$) $$E_0=\left[(N1)\beta +\frac{N}{2}\right]\omega ,$$ (13) $$\psi _0=\mathrm{exp}\left(\frac{\omega }{2}\underset{i=1}{\overset{N}{}}x_i^2\right)\underset{i=1}{\overset{N1}{}}(x_ix_{i+1})^\beta .$$ (14) A neat way of proving that we have indeed obtained the ground state can be given using the method of supersymmetric quantum mechanics . To this end, we define the operators $`Q_i`$ $`=`$ $`{\displaystyle \frac{d}{dx_i}}+\omega x_i+\beta \left[{\displaystyle \frac{1}{(x_{i1}x_i)}}{\displaystyle \frac{1}{(x_ix_{i+1})}}\right],(i=2,3,\mathrm{},N1),`$ $`Q_1`$ $`=`$ $`{\displaystyle \frac{d}{dx_1}}+\omega x_1\beta {\displaystyle \frac{1}{x_1x_2}},`$ $`Q_N`$ $`=`$ $`{\displaystyle \frac{d}{dx_N}}+\omega x_N+\beta {\displaystyle \frac{1}{x_{N1}x_N}},`$ (15) and their Hermitian conjugates $`Q_i^+`$. It is easy to see that the $`Q^{}s`$ annihilate the ground state as given by eq. (14). Further, the Hamiltonian (1) can be written in terms of these operators as $$HE_0=\frac{1}{2}\underset{i=1}{\overset{N}{}}Q_i^+Q_i,$$ (16) where $`E_0`$ is as given by eq. (13). Now since the operator on the right hand side is nonnegative and annihilates the ground state wavefunction as given by eq. (14), hence $`E_0`$ as given by eq. (13) must be the ground state energy of the system. On rewriting $`\psi _0`$ in terms of a new variable $$y_i\sqrt{\frac{\omega }{\beta }}x_i,$$ (17) one finds that the probability distribution for $`N`$ particles is given by $$\psi _0^2=C\mathrm{exp}\left(\beta \underset{i=1}{\overset{N}{}}y_i^2\right)\underset{i=1}{\overset{N1}{}}(y_iy_{i+1})^{2\beta }$$ (18) where $`C`$ is the normalization constant. We now observe that for $`\beta =1,2,4,`$ this $`\psi ^2`$ can be identified with the joint probability density function for the eigenvalues of SRDM with Gaussian orthogonal, unitary or symplectic ensembles respectively. We can therefore borrow the well-known results for these ensembles and obtain exact results about a many-body theory defined in the limit, $`N\mathrm{},\omega 0,N\omega =`$ finite which defines the density of the system. For example, as $`N\mathrm{}`$, the one-point function tends to a Gaussian for any $`\beta `$ and is given by $$R_1(x)=\frac{N}{\sqrt{2\pi \sigma ^2}}\mathrm{exp}(\frac{x^2}{2\sigma ^2}),$$ (19) where $`\sigma ^2=\frac{(\beta +1)}{\omega }`$. Other results about the many-body theory will be discussed in Secs. IV and V. Finally, let us discuss the polynomial solutions to the Laplace equation (9). So far, we have been able to obtain solutions in the following cases: (i) $`k=2,N2(ii)k=3,N3(iii)k=4,N4(iv)k=5,N5(v)k=6,N6`$. Besides we have also obtained solutions for $`k=4,5,6`$ in case $`N=3`$, and for $`k=5,6`$ in case $`N=4`$. We find that for $`k3`$, the demand that there be no pole in $`P_k(x)`$ alone does not require $`P_k(x)`$ to be completely symmetrical polynomial. However, for $`k=3,4`$ and $`N=3,4`$ it turns out that solution to Laplace eq. (9) exists only if $`P_k(x)`$ is a completely symmetric polynomial. We suspect that this may be true in general. On assuming completely symmetric $`P_k(x)`$ we find that in all the above cases we have a one-parameter family of solutions. In particular the various solutions are as follows (it is understood that the particle indices $`i,j,k,\mathrm{}`$ are always unequal unless mentioned otherwise). (i) $`k=2,N2`$ $$P_k(x)=a\underset{i=1}{\overset{N}{}}x_i^2+b\underset{i<j}{\overset{N}{}}x_ix_j,$$ (20) with $`\beta `$ given by $$\beta =\frac{aN}{(N1)(b2a)}.$$ (21) (ii) $`k=3,N3`$ $$P_k(x)=a\underset{i=1}{\overset{N}{}}x_i^3+b\underset{i,j=1}{\overset{N}{}}x_i^2x_j+c\underset{i<j<k}{\overset{N}{}}x_ix_jx_k,$$ (22) where $`c=3(ba)`$ and $`\beta `$ is given by $$\beta =\frac{3a+(N1)b}{(N1)(b3a)}.$$ (23) (iii) $`k=4,N4`$ $`P_k(x)=a{\displaystyle \underset{i=1}{\overset{N}{}}}x_i^4+b{\displaystyle \underset{i,j=1}{\overset{N}{}}}x_i^3x_j+c{\displaystyle \underset{i<j}{\overset{N}{}}}x_i^2x_j^2`$ $`+d{\displaystyle \underset{i,j<k}{\overset{N}{}}}x_i^2x_jx_k+e{\displaystyle \underset{i<j<k<l}{\overset{N}{}}}x_ix_jx_kx_l,`$ (24) where $`e=6(c2a),d=b+2c4a,`$ $`(N+4)b+2(N2)c4(N2)a+2(N1)(2a+bc)\beta =0,`$ (25) and $`\beta `$ is given by $$\beta =\frac{6a+(N1)c}{(N1)(b4a)}.$$ (26) (iv) $`k=5,N5`$ $`P_k(x)=a{\displaystyle \underset{i=1}{\overset{N}{}}}x_i^5+b{\displaystyle \underset{i,j=1}{\overset{N}{}}}x_i^4x_j+c{\displaystyle \underset{i,j=1}{\overset{N}{}}}x_i^3x_j^2+d{\displaystyle \underset{i,j<k}{\overset{N}{}}}x_i^3x_jx_k`$ $`+e{\displaystyle \underset{k,i<j}{\overset{N}{}}}x_i^2x_j^2x_k+f{\displaystyle \underset{i,j<k<l}{\overset{N}{}}}x_i^2x_jx_kx_l+g{\displaystyle \underset{i<j<k<l<m}{\overset{N}{}}}x_ix_jx_kx_lx_m,`$ (27) where $`e=5c5a3b,d=b+2c5a,`$ $`f=12c15a9b,g=30(cab),`$ $`(5N7)c3(N4)b5(N2)a`$ $`+(N1)(5a+3b2c)\beta =0,`$ (28) and $`\beta `$ is given by $$\beta =\frac{10a+(N1)c}{(N1)(b5a)}.$$ (29) (iv) $`k=6,N6`$ $`P_k(x)=a{\displaystyle \underset{i=1}{\overset{N}{}}}x_i^6+b{\displaystyle \underset{i,j=1}{\overset{N}{}}}x_i^5x_j+c{\displaystyle \underset{i,j=1}{\overset{N}{}}}x_i^4x_j^2+d{\displaystyle \underset{i,j<k}{\overset{N}{}}}x_i^4x_jx_k+e{\displaystyle \underset{i<j}{\overset{N}{}}}x_i^3x_j^3`$ $`+f{\displaystyle \underset{ij,k=1}{\overset{N}{}}}x_i^3x_j^2x_k+g{\displaystyle \underset{i,j<k<l}{\overset{N}{}}}x_i^3x_jx_kx_l+h{\displaystyle \underset{i<j<k}{\overset{N}{}}}x_i^2x_j^2x_k^2+p{\displaystyle \underset{i<j,<k<l}{\overset{N}{}}}x_i^2x_j^2x_kx_l`$ $`+q{\displaystyle \underset{i,j<k<l<m}{\overset{N}{}}}x_i^2x_jx_kx_lx_m+r{\displaystyle \underset{i<j<k<l<m<n}{\overset{N}{}}}x_ix_jx_kx_lx_mx_n,`$ (30) where $`3e=4b2c+6a+f,d=b+2c6a,g=2f+2c4b6a,`$ $`h=2f+9a4bc,p=5f+18a8b6c,`$ $`q=6(2f+9a4b3c),r=30(f+6a2b2c),`$ $`(5N9)f2(4N15)b+18(N5)a6(N5)c`$ $`+(N1)(8b+6c2f18a)\beta =0,`$ $`14b2c+6a+(N1)f+2(N1)(3a+2bc)\beta =0,`$ (31) and $`\beta `$ is given by $$\beta =\frac{15a+(N1)c}{(N1)(b6a)}.$$ (32) It would be nice if one can find solutions for higher values of $`k`$ and further check if solutions exist (if at all) only if $`P_k(x)`$ is a completely symmetric polynomial. While we are unable to prove it, we suspect that, subject to the solutions of the Laplace equation for higher k, we have obtained the complete spectrum for this problem. Finally it is worth enquiring if the bound state spectrum of the Hamiltonian (1) can also be obtained in case the oscillator potential is replaced by any other potential. It turns out that as in the Calogero case , in this case also the answer to the question is in affirmative. In particular, if instead the $`N`$ particles are interacting via the $`N`$-body potential as given by $$V(x_1,x_2,\mathrm{},x_N)=\alpha \underset{i=1}{\overset{N}{}}\frac{1}{\sqrt{_ix_i^2}},$$ (33) then also (most likely the entire) discrete spectrum can be obtained. This is because, after using the ansatz (7), eq. (8) is essentially the radial Schrödinger equation for an attractive Coulomb potential and it is well known that the only two problems which are analytically solvable for all partial waves are the Coulomb and the oscillator potentials. In particular the solution of (8) is then given by (note $`r^2=_{i=1}^Nx_i^2`$) $$\mathrm{\Phi }(r)=\mathrm{exp}(\sqrt{2E}r)L_n^b(2\sqrt{2E}r),$$ (34) and the corresponding energy eigenvalues are $$E_{n,k}=\frac{\alpha ^2}{2\left[n+k+\frac{N}{2}1+(N1)\beta \right]^2},$$ (35) when $`b=N+2k3+2(N1)\beta `$. It may be noted that whereas in the oscillator case the spectrum is linear in $`\beta `$, it is $`(E)^{1/2}`$ which is linear in $`\beta `$ in the case of the Coulomb-like potential. Secondly, as in any oscillator (Coulomb) problem, the energy depends on $`n`$ and $`k`$ only through the combination $`2n+k`$ ($`n+k`$). Is there any underlying reason why one is able to obtain the discrete spectrum for the $`N`$-body problem with either the oscillator or the Coulomb-like potential (33) ? Following it is easily shown that in both the cases one can write down an underlying $`SU(1,1)`$ algebra. Further, since the many-body potential $`W`$ in (1) is a homogeneous function of the coordinates of degree -2, i.e. it satisfies $$\underset{l=1}{\overset{N}{}}x_l\frac{W}{x_l}=2W,$$ (36) hence, following the arguments of , one can also establish a simple algebraic relationship between the energy eigenstates of the $`N`$-body problem (1) with the Coulomb-like potential (33) and the harmonic oscillator potential. It may be noted that the Hamiltonian (1) is not completely symmetric in the sense that whereas all other particles have two neighbours, particle 1 and $`N`$ have only one neighbour. Can one make it symmetric so that all particles will be treated on the same footing? One possible way is to add some extra terms in $`H`$. For example, consider $$H_1=H+H^{},$$ (37) where $`H`$ is as given by eq. (1) while $`H^{}`$ has the form $$H^{}=\frac{g}{(x_Nx_1)^2}G\left[\frac{1}{(x_Nx_1)(x_1x_2)}+\frac{1}{(x_{N1}x_N)(x_Nx_1)}\right].$$ (38) Clearly, by adding these extra terms, the problem has become cyclic invariant for any $`N`$ while for $`N=3`$ it is identical to the Calogero problem and hence is in fact completely symmetric under the interchange of any two of the three particle coordinates. It may be noted that in the thermodynamic limit, these extra terms are irrelevant. We can again start from the ansatz (5) (but with $`N1`$ replaced by $`N`$) in the Schrödinger equation $`H_1\psi =E\psi `$ and using eq. (2) we find that $`\varphi `$ again satisfies eq. (6) but with $`N1`$ in the second term being replaced by $`N`$. On further using the substitution as given by eq. (7) one finds that $`\mathrm{\Phi }`$ satisfies eq. (8) but with the coefficient of the $`2\beta `$ term being $`N`$ instead of $`N1`$ while $`P_k(x)`$ is again a homogeneous polynomial of degree $`k`$ (k=0,1,2,…) in the particle coordinates, which now satisfies instead of eq. (9) $$\left[\underset{i=1}{\overset{N}{}}\frac{^2}{x_i^2}+2\beta \underset{i=1}{\overset{N}{}}\frac{1}{(x_ix_{i+1})}\left(\frac{}{x_i}\frac{}{x_{i+1}}\right)\right]P_k(x)=0,$$ (39) with $`x_{N+1}=x_1`$. How do the solutions to the Laplace eqs. (9) and (39) compare? For $`N=3`$, eq. (39) is identical to that of Calogero and for this case Calogero has already obtained the solutions for any k. For $`N>3`$ and for $`k3`$, the demand that there be no pole in $`P_k(x)`$ alone does not require $`P_k(x)`$ to be completely symmetrical polynomial. However, for $`k=3,4`$ and $`N=4`$ it again turns out that solution to Laplace eq. (39) exists only if $`P_k(x)`$ is a completely symmetric polynomial. We suspect that this may be true in general. On assuming completely symmetric $`P_k(x)`$ we have been able to obtain a two-parameter family of solutions in case k=3,4,5,6 and $`Nk`$ (note that for eq. (9) we have obtained only one-parameter family of solutions). As an illustration, the solution for $`N4`$ and $`k=4`$ is given by (it is understood that the particle indices $`i,j,k,\mathrm{}`$ are always unequal) $`P_k(x)=a{\displaystyle \underset{i=1}{\overset{N}{}}}x_i^4+b{\displaystyle \underset{i,j=1}{\overset{N}{}}}x_i^3x_j+c{\displaystyle \underset{i<j}{\overset{N}{}}}x_i^2x_j^2`$ $`+d{\displaystyle \underset{i,j<k}{\overset{N}{}}}x_i^2x_jx_k+e{\displaystyle \underset{i<j<k<l}{\overset{N}{}}}x_ix_jx_kx_l,`$ (40) where $$e=2(2a+2d2bc),$$ (41) $$6a+(N1)c+\beta [8a2b+(2cd)(N2)]=0,$$ (42) $$6b+(N2)d+2\beta [2(N1)b2(N4)a+(N4)c2(N1)d]=0.$$ (43) Solution to the new $`\mathrm{\Phi }`$ equation can be easily written down in case $`V(r)=\frac{\omega ^2r^2}{2}`$ or if it is given by eq. (33). For example, it is easily checked that in the former case the solution is again given by eq. (11) but the energy eigenvalues are now given by $$E_{n,k}=[2n+k+\frac{N}{2}+N\beta ]\omega =E_0+(2n+k)\omega .$$ (44) Similarly, in the later case, the solution is as given by eq. (35) except that in the term containing $`\beta `$, $`N1`$ must be replaced by $`N`$. Apart from these two potentials, where we have obtained the entire bound state spectrum, there are several other potentials which are quasi-exactly solvable. For example, for the potential $$V\left(x_i^2\right)=A\underset{i=1}{\overset{N}{}}x_i^2B\left(\underset{i=1}{\overset{N}{}}x_i^2\right)^2+C\left(\underset{i=1}{\overset{N}{}}x_i^2\right)^3,$$ (45) it is easily shown that the ground state energy and eigenfunctions are $$E=\frac{B}{4\sqrt{C}}[N+2(N1)\beta ],$$ (46) $$\psi _0=\mathrm{exp}\left[\frac{\sqrt{C}}{4}(\underset{i=1}{\overset{N}{}}x_i^2)^2+\frac{B}{4\sqrt{C}}\underset{i=1}{\overset{N}{}}x_i^2\right]\underset{i=1}{\overset{N1}{}}(x_ix_{i+1})^\beta ,$$ (47) provided $`A,B,C`$ are related by $$A=\frac{B^2}{4C}[N+2+2(N1)\beta ]\sqrt{C}.$$ (48) It is worth enquiring if the probability distribution for $`N`$ particles corresponding to (47) can be mapped to some matrix model. In this context let us point out that the corresponding (long-ranged) Calogero problem was in fact mapped to the matrix model corresponding to branched polymers . So far as we are aware of, answer to this question is not known in our case. ## 3 N-body problem in one-dimension with periodic boundary condition Soon after the seminal papers of Calogero and Sutherland where they considered an $`N`$-body problem on full line, Sutherland also considered an $`N`$-body problem with long-ranged interaction and with periodic boundary condition. He obtained the exact ground state energy and showed that the corresponding $`N`$-particle probability density function is related to the random matrix with circular ensemble . Using the known results for the random matrix theory , he was able to obtain the static correlation functions of the corresponding many body theory. It is then natural to enquire if one can also obtain the exact ground state of an $`N`$-body problem with nearest and next-to-nearest neighbour interaction with periodic boundary condition (PBC). Further, one would like to enquire if the corresponding $`N`$-particle probability density can be mapped to some known matrix model. The hope is that in this case one may be able to obtain the correlation functions of a related many-body theory in the thermodynamic limit. We now show that the answer to the question is in the affirmative. Let us start from the Hamiltonian (1). We wish to find the ground state of the system subject to the periodic boundary condition (PBC) $$\psi (x_1,\mathrm{},x_i+L,\mathrm{},x_N)=\psi (x_1,\mathrm{},x_i,\mathrm{},x_N).$$ (49) For this, we start with a trial wave function of the form $$\mathrm{\Psi }_0=\underset{i=1}{\overset{N}{}}\mathrm{sin}^\beta \left[\frac{\pi }{L}(x_ix_{i+1})\right],(x_{N+1}=x_1).$$ (50) In this section, we restrict the coordinates $`x_i`$ to the sector $`Lx_1x_2\mathrm{}x_N0`$, so that eq. (50) makes sense even for noninteger $`\beta `$. The extension to the full configuration space will be made in Sec. 5. On substituting eq. (50) in the Schrödinger equation for the Hamiltonian (1), we find that it is indeed a solution provided $`g`$ and $`G`$ are again related to $`\beta `$ by eq. (2). The corresponding ground state energy turns out to be $$E_0=\frac{N\beta ^2\pi ^2}{L^2}.$$ (51) The fact that this is indeed the ground state energy can be neatly proved by using the operators $$Q_i=\frac{d}{dx_i}+\beta \frac{\pi }{L}\left[\mathrm{cot}(x_{i1}x_i)\mathrm{cot}(x_ix_{i+1})\right],$$ (52) and their Hermitian conjugates $`Q_i^+`$. It is easy to see that the $`Q^{}s`$ annihilate the ground state as given by eq. (50). The Hamiltonian (1) can be rewritten in terms of these operators as $$HE_0=\frac{1}{2}\underset{i}{}Q_i^+Q_i,$$ (53) where $`E_0`$ is as given by eq. (51). Hence $`E_0`$ must be the ground state energy of the system. Thus unlike the Calogero-Sutherland type of models, our models (both of Sec. II and here) have good thermodynamic limit, i.e., the ground state energy per particle $`(=E_0/N)`$ is finite as $`N\mathrm{}`$. Having obtained the exact ground state, it is natural to enquire if the corresponding $`N`$-particle probability density can be mapped to the joint probability distribution of some SRCDM so that we can obtain some exact results for the corresponding many-body theory. It turns out that indeed the square of the ground-state wave function is related to the joint probability distribution function for the SRCDM from where we conclude that the density is a constant if $`0x\frac{N}{L}`$, and zero outside. Other exact results for the many-body theory will be discussed in the next two sections. ## 4 Some exact results for the many-body problem The square of the ground-state wavefunction of the many-body problem introduced in Sec.II (Sec.III) can be identified with the joint probability distribution function of eigenvalues of the SRDM (SRCDM). Using SRCDM, Pandey and Bogomolny et al. have shown that for any $`\beta `$, the two-point correlation function has the form $$R_2^{(\beta )}(s)=\underset{n=1}{\overset{\mathrm{}}{}}P^{(\beta )}(n,s),$$ (54) where $`s`$ is the separation of two levels (or distance between two particles in the many-body theory considered here) and $$P^{(\beta )}(n,s)=\frac{(\beta +1)^{n(\beta +1)}}{\mathrm{\Gamma }[n(\beta +1)]}s^{n(\beta +1)1}e^{(\beta +1)s}.$$ (55) ¿From this expression it is not very easy to compute $`R_2(s)`$ for arbitrary $`\beta `$. However, it is easy to obtain the Laplace transform of $`R_2(s)`$ for any $`\beta `$. In particular, if $$g_2(t)=_0^{\mathrm{}}R_2(s)e^{ts}𝑑s,$$ (56) then $$g_2(t)=\underset{n=1}{\overset{\mathrm{}}{}}g(n,t),$$ (57) where $`g(n,t)`$ is the Laplace transform of $`P(n,s)`$, i.e., $$g(n,t)=_0^{\mathrm{}}P(n,s)e^{ts}𝑑s.$$ (58) On using $`P^{(\beta )}(n,s)`$ as given by eq. (55) in eq. (58) it is easily shown that $$g^{(\beta )}(n,t)=\left(\frac{\beta +1}{t+\beta +1}\right)^{(\beta +1)n}.$$ (59) Hence $$g_2^{(\beta )}(t)=\underset{n=1}{\overset{\mathrm{}}{}}g^{(\beta )}(n,t)=\frac{1}{(\frac{t+\beta +1}{\beta +1})^{\beta +1}1},$$ (60) from which one has to compute $`R_2^{(\beta )}(s)`$ by the Laplace inversion. For integer $`\beta `$, it is possible to perform the Laplace inversion by making use of the fact that $$\frac{1}{x^n1}=\frac{1}{n}\underset{k=0}{\overset{n1}{}}\frac{e^{2ik\pi /n}}{xe^{2ik\pi /n}},$$ (61) yielding $$R_2^{(\beta )}(s)=\underset{k=0}{\overset{\beta }{}}\mathrm{\Omega }^ke^{(\beta +1)s(\mathrm{\Omega }^k1)}$$ (62) where $$\mathrm{\Omega }=e^{2\pi i/(\beta +1)}.$$ (63) For $`\beta =1`$, which corresponds to the orthogonal ensemble, the result is already known : $`R_2^{(1)}(s)=1e^{4s}`$. It is interesting to mention that $`R_2^{(1)}(s)`$ agrees very well with some of the pseudo-integrable billiards (e.g., the $`\frac{\pi }{3}`$-rhombus billiard). It is important here to note that for rhombus billiards , the Hamiltonian matrix has elements which fall in their magnitude away from the principal diagonal. Thus, beyond a certain bandwidth, the elements are insignificant and the matrix is effectively banded. Immediately then, the results of banded matrices become applicable. Although there seems to be good agreement of the results from this random matrix theory as shown in , in it is also shown that there are other polygonal billiards for which $`R_2^{(1)}(s)`$ is not an appropriate correlator. It is possible that for different bandwidths, and, by an inclusion of interactions beyond nearest neighbours in the short-range Dyson model, a family of random matrices result. This may, eventually, explain the entire family of systems exhibiting intermediate spectral statistics. Coming back to the two-point correlation function, depending on if $`\beta `$ is an odd or an even integer, $`R_2(s)`$, as given by eq. (62), can be written in a closed form which shows that $`R_2(s)`$ is indeed real and further, it clearly exhibits oscillations for large $`s`$. In particular, it is easily shown that $`R_2(\beta =2p+1,s)`$ $`=`$ $`1e^{2(2p+2)s}+2e^{(2p+2)s}{\displaystyle \underset{m=1}{\overset{p}{}}}e^{(2p+2)s\mathrm{cos}(\frac{m\pi }{p+1})}`$ (64) $`\mathrm{cos}\left[{\displaystyle \frac{m\pi }{p+1}}+(2p+2)s\mathrm{sin}({\displaystyle \frac{m\pi }{p+1}})\right],`$ $`R_2(\beta =2p,s)`$ $`=`$ $`1+2e^{(2p+1)s}{\displaystyle \underset{m=1}{\overset{p}{}}}e^{(2p+1)s\mathrm{cos}(\frac{2m\pi }{2p+1})}`$ (65) $`\mathrm{cos}\left[{\displaystyle \frac{2m\pi }{2p+1}}+(2p+1)s\mathrm{sin}({\displaystyle \frac{2m\pi }{2p+1}})\right].`$ For illustration, we give below explicit expressions for $`\beta =2,3,4`$ $`R_2^{(2)}(s)`$ $`=`$ $`12e^{\frac{9s}{2}}\mathrm{cos}\left({\displaystyle \frac{3\sqrt{3}s}{2}}{\displaystyle \frac{\pi }{3}}\right);`$ $`R_2^{(3)}(s)`$ $`=`$ $`1e^{8s}2e^{4s}\mathrm{sin}(4s);`$ $`R_2^{(4)}(s)`$ $`=`$ $`1+2e^{5s(1+\mathrm{cos}(2\pi /5))}\mathrm{cos}\left[{\displaystyle \frac{2\pi }{5}}+5s\mathrm{sin}\left({\displaystyle \frac{2\pi }{5}}\right)\right]`$ (66) $`+`$ $`2e^{5s(1+\mathrm{cos}(4\pi /5))}\mathrm{cos}\left[{\displaystyle \frac{4\pi }{5}}+5s\mathrm{sin}\left({\displaystyle \frac{4\pi }{5}}\right)\right].`$ In Fig. 1, we have plotted $`R_2^{(\beta )}(s)`$ as a function of $`s`$ for $`\beta =1,2,3,4`$. These results show that, for integer $`\beta `$, there is no long-range order in the corresponding many-body theory. Similarly, if $`\beta `$ is half-integral, i.e., $`\beta =(2n+1)/2`$ then it is easily shown that $$R_2^{((2n+1)/2)}(s)=\frac{1}{2}\underset{k=0}{\overset{2n}{}}\mathrm{\Omega }^{2k}e^{\frac{2n+1}{2}s(1\mathrm{\Omega }^{2k})}\left[1+\mathrm{erf}\left(\sqrt{\frac{(2n+1)s}{2}}\mathrm{\Omega }^k\right)\right],$$ (67) where $`\mathrm{\Omega }`$ is as given by eq. (63). For arbitrary $`\beta `$, however, we are unable to perform the Laplace inversion and hence we do not have a closed expression for $`R_2(s)`$. However, one can numerically calculate it by using eqs. (54) and (55). In Fig. 2, we have plotted $`R_2^{(\beta )}(s)`$ as a function of $`s`$ for $`\beta =1,4/3,3/2,5/3,2,7/3,5/2`$. From this figure it is clear that even for fractional $`\beta `$, there is no long-range order. ## 5 Off-diagonal long-range order So far, nothing has been specified regarding the statistical character of the particles involved in the N-body problem of Sec. III. We now do that by first symmetrizing the Hamiltonian, that is by rewriting it as $`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{^2}{x_i^2}}`$ (68) $`+`$ $`{\displaystyle \underset{P\epsilon S_N}{}}\mathrm{\Theta }(x_{P(1)}x_{P(2)})\mathrm{}\mathrm{\Theta }(x_{P(N1)}x_{P(N)})W(x_{P(1)},\mathrm{},x_{P(N)}),`$ where $`\theta `$ is the step function and $`W(x_1,\mathrm{},x_N)`$ is the N-body potential of eq. (3). Next, relying on the solution given in eq. (50), we introduce the (not normalized) wave function: $$\psi _N(x_1,\mathrm{},x_N)=\epsilon _P\varphi _N(x_{P(1)},\mathrm{},x_{P(N)}),$$ (69) where P is the permutation in $`S_N`$ such that $`1>x_{P(1)}>x_{P(2)}>\mathrm{}>x_{P(N)}>0,\epsilon _P=1(\epsilon _P=sign(P))`$ in the N-boson (N-fermion) case and $$\varphi _N(x_1,\mathrm{},x_N)=\underset{n=1}{\overset{N}{}}sin\pi (x_nx_{n+1})^\beta ;(x_{N+1}=x_1),$$ (70) (we have set the scale factor $`L`$ equal to 1). Primitively, the function (69) is defined on the hypercube $`[0,1]^N`$. The following properties of $`\psi _N`$ are easily verified, provided that $`\beta 2`$: 1. In the bosonic case, $`\psi _N`$ can be continued to a multi-periodic function in the whole space $`^N`$ (or equivalently on the torus $`T^N`$): $$\psi _N(x_1,\mathrm{},x_i+1,\mathrm{},x_N)=\psi _N(x_1,\mathrm{},x_i,\mathrm{},x_N);(i=1,\mathrm{},N),$$ (71) which belongs to $`C^2`$ (i.e. is twice continuously differentiable). Owing to this property and the results of Sec.3, $`\psi _N`$ then obeys the Schrödinger equation (with Hamiltonian (3) and energy as given by eq. (51)) not only in the sector $`x_1>x_2>\mathrm{}>x_N`$ but every where. Thus, $`\psi _N`$ describes the ground state wave function of the N-boson system. Moreover, it is translation invariant (on $`^N`$): $$\psi _N(x_1+a,x_2+a,\mathrm{},x_N+a)=\psi _N(x_1,x_2,\mathrm{},x_N);Va\epsilon .$$ (72) 2. In the fermionic case, the continuation by periodicity is possible only for odd $`N`$, in which case eq. (71) still holds with $`\psi _N\epsilon C^2`$. For even $`N`$ on the contrary, enforcing the periodicity (71) leads to a discontinuous function $`\psi _N`$, so that the Schrödinger equation is no longer satisfied on the configuration space $`T^N`$. Therefore, in the following we shall implicitly restrict ourselves to odd values of $`N`$ when treating fermions. The translation invariance (72) then remains valid. We are interested in the one-particle reduced density matrix, given by $$\rho _N(xx^{})=\frac{N}{C_N}_0^1𝑑x_1\mathrm{}_0^1𝑑x_{N1}\psi _N(x_1,\mathrm{},x_{N1},x)\psi _N(x_1,\mathrm{},x_{N1},x^{}),$$ (73) where $`C_N`$ stands for the squared norm of the wave function: $$C_N=_0^1𝑑x_1\mathrm{}_0^1𝑑x_N\psi _N(x_1,\mathrm{},x_N)^2.$$ (74) That the R.H.S. of eq. (73) defines a (periodic) function of $`(xx^{})`$ is an easy consequence of eqs. (71) and (72). The normalization of $`\rho _N`$ is such that $`\rho _N(0)=N`$, the particle density. Further, the function $`\rho _N(\xi )`$ is manifestly of positive type on the $`U(1)`$ group, which implies that its Fourier coefficients, $$\rho _N^{(n)}=_0^1𝑑\xi e^{2i\pi n\xi }\rho _N(\xi );(n=0,\pm 1,\pm 2,\mathrm{}),$$ (75) are non-negative (Bochner’s theorem). In fact, this directly appears if one writes their explicit expression $`\rho _N^{(n)}`$ $`=`$ $`{\displaystyle \frac{N}{C_N}}{\displaystyle _0^1}𝑑x_1\mathrm{}{\displaystyle _0^1}𝑑x_{N1}\psi _N(x_1,\mathrm{},x_{N1},0)X`$ (76) $`X`$ $`{\displaystyle _0^1}𝑑xe^{2i\pi nx}\psi _N(x_1,\mathrm{},x_{N1},x),`$ in the form (obtained by using the periodicity property): $$\rho _N^{(n)}=\frac{N}{C_N}_0^1𝑑x_1\mathrm{}_0^1𝑑x_{N1}_0^1𝑑xe^{2i\pi nx}\psi _N(x_1,\mathrm{},x_{N1},x)^2.$$ (77) In the bosonic case, since the function $`\rho _N`$ is not only of positive type but also positive (like $`\psi _N`$), eq. (75) shows us that $$\rho _N^{(0)}\rho _N^{(n)};(n=\pm 1,\pm 2,\mathrm{}).$$ (78) In the fermionic case, eq. (78) is not necessarily true (because $`\psi _N`$ changes sign on $`T^N`$) and it is not an easy matter to determine the largest Fourier coefficient. Notice that the coefficients $`\rho _N^{(n)}`$, which physically represent the expectation values of the number of particles having momentum $`k_n=2\pi n`$ in the ground state, are nothing but the eigenvalues of the one-particle reduced density matrix (diagonal in the $`k_n`$ representation). According to the Onsager-Penrose criterion , no condensation can occur in the system (at least for Bose particles) if the largest of these eigenvalues is not an extensive quantity in the thermodynamic limit, that is, if $$\underset{N\mathrm{}}{lim}\frac{\rho _N^{(0)}}{N}=0.$$ (79) For Fermi particles, this criterion is not sufficient, and one has to look also at the largest eigenvalue of the two-particle reduced density matrix . Since we are presently unable to determine the largest eigenvalue of $`\rho _N`$ itself in the fermionic case, we shall not discuss extensively the latter here. Nevertheless, we shall look for the large $`N`$ behaviour of $`\rho _N^{(0)}`$ for bosons and fermions at a time, as this does not require much extra work and can give some indications in the fermionic case too. Let us write: $$\frac{\rho _N^{(0)}}{N}=\frac{A_N}{C_N},$$ (80) where $`C_N`$ is given by eq. (74) and $$A_N=_0^1𝑑x_1\mathrm{}_0^1𝑑x_{N1}\psi _N(x_1,\mathrm{},x_{N1},0)_0^1𝑑x\psi _N(x_1,\mathrm{},x_{N1},x),$$ (81) (the expression (76) of $`\rho _N^{(0)}`$ is more convenient than (77) for our purpose). Because of the special form (69)-(70) of the wave function, the computation of the squared norm $`C_N`$ is already not a trivial task, in sharp contrast to the case of N free, impenetrable particles. As a consequence, the (mainly algebraic) method introduced long ago by Lenard to deal with the latter case does not apply here, and we have to resort to another device. For conciseness, we introduce the notation: $$S(x_nx_{n1})_=\mathrm{sin}\pi (x_nx_{n+1})^\beta ,$$ (82) and define: $$S_2(\mathrm{})=_0^{\mathrm{}}𝑑xS(x)S(\mathrm{}x);(0\mathrm{}1).$$ (83) Our starting point will be the following representations of $`C_N`$ and $`A_N`$: $$C_N=(N1)!\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑xe^{ix}\stackrel{~}{F}(x)^N,$$ (84) $$A_N=(N1)!\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑xe^{ix}\stackrel{~}{F}(x)^{N3}\left[\stackrel{~}{F}(x)\stackrel{~}{G}(x)+\eta _N\stackrel{~}{H}(x)^2\right],$$ (85) where $`\stackrel{~}{F}(x)`$ $`=`$ $`{\displaystyle _0^1}𝑑\mathrm{}e^{i\mathrm{}x}S(\mathrm{})^2,`$ $`\stackrel{~}{G}(x)`$ $`=`$ $`{\displaystyle _0^1}𝑑\mathrm{}e^{i\mathrm{}x}S_2(\mathrm{})^2,`$ $`\stackrel{~}{H}(x)`$ $`=`$ $`{\displaystyle _0^1}𝑑\mathrm{}e^{i\mathrm{}x}S(\mathrm{})S_2(\mathrm{})^2,`$ (86) and $$\eta _N=\begin{array}{cc}(N2)\hfill & \text{for bosons}\hfill \\ 1\hfill & \text{for fermions}.\hfill \end{array}$$ (87) The representations (84)-(87) follow from the convolution structure of the expressions (74) and (81) of $`C_N`$ and $`A_N`$, when written in terms of appropriate variables. Their proof is given in the Appendix. Our aim is to extract from them the large $`N`$ behaviour of $`C_N`$ and $`A_N`$. Their form is especially suited for that purpose, because the integrands in eqs. (84) and (85) are entire functions, as polynomial combinations of Fourier transforms of functions with compact support (eq. (5)). Indeed, we are then allowed to, first, shift the integration path and then apply the residue theorem to meromorphic pieces of the integrands. However, it turns out that the calculations needed for arbitrary (integer) values of $`\beta `$ are quite cumbersome. So, in order to keep the argument clear enough, we shall content ourselves to present below these calculations in the simplest case, namely $`\beta `$ = 1 (recall that, strictly speaking, this value is not allowed), being understood that similar results are obtained for all integers $`\beta 2`$. For $`\beta =1,S(\mathrm{})=\mathrm{sin}\pi \mathrm{}`$, and eq. (5) gives, after reductions: $`\stackrel{~}{F}(x)`$ $`=`$ $`{\displaystyle \frac{2\pi ^2}{i}}{\displaystyle \frac{1e^{ix}}{x(x^24\pi ^2)}},`$ $`\stackrel{~}{G}(x)`$ $`=`$ $`{\displaystyle \frac{4\pi ^4}{i}}{\displaystyle \frac{5x^24\pi ^2}{x^3(x^24\pi ^2)^3}}+e^{ix}R^{(1)}(x),`$ $`\stackrel{~}{H}(x)`$ $`=`$ $`{\displaystyle \frac{4\pi ^3}{i}}{\displaystyle \frac{1}{x(x^24\pi ^2)^2}}+e^{ix}R^{(2)}(x),`$ (88) where $`R^{(n)}(x)`$ is a generic notation for rational functions behaving like $`x^n`$ when $`x\mathrm{}`$, and the precise form of which will be eventually of no importance. This produces, for the functions to be integrated in eqs. (84) and (85): $$\stackrel{~}{F}(x)^N=(\frac{2\pi ^2}{i})^N\left[\frac{1}{[x(x^24\pi ^2)]^N}+\underset{n=1}{\overset{N}{}}e^{inx}R_n^{(3N)}(x)\right],$$ (89) $`\stackrel{~}{F}(x)^{N3}[\stackrel{~}{F}(x)\stackrel{~}{G}(x)+\eta _N\stackrel{~}{H}(x)^2]`$ $`=i\left({\displaystyle \frac{2\pi ^2}{i}}\right)^N\left\{{\displaystyle \frac{(5+2\eta _N)x^24\pi ^2}{[x(x^24\pi ^2)]^{N+1}}}+{\displaystyle \underset{n=1}{\overset{N+1}{}}}e^{inx}R_n^{(3N1)}(x)\right\}.`$ (90) Let us stress again that these functions, when analytically continued, are holomorphic in the whole complex plane (the poles appearing in the first term are exactly canceled by the remaining ones). We consider first $`C_N`$, now given by $$C_N=(N1)!\left(\frac{2\pi ^2}{i}\right)^N\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑xe^{ix}\left\{\frac{1}{[x(x^24\pi ^2)]^N}+\underset{n=1}{\overset{N}{}}e^{inx}R_n^{(3N)}(x)\right\}.$$ (91) Since the function within the curly bracket is an entire one, we can shift the integration path to $`I\{z=x+iax\epsilon \}`$. Let us choose $`a>0`$. Then, by Cauchy theorem $$_I𝑑ze^{iz}\underset{n=1}{\overset{N}{}}e^{inz}R_n^{(3N)}(z)=0.$$ (92) Indeed, the integrand is holomorphic above $`I`$ and is bounded there by const. $`z^{3N}`$, which allows us to close the integration path at infinity in the upper complex plane. We end up with $$C_N=(N1)!\left(\frac{2\pi ^2}{i}\right)^N\frac{1}{2\pi }_I𝑑z\frac{e^{iz}}{z^N(z^24\pi ^2)^N}.$$ (93) Similarly, we are allowed to close the integration path at infinity in eq. (93), but this time in the lower complex plane. The integrand has now poles at $`z=0,\pm 2\pi `$, and applying the residue theorem leads to explicit expressions for $`C_N`$. Unfortunately, these expressions turn out to appear as (finite) sums with alternating signs, the terms of which become very close to each other for large $`N`$. They are therefore useless for determining the asymptotic behaviour of $`C_N`$, and we have to proceed differently. Let us write $`{\displaystyle _I}𝑑z{\displaystyle \frac{e^{iz}}{z^N(z^24\pi ^2)^N}}={\displaystyle \frac{1}{(N1)!}}{\displaystyle \frac{d^{N1}}{d\alpha ^{N1}}}_{\alpha =4\pi ^2}{\displaystyle _I}𝑑z{\displaystyle \frac{e^{iz}}{z^N(z^2\alpha )}}`$ $`={\displaystyle \frac{2i\pi }{(N1)!}}{\displaystyle \frac{d^{N1}}{d\alpha ^{N1}}}_{\alpha =4\pi ^2}[R_+(\alpha )+R_{}(\alpha )+R_0(\alpha )],`$ (94) where $`R_\pm (\alpha )`$ and $`R_0(\alpha )`$ are the residues of the last integrand at $`z=\pm \sqrt{\alpha }`$ and $`z=0`$ respectively. They are readily computed, assuming first that $`N=2M+1`$ is odd: $`R_+(\alpha )+R_{}(\alpha )`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}\sqrt{\alpha }}{\alpha ^{M+1}}}={\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^r}{(2r)!}}\alpha ^{rM1},`$ $`R_0(\alpha )`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{M}{}}}{\displaystyle \frac{(1)^r}{(2r)!}}\alpha ^{rM1}.`$ (95) Hence $$R_+(\alpha )+R_{}(\alpha )+R_0(\alpha )=(1)^{M+1}\underset{s=0}{\overset{\mathrm{}}{}}\frac{(1)^s}{(2M+2s+2)!}\alpha ^s$$ (96) Using eqs. (93), (5) and (96) we then obtain $`C_N`$ $`=`$ $`\left({\displaystyle \frac{2\pi ^2}{i}}\right)^N(1)^{M+1}(i){\displaystyle \frac{d^{N1}}{d\alpha ^{N1}}}_{\alpha =4\pi ^2}{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^s}{(2M+2s+2)!}}\alpha ^s`$ (97) $`=`$ $`(2\pi ^2)^N{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(N+n1)!}{n!(3N+2n1)!}}(4\pi ^2)^n.`$ The result is exactly the same for even $`N`$. It suffices now to observe that the last series alternates in sign and is decreasing to deduce $$C_N=(2\pi ^2)^N\frac{(N1)!}{(3N1)!}[1+O(\frac{1}{N})].$$ (98) Our procedure for evaluating $`A_N`$ is quite similar, and we give below only the main steps. From eqs. (85) and (5) we get $`A_N=(N1)!({\displaystyle \frac{2\pi ^2}{i}})^N{\displaystyle \frac{i}{2\pi }}{\displaystyle _I}𝑑ze^{iz}{\displaystyle \frac{(5+2\eta _N)z^24\pi ^2}{Z^{N+1}(z^24\pi ^2)^{N+1}}}`$ $`={\displaystyle \frac{1}{N}}({\displaystyle \frac{2\pi ^2}{i}})^N{\displaystyle \frac{i}{2\pi }}{\displaystyle \frac{d^N}{d\alpha ^N}}_{\alpha =4\pi ^2}{\displaystyle _I}𝑑ze^{iz}\left[{\displaystyle \frac{5+2\eta _N}{z^{N1}(z^2\alpha )}}{\displaystyle \frac{4\pi ^2}{z^{N+1}(z^2\alpha )}}\right],`$ (99) and, after computing the residues at $`z=\pm \sqrt{\alpha }`$ and $`z=0`$, we get $`A_N={\displaystyle \frac{(2\pi ^2)^N}{N}}{\displaystyle \frac{d^N}{d\alpha ^N}}_{\alpha =4\pi ^2}{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{5+2\eta _N}{(N+2s)!}}{\displaystyle \frac{4\pi ^2}{(N+2s+2)!}}\right](\alpha )^s`$ $`={\displaystyle \frac{(2\pi ^2)^N}{N}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(N+n)!}{n!}}\left[{\displaystyle \frac{5+2\eta _N}{(3N+2n)!}}{\displaystyle \frac{4\pi ^2}{(3N+2n+2)!}}\right](4\pi ^2)^n.`$ (100) Again, the last series alternates in sign and decreases, which entails $$A_N=(5+2\eta _N)(2\pi ^2)^N\frac{(N1)!}{(3N)!}[1+O(\frac{1}{N}].$$ (101) Finally, using eqs. (80), (98), (101) and (87) we obtain $$\frac{\rho _N^{(0)}}{N}=\frac{5+2\eta _N}{3N}\left[1+O\left(\frac{1}{N}\right)\right]=\begin{array}{cc}\frac{2}{3}[1+O(\frac{1}{N})]\hfill & \text{for bosons}\hfill \\ \frac{1}{N}[1+O(\frac{1}{N})]\hfill & \text{for fermions}\hfill \end{array}$$ (102) The same procedure applies for all integer values of $`\beta `$, although the algebra becomes quite involved. The general result for bosons (and for any integer $`\beta `$) is: $$\underset{N\mathrm{}}{lim}\frac{\rho _N^{(0)}}{N}=\frac{(\beta !)^4[(3\beta +1)!]^2}{[(2\beta )!]^2[(2\beta +1)!]^3}.$$ (103) Our method does not adapt straight forwardly to the case of non-integer values of $`\beta `$, but there is clearly no reason to expect a different outcome for such intermediate values. Therefore, the Onsager-Penrose criterion (79) is not met for bosons, and we reach the conclusion that Bose-Einstein condensation is possible in the bosonic version of the $`N`$-body model discussed in Sect.3. In the fermionic version, the result (102) is not conclusive, as explained after eq. (79). It only points (not too surprisingly) to the absence of quantum phase in the system. ## 6 The $`B_N`$ model in one dimension Subsequent to the seminal work of Calogero and Sutherland for the $`A_{N1}`$ system, the entire bound state spectrum of the Calogero model was obtained for $`BC_N,D_N`$ root systems . It is then natural to enquire if in our case, can one at least obtain the exact ground state and radial excitation spectrum in the $`BC_N`$ or $`D_N`$ case? We now show that the answer to this question is in the affirmative. Consider the $`BC_N`$ Hamiltonian, $`H={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{^2}{x_i^2}}+V\left({\displaystyle \underset{i=1}{\overset{N}{}}}x_i^2\right)+g{\displaystyle \underset{i=1}{\overset{N1}{}}}\left[{\displaystyle \frac{1}{(x_ix_{i+1})^2}}+{\displaystyle \frac{1}{(x_i+x_{i+1})^2}}\right]`$ $`G{\displaystyle \underset{i=2}{\overset{N1}{}}}\left[\left({\displaystyle \frac{1}{x_{i1}x_i}}{\displaystyle \frac{1}{x_{i1}+x_i}}\right)\left({\displaystyle \frac{1}{x_ix_{i+1}}}+{\displaystyle \frac{1}{x_i+x_{i+1}}}\right)\right]`$ $`+g_1{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{1}{x_i^2}},`$ (104) of which $`B_N,C_N`$ and $`D_N`$ are the special cases. We again restrict our attention to the sector of configuration space corresponding to a definite ordering of the particles as given by eq. (4). We start with the ansatz $$\psi =P_{2k}(x)\varphi (r)\left(\underset{i=1}{\overset{N}{}}(x_i^2)^{\gamma /2}\right)\underset{i=1}{\overset{N1}{}}(x_i^2x_{i+1}^2)^\beta ,$$ (105) where $`r^2=_{i=1}^Nx_i^2`$. On substituting it in the Schrödinger equation for the $`B_N`$-Hamiltonian (6) we find that $`\varphi `$ satisfies $$\mathrm{\Phi }^{^{\prime \prime }}(r)+[N+4k1+2N\gamma +4(N1)\beta ]\frac{1}{r}\mathrm{\Phi }^{}(r)+2\left[EV(r)\right]\mathrm{\Phi }(r)=0,$$ (106) provided $`g`$ and $`G`$ are again related to $`\beta `$ by eq. (2) while $`g_1`$ is related to $`\gamma `$ by $$g_1=\frac{\gamma }{2}(\gamma 1).$$ (107) Here, $`P_{2k}(x)`$ is a homogeneous polynomial of degree $`2k`$ ($`k=0,1,2,\mathrm{}`$) in the particle-coordinates and satisfies the generalized Laplace equation $$\left[\underset{i=1}{\overset{N}{}}\frac{^2}{x_i^2}+2\gamma \underset{i=1}{\overset{N}{}}\frac{1}{x_i}\frac{}{x_i}+4\beta \underset{i=1}{\overset{N1}{}}\frac{1}{(x_i^2x_{i+1}^2)}\left(x_i\frac{}{x_i}x_{i+1}\frac{}{x_{i+1}}\right)\right]P_{2k}(x)=0.$$ (108) Let us now specialize to the case of the oscillator potential, i.e., $`V(r)=\frac{\omega ^2}{2}r^2`$. In this case, (106) is the well known radial equation for the oscillator problem in more than one dimension and its solution is $$\mathrm{\Phi }(r)=\mathrm{exp}(\omega r^2/2)L_n^a(\omega r^2),n=0,1,2,\mathrm{}.$$ (109) where $`L_n^a(x)`$ is the associated Laguerre polynomial while the energy eigenvalues are given by $$E_n=\left[2n+2k+\frac{N}{2}+N\gamma +2(N1)\beta \right]\omega ,$$ (110) with $`a=\frac{E}{\omega }2n1`$. The exact ground state is obtained from here when $`n=k=0`$. The fact that $`n=k=0`$ gives the exact ground state energy of the system can be easily shown a la $`A_{N1}`$ case by the method of supersymmetric quantum mechanics. It may be noted that for large $`N`$, the energy $`E`$ is proportional to $`N`$ so that like the $`A_{N1}`$ case, the $`B_N`$ model also has a good thermodynamic limit. In contrast, notice that the long-ranged $`B_N`$ Calogero model does not have a good thermodynamic limit. Are there homogeneous polynomial solutions of eq. (108) of degree $`2k`$ ($`k1`$)? While we are unable to answer this question for any $`k`$, at least for small values of $`k`$ ($`k>0`$) there does not seem to be any solution to eq. (108). For example, we have failed to find any polynomial solution of degree 2,4 and 6. Thus it appears that unlike the $`A_{N1}`$ case, in the $`BC_N`$ case one is only able to obtain the ground state and radial excitations over it. Proceeding in the same way, the energy eigenvalues and eigenfunctions in the case of the Coulomb-like potential (33) are $$E=\frac{\alpha ^2}{2\left[n+2k+\frac{N1}{2}+N\gamma +2(N1)\beta \right]^2}$$ (111) $$\mathrm{\Phi }=e^{\sqrt{2Er}}L_n^b(2\sqrt{2E}r)$$ (112) where $`b=N2+4k+2N\gamma +4(N1)\beta `$. Again, so far we have been able to obtain solutions only in case k=0. As in Sec.II, in the $`BC_N`$ Hamiltonian (6), all the particles are not being treated on the same footing. Again, one possibility is to add extra terms. Consider for example, $$H_1=H+H^{},$$ (113) where $`H`$ is as given by eq. (6) while $`H^{}`$ has the form $`H^{}`$ $`=`$ $`g\left[{\displaystyle \frac{1}{(x_Nx_1)^2}}+{\displaystyle \frac{1}{(x_N+x_1)^2}}\right]`$ (114) $``$ $`G[({\displaystyle \frac{1}{x_Nx_1}}{\displaystyle \frac{1}{x_N+x_1}})({\displaystyle \frac{1}{x_1x_2}}+{\displaystyle \frac{1}{x_1+x_2}})`$ $`+`$ $`({\displaystyle \frac{1}{x_{N1}x_N}}{\displaystyle \frac{1}{x_{N1}+x_N}})({\displaystyle \frac{1}{x_Nx_1}}+{\displaystyle \frac{1}{x_N+x_1}})].`$ One can now run through the arguments as given above and show that the eigenstates for both the oscillator and Coulomb-like potentials have the same form as given above except that in the term multiplying $`\beta `$, $`N1`$ gets replaced by $`N`$ at all places including in the Laplace eq. (108). However, now we find that there are indeed solutions to the Laplace eq. (108) (with $`N1`$ replaced by $`N`$). In particular, the solution for any $`N(4`$) and $`k=4`$ is given by $$P_{k=4}(x)=a\underset{i=1}{\overset{N}{}}x_i^4+b\underset{i<j}{\overset{N}{}}x_i^2x_j^2,$$ (115) where $$\frac{b}{a}=2\left[\frac{3+8\beta +2\gamma }{N1+2(N1)\gamma +4(N2)\beta }\right].$$ (116) As in the $`A_{N1}`$ case, we again find that even though the Laplace eq. (108) is only invariant under cyclic permutations, the solution is in fact invariant under the permutation of any two coordinates. It will be interesting to try to find solutions for higher values of $`k`$ and study the full degeneracy of the spectrum. Besides these two, one can obtain a part of the spectra including the ground state for several other potentials but we shall not discuss them here. ## 7 $`BC_N`$ model in one dimension with periodic boundary condition Following the work of Sutherland on the $`A_{N1}`$ root system, the exact ground state as well as the excitation spectrum was also obtained in the case of the $`BC_N,D_N`$ root systems . It is then worth enquiring if, in our case, one can also obtain the ground state and the excitation spectrum. As a first step in that direction, we shall obtain the exact ground state of the $`BC_N`$ model with periodic boundary condition. The Hamiltonian for the $`BC_N`$ case is given by $`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{^2}{x_i^2}}+g{\displaystyle \frac{\pi ^2}{L^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left[{\displaystyle \frac{1}{\mathrm{sin}^2\frac{\pi }{L}(x_ix_{i+1})}}+{\displaystyle \frac{1}{\mathrm{sin}^2\frac{\pi }{L}(x_i+x_{i+1})}}\right]`$ (117) $`+`$ $`g_1{\displaystyle \frac{\pi ^2}{L^2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{\mathrm{sin}^2\frac{\pi }{L}x_i}}+g_2{\displaystyle \frac{\pi ^2}{L^2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{\mathrm{sin}^2\frac{2\pi }{L}x_i}}G{\displaystyle \frac{\pi ^2}{L^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}[\mathrm{cot}{\displaystyle \frac{\pi }{L}}(x_{i1}x_i)`$ $``$ $`\mathrm{cot}{\displaystyle \frac{\pi }{L}}(x_{i1}+x_i)\left]\right[\mathrm{cot}{\displaystyle \frac{\pi }{L}}(x_ix_{i+1})+\mathrm{cot}{\displaystyle \frac{\pi }{L}}(x_i+x_{i+1})].`$ We again restrict our attention to the sector of the configuration space corresponding to a definite ordering of the particles as given by eq. (2). For this case, we start with a trial wave function of the form $$\mathrm{\Psi }_0=\underset{i=1}{\overset{N}{}}\mathrm{sin}^\gamma \theta _i\underset{i=1}{\overset{N}{}}(\mathrm{sin}^22\theta _i)^{\gamma _1/2}\underset{i=1}{\overset{N}{}}[\mathrm{sin}^2(\theta _i\theta _{i+1})]^{\beta /2}\underset{i=1}{\overset{N}{}}[\mathrm{sin}^2(\theta _i+\theta _{i+1})]^{\beta 2},$$ (118) ($`\theta _i=\pi x_i/L`$) and substitute it in the Schrödinger equation for the Hamiltonian (117). We find that it is indeed a solution provided $`g`$ and $`G`$ are again related to $`\beta `$ by eq. (2) while $`g_1,g_2`$ are related to $`\gamma ,\gamma _1`$ by $$g_1=\frac{\gamma }{2}[\gamma +2\gamma _11],g_2=2\gamma _1(\gamma _11).$$ (119) The corresponding ground state energy turns out to be $$E_0=\frac{N\pi ^2}{2L^2}(\gamma +\gamma _1+2\beta )^2.$$ (120) The fact that this is indeed the ground state energy can be easily proved as in Secs. II and III. ## 8 $`N`$-body problem in $`D`$-dimensions Having obtained some results for the $`N`$-body problem (1) in one dimension, we study generalization to higher dimensions. Let us consider the following model in $`D`$-dimensions : $`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\stackrel{}{}_i^2+g{\displaystyle \underset{i=1}{\overset{N1}{}}}{\displaystyle \frac{1}{(\stackrel{}{r}_i\stackrel{}{r}_{i+1})^2}}`$ (121) $``$ $`G{\displaystyle \underset{i=2}{\overset{N1}{}}}{\displaystyle \frac{(\stackrel{}{r}_{i1}\stackrel{}{r}_i).(\stackrel{}{r}_i\stackrel{}{r}_{i+1})}{(\stackrel{}{r}_{i1}\stackrel{}{r}_i)^2(\stackrel{}{r}_i\stackrel{}{r}_{i+1})^2}}+V\left({\displaystyle \underset{i=1}{\overset{N}{}}}\stackrel{}{r}_i^2\right).`$ On using the ansatz, $$\psi =\left(\underset{i=1}{\overset{N1}{}}\stackrel{}{r}_i\stackrel{}{r}_{1+1}^\beta \right)\varphi (r),r^2=\underset{i=1}{\overset{N}{}}\stackrel{}{r}_i^2,$$ (122) in the Schrödinger equation for the Hamiltonian (121), it can be shown that $`\varphi (r)`$ satisfies $$\varphi ^{^{\prime \prime }}(r)+[DN1+2(N1)\beta ]\frac{1}{r}\varphi ^{}(r)+2(EV(r))\varphi (r)=0,$$ (123) provided $`g`$ and $`G`$ are related to $`\beta `$ by $$g=\beta ^2+(D2)\beta ,G=\beta ^2.$$ (124) Equation (123) is easily solved in the case of the oscillator potential (i.e., $`V(r)=\frac{\omega ^2}{2}r^2)`$ yielding the energy eigenstates as $$\varphi (r)=\mathrm{exp}\left(\frac{\omega }{2}r^2\right)L_n^b(\omega r^2),$$ (125) $$E_n=\left[2n+(N1)\beta +\frac{DN}{2}\right]\omega .$$ (126) Here $`b=\frac{E}{\omega }2n1`$. It may be noted that as in all other higher dimensional many-body problems, one has only obtained a part of the energy eigenvalue spectrum which however includes the ground state. In particular, the ground state energy eigenvalue and eigenfunction is given by $$E_0=\left[(N1)\beta +\frac{DN}{2}\right]\omega ,$$ (127) $$\psi _0=\mathrm{exp}\left(\frac{\omega }{2}\underset{i=1}{\overset{N}{}}r_i^2\right)\underset{i=1}{\overset{N1}{}}\stackrel{}{r}_i\stackrel{}{r}_{i+1}^\beta .$$ (128) As expected, for $`D`$ = 1 these results go over to those obtained in Sec. II. The fact that this is indeed the ground state energy can be easily proved by using again a supersymmetric formulation . At this point it is worth asking if the probability distribution for $`N`$ particles (at least for some $`D(>1))`$ can be mapped to some known random matrix ensemble ? In this context we recall that in the case of the Calogero-type model, it has been shown that in two space dimensions $`\psi _0^2`$ can be mapped to complex random matrix . Using this identification one was able to calculate all the correlation functions of the corresponding many-body theory and show the absence of long-range order but the presence of an off-diagonal long-range order in that theory. Unfortunately, so far as we are aware of, answer to this question is unknown in this particular case. We hope that at least in the case of two space dimensions, where $`\psi _0^2`$ for our model is given by $$\|\psi _0(z_i)\|^2=C\mathrm{exp}\left(\omega \underset{i=1}{\overset{N}{}}\|z_i\|^2\right)\underset{i=1}{\overset{N1}{}}\|z_iz_{i+1}\|^{2\beta },$$ (129) $`\psi _0^2`$ can be mapped to some variant of the short-range Dyson model. Finally, we observe that the ground state and a class of excited states can also be obtained in $`D`$-dimensions in case the oscillator potential is replaced by the $`N`$-body Coulomb-like potential $`V(r)=\alpha /\sqrt{𝐫_i^2}`$, because the resulting equation (123) is essentially the radial equation for the Coulomb potential. In particular, the energy eigenvalues and eigenfunctions are given by $$E_n=\frac{\alpha ^2}{2\left[n+\frac{DN1}{2}+(N1)\beta \right]^2},$$ (130) $$\psi _n=\mathrm{exp}(\sqrt{2\|E\|}r)L_n^b^{}(2\sqrt{2\|E\|r})\left(\underset{i=1}{\overset{N1}{}}\|𝐫_i𝐫_{i+1}\|^\beta \right),$$ (131) where $`b^{}=DN2+2(N1)\beta `$. It may again be noted that whereas the ground state energy is linear in $`\beta `$ in the oscillator case, it is not so in the case of the Coulomb-like $`N`$-body potential. ## 9 Short-range model in two dimensions with novel correlations Few years back, Murthy et al. considered a model in two dimensions with two-body and three-body long-ranged interactions and obtained the exact ground state and a class of excited states. The interesting feature of this model was that all these states had a built-in novel correlation of the form $`X_{ij}^g`$ where $$X_{ij}=x_iy_jx_jy_i.$$ (132) It is then natural to enquire if one can construct a model in two dimensions and obtain ground and few excited states of the system all of which would have a built-in short-range correlation of the form $$X_{j,j+1}=x_jy_{j+1}y_jx_{j+1}.$$ (133) We now show that this is indeed possible. Let us consider the following Hamiltonian $$H=\frac{1}{2}\underset{i=1}{\overset{N}{}}\stackrel{}{}_i^2+\frac{\omega ^2}{2}\underset{i=1}{\overset{N}{}}\stackrel{}{r}_i^2+g\underset{i=1}{\overset{N1}{}}\frac{\stackrel{}{r}_i^2+\stackrel{}{r}_{i+1}^2}{X_{i,i+1}^2}G\underset{i=2}{\overset{N1}{}}\frac{\stackrel{}{r}_{i1}\stackrel{}{r}_{i+1}}{X_{i1,i}X_{i,i+1}}$$ (134) where $`X_{i,i+1}`$ is as given by eq. (133). We start with the ansatz $$\psi (x_i,y_i)=\left[\underset{i=1}{\overset{N1}{}}X_{i,i+1}^\beta \right]\mathrm{exp}\left(\frac{\omega }{2}\underset{i}{}\stackrel{}{r}_i^2\right)\varphi (x_i,y_i).$$ (135) On substituting the ansatz in the Schrödinger equation $`H\psi =E\psi `$, one finds that $`\varphi `$ satisfies the equation $`[`$ $``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\stackrel{}{}_i^2+\omega {\displaystyle \underset{i=1}{\overset{N}{}}}\stackrel{}{r}_i\dot{\stackrel{}{}}_i+\beta {\displaystyle \underset{i=1}{\overset{N1}{}}}{\displaystyle \frac{1}{X_{i,i+1}}}(x_{i+1}{\displaystyle \frac{}{y_i}}y_{i+1}{\displaystyle \frac{}{x_i}}`$ (136) $`+`$ $`y_i{\displaystyle \frac{}{x_{i+1}}}x_i{\displaystyle \frac{}{y_{i+1}}})]\varphi =(E[N+2(N1)\beta ]\omega )\varphi ,`$ provided $`g`$ and $`G`$ are related by (2). It is interesting to note that even though we are considering the novel correlation model in two-dimensions, the relationship between $`g`$ and $`G`$ is as in the case of our one-dimensional model. We do not know if this has any deep significance. We conclude from here that $`\psi `$, as given by eq. (135), with $`\varphi `$ being a constant is the ground state of the system with the corresponding ground state energy being $$E_0=[N+2(N1)\beta ]\omega .$$ (137) Let us remark that, like the relationship between coupling constants, the ground state energy too has essentially the same form as that of the one-dimensional short-range $`A_{N1}`$ model as given by eq. (13). That one has indeed obtained the ground state can be proved as before. As in other many-body problems in two and higher dimensions, we are unable to find the complete excited-state spectrum. However, a class of excited states can be obtained from (136). To that end we introduce the complex coordinates $$z=x+iy,z^{}=xiy,\frac{}{z}=\frac{1}{2}\left(\frac{}{x}i\frac{}{y}\right),^{}\frac{}{z^{}}=\frac{1}{2}\left(\frac{}{x}+i\frac{}{y}\right).$$ (138) In terms of these coordinates, the differential eq. (136) takes the form $`[2{\displaystyle \underset{i=1}{\overset{N}{}}}_i_i^{}`$ $`+`$ $`2\beta {\displaystyle \underset{i=1}{\overset{N1}{}}}{\displaystyle \frac{\left(z_{i+1}_iz_i_{i+1}+z_i^{}_{i+1}^{}z_{i+1}^{}_i^{}\right)}{\left(z_iz_{i+1}^{}z_i^{}z_{i+1}\right)}}`$ (139) $`+`$ $`\omega {\displaystyle \underset{i=1}{\overset{N}{}}}(z_i_i+z_i^{}_i^{})(EE_0)]\varphi =0.`$ Now it is readily proved shown that the Hamiltonian H commutes with the total angular momentum operator $`L=_{i=1}^N(z_i_iz_i^{}_i^{}),`$ so that one can classify solutions according to their angular momentum: $`L\varphi =l\varphi `$. On defining $`t=\omega _iz_iz_i^{}`$ and let $`\varphi \varphi (t)`$ it is easily shown that $`\varphi (t)`$ satisfies $$t\varphi ^{^{\prime \prime }}(t)+\left[\frac{E_0}{\omega }t\right]\varphi ^{}(t)+\left(\frac{EE_0}{2\omega }\right)\varphi (t)=0,$$ (140) where $`E_0`$ is as given by eq. (137). Hence the allowed solutions with $`l`$ = 0 are $$E=E_0+2n\omega ,\varphi (t)=L_n^{\frac{E_0}{\omega }1}(t).$$ (141) Solutions with angular momentum $`l>0`$ or $`l<0`$ can similarly be obtained by introducing $`t_z=\omega _iz_i^2`$ or $`t_z^{}=\omega _i(z_i^{})^2`$. For example, let $`\varphi =\varphi (t_z)`$. Then eq. (139) reduces to $$2\omega t_z\frac{d\varphi }{dt_z}=(EE_0)\varphi .$$ (142) This is the well known Euler equation whose solutions are just monomials in $`t_z`$. The solution is given by $`\varphi (t_z)=t_z^m(m>0)`$, and hence the angular momentum $`l=2m`$ while the energy eigenvalues are $`E=E_0+2m\omega =E_0+l\omega `$. Further, we can combine these solutions with the $`l=0`$ solutions obtained above and obtain a tower of excited states. For example, let us define $`\varphi (z_i,z_i^{})=\varphi _1(t)\varphi _2(t_z)`$, where $`\varphi _1`$ is a solution with $`l=0`$, while $`\varphi _2`$ is the solution with $`l>0`$. On using $`\varphi _2(t_z)=t_z^m`$ it is easily shown that $`\varphi _1`$ again satisfies a confluent hypergeometric equation, $$t\varphi _1^{^{\prime \prime }}(t)+\left[\frac{E_0}{\omega }+2mt\right]\varphi _1^{}(t)+\left(\frac{EE_0}{2\omega }+m\right)\varphi _1(t)=0.$$ (143) Hence the energy eigenvalues are given by $`EE_0=(2n_r+2m)\omega `$. One may repeat the procedure to obtain exact solutions for a tower of states with $`l<0`$. ## 10 Summary In this paper we have discussed an $`N`$-body problem in one dimension and presented its exact ground state on a circle and most likely the entire spectrum on a real line. There are several similarities as well as differences between the model discussed here and Calogero-Sutherland (CSM) type of models and it might be worthwhile to compare the salient features of the two. 1. Whereas in CSM the interaction is between all neighbours, in our case the interaction is only between nearest and next-to-nearest neighbours. Note however that in both the cases it is an inverse square interaction. 2. Whereas in CSM (in one dimension) there is only two-body interaction, both two- and three-body interactions are required in our model for partial (or possibly exact) solvability on a real line. 3. Whereas the complete bound state spectrum is obtained in the Sutherland model (periodic potential) or if there is external harmonic or Coulomb-like $`N`$-body potential as given by eq. (33) and in the case of both $`A_{N1}`$ and $`BC_N`$ root systems, it is not clear if this is so in our case even though it is likely that this may be so in the $`A_{N1}`$ case. 4. Whereas our system, both on a line and on a circle, has good thermodynamic limit (i.e. $`E/N`$ is finite for large $`N`$), CSM does not have good thermodynamic limit in either case and $`E/N`$ diverges like $`N`$ for large $`N`$. 5. In both the cases, the norm of the ground state wavefunction can be mapped to the joint probability density function of the eigenvalues of some random matrix. Using this correspondence, in both the cases, one is able to calculate one- and two-point functions. However, whereas in the CSM this is possible only at three values of the coupling (corresponding to orthogonal, unitary or simplictic random matrices), in our case the correlation functions can be computed analytically for any integral or half-integral values of the coupling while numerically it can be done for any positive $`\beta `$. 6. In the CSM case with an external potential of the form $$V\left(\underset{i}{}x_i^2\right)=A\underset{i=1}{\overset{N}{}}x_i^2+B\left(\underset{i}{}x_i^2\right)^2+C\left(\underset{i}{}x_i^2\right)^3$$ (144) it has been shown that the norm of the ground state wave function can be mapped to a random matrix corresponding to branched polymers. It is not known if a similar mapping is possible in our case. 7. A multi-species generalization of CSM has been done , it is not clear if a similar generalization is possible in our case or not. 8. Generalization to $`D`$-dimensions ($`D>1`$) is possible in CSM as well as in our model and in both the cases one is able to obtain only a partial spectrum including the ground state. In both the cases, both two- and three-body interactions are required. Whereas our system has a good thermodynamic limit in any dimension $`D`$, the CSM does not have a good thermodynamic limit in any dimension. However, whereas the norm of the ground state wave function can be mapped to complex random matrices in the CSM case in two dimensions , no such mapping has so far been possible in our case for $`D>1`$. 9. Model with novel correlations is possible in two dimensions in both the cases but unlike CSM, our system has a good thermodynamic limit. 10. In the CSM, it has been possible to obtain the entire spectrum algebraically by using supersymmetry and shape invariance . It would be nice if similar thing can also be done in our model. Further, in the CSM, one has also written down the supersymmetric version of the model . It would be worth enquiring if a similar thing can also be done in our model. 11. In the CSM type models, one knows the various exactly solvable problems in which the $`N`$-particles interact pairwise by two body interaction . The question one would like to ask in our context is: what are the various exactly solvable problems in one dimension in which the $`N`$ particles have only nearest- and next-to-nearest neighbour interactions? 12. In the CSM, not only one- and two-point but even n-point correlation functions are known. It would be nice if the same is also possible in the present context. 13. A la Haldane-Shastry spin models , can we also construct spin models in the context of our model? 14. Unlike CSM, in our case the off-diagonal long-range order is nonzero in the bosonic version of the many-body theory in one dimension. Note however that the off-diagonal long-range order is nonzero in the CSM in two dimensions. Appendix 1. Proof of the representation (84) of $`C_N`$ By construction, the square of the wave function (69) is a symmetrical function of all its arguments, so that we can write eq. (74) as well: $$C_N=N!_0^1𝑑x_1_0^{x_1}𝑑x_2\mathrm{}_0^{x_{N1}}𝑑x_N\psi _N(x_1,\mathrm{},x_N)^2,$$ (145) where the particle coordinates are now properly ordered. We are thus allowed to substitute $`\varphi _N`$ for $`\psi _N`$ in (145) and obtain from eqs. (70) and (82) $$C_N=N!_0^1𝑑x_1_0^{x_1}𝑑x_2\mathrm{}_0^{x_{N1}}𝑑x_N\underset{n=1}{\overset{N}{}}S(x_nx_{n+1})^2.$$ (146) Changing the integration variables $`(x_1,x_2,\mathrm{},x_N)`$ to $`(\mathrm{}_1,\mathrm{}_2,\mathrm{},\mathrm{}_{N1},x_N)`$, where $$\mathrm{}_n=x_nx_{n+1};(n=1,\mathrm{},N1),$$ (147) one easily gets $`C_N`$ $`=`$ $`N!{\displaystyle _0^1}𝑑\mathrm{}_1{\displaystyle _0^{1\mathrm{}_1}}𝑑\mathrm{}_2\mathrm{}{\displaystyle _0^{1\mathrm{}_1\mathrm{}\mathrm{}_{N2}}}𝑑\mathrm{}_{N1}X`$ (148) $`X`$ $`{\displaystyle _0^{1_{p=1}^{N1}\mathrm{}_p}}𝑑x_N{\displaystyle \underset{n=1}{\overset{N1}{}}}S(\mathrm{}_n)^2S(x_Nx_1)^2.`$ Since $`x_Nx_1=_{p=1}^{N1}\mathrm{}_p`$ is in fact independent of $`x_N`$ in the new set of variables, eq. (148) becomes, using also $`S(x)=S(1x)`$: $`C_N`$ $`=`$ $`N!{\displaystyle _0^1}𝑑\mathrm{}_1{\displaystyle _0^{1\mathrm{}_1}}𝑑\mathrm{}_2\mathrm{}{\displaystyle _0^{1\mathrm{}_1\mathrm{}\mathrm{}_{N2}}}𝑑\mathrm{}_{N1}X`$ (149) $`X`$ $`(1{\displaystyle \underset{p=1}{\overset{N1}{}}}\mathrm{}_p){\displaystyle \underset{n=1}{\overset{N1}{}}}S(\mathrm{}_n)^2S(1{\displaystyle \underset{p=1}{\overset{N1}{}}}\mathrm{}_p)^2.`$ It is now convenient to introduce the extra variable $$\mathrm{}_N=1\underset{p=1}{\overset{N1}{}}\mathrm{}_p,$$ (150) and to recast eq. (149) in the form $`C_N`$ $`=`$ $`N!{\displaystyle _0^1}𝑑\mathrm{}_1{\displaystyle _0^1}𝑑\mathrm{}_2\mathrm{}{\displaystyle _0^1}𝑑\mathrm{}_{N1}{\displaystyle _0^1}𝑑\mathrm{}_N\delta (1{\displaystyle \underset{p=1}{\overset{N}{}}}\mathrm{}_p)X`$ (151) $`X`$ $`\mathrm{}_N{\displaystyle \underset{n=1}{\overset{N}{}}}S(\mathrm{}_N)^2`$ $`=`$ $`N!{\displaystyle _0^1}𝑑\mathrm{}_1\mathrm{}{\displaystyle _0^1}𝑑\mathrm{}_N\delta (1{\displaystyle \underset{p=1}{\overset{N}{}}}\mathrm{}_p){\displaystyle \frac{1}{N}}{\displaystyle \underset{m=1}{\overset{N}{}}}\mathrm{}_m{\displaystyle \underset{n=1}{\overset{N}{}}}S(\mathrm{}_n)^2`$ $`=`$ $`(N1)!{\displaystyle _0^1}𝑑\mathrm{}_1\mathrm{}{\displaystyle _0^1}𝑑\mathrm{}_N\delta (1{\displaystyle \underset{p=1}{\overset{N}{}}}\mathrm{}_p){\displaystyle \underset{n=1}{\overset{N}{}}}S(\mathrm{}_n)^2.`$ In the second equality, we have used the fact that, apart from the factor $`\mathrm{}_N`$, the integrand and the integration range are completely symmetrical in the variables $`(\mathrm{}_1,\mathrm{},\mathrm{}_N)`$. Finally, the integration over these variables factorizes after introducing the representation $$\delta (1\underset{p=1}{\overset{N}{}}\mathrm{}_p)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑xe^{ix(1_{p=1}^N\mathrm{}_p)},$$ (152) and interchanging the x- and $`\mathrm{}`$integrations. This produces eq. (84). 2. Proof of the representation (85) of $`A_N`$ Proceeding along the same lines, we first put the expression (81) of $`A_N`$ in the form $`A_N`$ $`=`$ $`(N1)!{\displaystyle _0^1}𝑑x_1{\displaystyle _0^{x_1}}𝑑x_2\mathrm{}{\displaystyle _0^{x_{N2}}}𝑑x_{N1}\varphi _N(x_1,\mathrm{},x_{N1},0)X`$ (153) $`X`$ $`R_N(x_1,\mathrm{},x_{N1}),`$ where $`R_N(x_1,\mathrm{},x_{N1})={\displaystyle _0^{x_{N1}}}𝑑x\varphi _N(x_1,\mathrm{},x_{N1},x)`$ $`\pm {\displaystyle _{x_{N1}}^{x_{N2}}}𝑑x\varphi _N(x_1,\mathrm{},x,x_{N1})+\mathrm{}+{\displaystyle _{x_1}^1}𝑑x\varphi _N(x,x_1,\mathrm{},x_{N1})`$ $`={\displaystyle _0^{x_{N1}}}𝑑x\varphi _N(x_1,\mathrm{},x_{N1},x)+{\displaystyle _{x_1}^1}𝑑x\varphi _N(x,x_1,\mathrm{},x_{N1})`$ $`+{\displaystyle \underset{p=1}{\overset{N2}{}}}\nu _p{\displaystyle _{x_{p+1}}^{x_p}}𝑑x\varphi _N(x_1,\mathrm{},x_p,x,x_{p+1},\mathrm{},x_N).`$ (154) Here, $`\nu _p=1(\nu _p=(1)^p)`$ for bosons (fermions) and we have used the restriction to odd $`N`$ in the second case. Thanks to the periodicity and the cyclic symmetry of $`\varphi _N`$ , the first two terms in the last expression above can be collected to give $$_{x_11}^{x_{N1}}𝑑x\varphi _N(x,x_1,\mathrm{},x_{N1}).$$ Hence $`R_N`$ becomes (with $`x_N=x_11`$) $`R_N(x_1,\mathrm{},x_{N1})={\displaystyle \underset{p=1}{\overset{N1}{}}}\nu _p{\displaystyle _{x_{p+1}}^{x_p}}𝑑x\varphi _N(x_1,\mathrm{},x_p,x,x_{p+1},\mathrm{},x_{N1})`$ $`={\displaystyle \underset{p=1}{\overset{N1}{}}}\nu _p{\displaystyle \underset{n=1}{\overset{N1}{}}}S(x_nx_{n+1}){\displaystyle _{x_{p+1}}^{x_p}}𝑑xS(xx_{p+1})S(x_px),(np)`$ $`={\displaystyle \underset{p=1}{\overset{N1}{}}}\nu _p{\displaystyle \underset{n=1}{\overset{N1}{}}}S(x_nx_{n+1})S_2(x_px_{p+1}),(np)`$ (155) according to the definition (83). We also have: $$\varphi _N(x_1,\mathrm{},x_{N1},0)=\underset{m=1}{\overset{N2}{}}S(x_mx_{m+1})S(x_{N1})S(x_1).$$ (156) Inserting eqs. (2) and (156) in eq. (153) and introducing as before the new integration variables $`\mathrm{}_nx_nx_{n+1}`$ ($`n=1,\mathrm{},N2)`$ and $`x_{N1}`$, we obtain $`A_N`$ $`=`$ $`(N1)!{\displaystyle _0^1}d\mathrm{}_1{\displaystyle _0^{1\mathrm{}_1}}d\mathrm{}_2...{\displaystyle _0^{1\mathrm{}_1\mathrm{}\mathrm{}_{N3}}}d\mathrm{}_{N2}X`$ (157) $`X`$ $`{\displaystyle _0^{\mathrm{}_{N1}}}𝑑x_{N1}{\displaystyle \underset{m=1}{\overset{N2}{}}}S(\mathrm{}_m)S(x_{N1})S(x_{N1}\mathrm{}_{N1})X`$ $`X`$ $`{\displaystyle \underset{p=1}{\overset{N1}{}}}\nu _p{\displaystyle \underset{n=1}{\overset{N1}{}}}S(\mathrm{}_n)S_2(\mathrm{}_p);(np)`$ where $`\mathrm{}_{N1}=1_{p=1}^{N2}\mathrm{}_p`$. The integration over $`x_{N1}`$ gives the factor $`S_2(\mathrm{}_{N1})`$ in place of $`S(x_{N1})S(x_{N1}\mathrm{}_{N1})`$, so that $`A_N=(N1)!{\displaystyle _0^1}d\mathrm{}_1\mathrm{}{\displaystyle _0^1}d\mathrm{}_{N1}\delta (1{\displaystyle \underset{p=1}{\overset{N1}{}}}\mathrm{}_p)[{\displaystyle \underset{m=1}{\overset{N2}{}}}X`$ $`XS(\mathrm{}_m)S_2(\mathrm{}_{N1}){\displaystyle \underset{p=1}{\overset{N1}{}}}\nu _p{\displaystyle \underset{n=1}{\overset{N1}{}}}S(\mathrm{}_n)S_2(\mathrm{}_p)],(np).`$ (158) On taking into account the complete symmetry of the integration measure, one finds that the square bracket in eq. (2) can be replaced by $`[\mathrm{}]`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{N2}{}}}S(\mathrm{}_m)^2S_2(\mathrm{}_{N1})^2+\eta _N{\displaystyle \underset{m=1}{\overset{N3}{}}}S(\mathrm{}_m)^2X`$ (159) $`X`$ $`[S(\mathrm{}_{N2})S_2(\mathrm{}_{N2})][S(\mathrm{}_{N1})S_2(\mathrm{}_{N1})],`$ where $`\eta _N`$ is as defined in eq. (87). Finally, one obtains the factorization of the multiple integral in eq. (2) by using again the representation (152) of the $`\delta `$ measure (with $`(N1)`$ in place of $`N`$). This entails eq. (85). A last remark may be in order. Alternative, equivalent forms of the representations (84) and (85) would be obtained by relying on Fourier expansions instead of Fourier integrals, that is by considering the integrands in eqs. (149) and (157) not as functions with compact supports $`[0,1]^N^N`$, resp. $`[0,1]^{N1}^{N1}`$, but as periodic functions (this would amount to modifying eq. (152) accordingly). It turns out however that the resulting representations of $`C_N`$ and $`A_N`$ (as Fourier series) are much less convenient for the explicit or asymptotic evaluations of these quantities. Acknowledgements SRJ acknowledges the warm hospitality of the Institute of Physics, Bhubaneswar where this work was initiated while AK would like to thank the members of the Laboratoire de Physique Mathématique of Montpellier University for warm hospitality during his trip there as a part of the Indo-French Collaboration Project 1501-1502. Figure Legends Fig. 1 The two-point correlation function for four integer values of $`\beta `$ (from left to rightmost are increasing values from 1 to 4) shows clearly an absence of long-range order. Fig. 2 The two-point correlation function for some fractional values of $`\beta `$ plotted alongwith $`\beta `$ equal to 1 and 2. ¿From left to rightmost are increasing values from 1, 4/3, 3/2, 5/3, 2, 7/3, and 5/2. Thus, even for fractional values, there is no long-range order.
warning/0004/cond-mat0004077.html	ar5iv	text	# Ga NMR study of the local susceptibility in SrCr8Ga4O19: pseudogap and paramagnetic defects ## Abstract We present the first Ga(4f) NMR study of the Cr susceptibility in the archetype of Kagomé based frustrated antiferromagnets, SrCr<sub>8</sub>Ga<sub>4</sub>O<sub>19</sub>. Our major finding is that the susceptibility of the frustrated lattice goes through a maximum around 50 K. Our data also supports the existence of paramagnetic “clusters” of spins, responsible for the Curie behavior observed in the macroscopic susceptibility at low $`T`$. These results set novel features for the constantly debated physics of geometrically frustrated magnets. The interest in triangular-based antiferromagnets (AF) was raised long ago by Anderson’s suggestion for a resonating valence bond state (R.V.B) as an open alternative to the classical Néel state . After a decade of extensive work on 2D geometrically frustrated AF, there is now a growing theoretical consensus that the $`S=\frac{1}{2}`$ Kagomé Heisenberg AF is a very promising candidate for such an RVB state at low $`T`$. Instead of sharing bonds, as in a triangular lattice, the frustrated triangles share sites with a smaller coordinance. From a classical point of view, this generates a very high degeneracy of the ground state, which translates into a huge density of low energy excitations, an absence of long range order at $`T=0`$ and a very short magnetic correlation length characteristic of a so-called ”spin liquid state”. Quantum approaches for $`S=1/2`$ spins suggest that a singlet ground state is energetically favoured . Only some of these theoretical features were demonstrated experimentally on the archetypal system SrCr<sub>9p</sub>Ga<sub>12-9p</sub>O<sub>19</sub> (SCGO:$`p`$) Kagomé-based AF. Susceptibility measurements showed the strong AF character of the interactions through the occurrence of a Curie-Weiss (CW) temperature as high as $`\theta \mathrm{\hspace{0.17em}500}600`$ K (for $`p0.9`$). The extension of the CW law well below $`\theta `$ and the smallness of the spin glass like ordering detected in susceptibility only at low $`T\theta ,`$ ($`T_g=35`$ K) are taken as convincing signatures of the high frustration. In addition, neutron data and low $`T`$ $`\mu `$SR studies altogether prove that the freezing is quite marginal and involves only a 20-30% fraction of the moment(s). The magnetic correlation length is also found of the order of the inter-Cr spacing. However, due to a Curie upturn which dominates the macroscopic susceptibility, $`\chi _{macro}`$, below 60 K there is no experimental determination of the $`T\theta `$ Kagomé-based lattice susceptibility. The Curie upturn received many interpretations still highly debated, encompassing magnetic disorder due to non magnetic substitutions in the Kagomé plane or original geometrically frustrated spin glass like order . In this Letter, we present the first NMR study of local susceptibility in SCGO, where we can discriminate between the susceptibility associated with the Kagomé based frustrated magnetism, $`\chi _{frust}`$ , and the purely paramagnetic susceptibility, $`\chi _{def}`$, most likely induced by Ga/Cr substitutions (defects). Strikingly, we find a gap-like downturn of $`\chi _{frust}`$ below $`T`$ as high as 50 K. This sets a novel energy scale, between $`\theta `$ and $`T_g`$ for the relevant physics of SCGO. From $`\chi _{macro}`$ measurements, the SCGO:$`p=0.9`$ sample was found to be typical for such a Cr content , with $`\theta =560`$ K and a high-$`T`$ Curie constant, $`C=2.2`$ emu/mol. An extended NMR analysis of the spectrum was presented in a previous paper . There, it was clearly established that 3 different Ga NMR sites could be resolved, Ga($`4e`$), Ga($`4f`$) and Ga substituted on Cr sites, Ga($`sub`$). We focus here on the Ga($`4f`$) site, whose line shift probes both the Cr($`12k`$) Kagomé and Cr($`2a`$) triangular planes susceptibility through the neighboring O. Interestingly, these sites form the so-called pyrochlore slab (fig.1). Below $`T=200`$ K the <sup>71</sup>Ga spectra were recorded by sweeping the field, using a conventional $`\pi /2`$ \- $`\pi `$ spin echo sequence to detect the <sup>71</sup>Ga($`4f`$) and <sup>69</sup>Ga($`4f`$) nuclear transitions. Compared to Ref. , a much lower NMR frequency $`\nu _040`$ MHz was used to ensure a better suppression of the Ga ($`4e`$) contribution. This is achieved by taking advantage of the quadrupole line broadening ($`\nu _Q^2/\nu _0`$) which is strongly different for the two sites which have a different local charge environment ($`{}_{}{}^{71}\nu _{Q}^{}(4f,4e)2.9,20.5`$ MHz). For $`220<T<450`$ K, the NMR specta were taken by sweeping the frequency in a fixed 7.5 T field. A typical set of the field sweep spectra, recorded around $`3`$T, is reported in Fig. 2. The high $`T`$ part (upper panel), show that upon cooling the NMR lines shift to the right (lower $`H`$), without any appreciable broadening. In contrast, the low $`T`$ part (lower panel), shows that upon further cooling the lines shift to the left (higher $`H`$), and broaden. This crossover (taking place at $`50`$ K) in the $`T`$ \- dependence of the shift, reflecting the local magnetic susceptibility, is the major finding of this letter. Another feature seen in Fig. 2 is a wipe-out of the intensity due to fast nuclear relaxation when the dynamics of electronic spins slows down. This is common in various systems ranging from spin glasses to AF correlated systems . From Ref. and the $`\mu `$SR relaxation data reported on the same compound , one would naively expect this effect to occur in SCGO only in the vicinity of $`T_g`$. Carefully measured integrated intensity of the <sup>69,71</sup>Ga spectra shows no variation above 15 K (fig.3), clearly demonstrating that our data very reliably reflects the behavior of all the electronic spins in the system, including the $`T`$-range 15-50 K where the shift direction changes. On the contrary, it is worth noticing that below $`T=10`$ K$`3T_g`$, more than 50% of the sites are wiped out of our experimental window, indicating the occurrence of an inhomogeneous dynamics of the spin system in an unusually high $`T`$\- range as compared to $`T_g`$. The shift $`K`$ and the width $`\mathrm{\Delta }H`$ are extracted from the NMR line at all $`T`$. $`K`$ is related to the average field at the Ga($`4f`$) site, and directly probes the (homogeneous) susceptibility of the Cr$`(12k)`$ and Cr($`2a`$) ions. $`\mathrm{\Delta }H`$ originates from a distribution of internal fields on the nuclear Ga($`4f`$) site, which is naturally associated with an inhomogeneous susceptibility of the Cr spin system. For $`T>120`$ K, the line broadening is small enough that, for practical purpose, we extract $`K`$ from the shift of the line edges. This method is not adequate at lower $`T`$. Below 120 K, the line is symmetrically broadened, and to deduce $`K`$, we used either the centre of gravity or a partial Gaussian fit of the Ga(4f) contribution. Independently of the type of analysis, $`K`$ was found to decrease at low $`T`$. $`\mathrm{\Delta }H`$ was extracted from Gaussian fits. First we discuss the temperature dependence of $`K`$ which is presented in fig. 4, where we also include results taken at various applied fields. From the high $`T`$-data (inset), we can extract a Néel temperature $`\theta _{NMR}470`$ K of the same order as $`\theta _{macro}=560`$ K. This confirms that $`K`$ reflects the physics of the frustrated unit. As mentioned before, $`K`$ first increases with decreasing $`T`$ down to $`50`$ K, but below, $`K`$ flattens and even decreases by $`20`$%. The sharp contrast between the temperature dependence of $`K`$ and $`\chi _{macro}`$, below $`50`$ K, is emphasized by the dashed arrow in the figure. It reveals that 2 different types of Cr have to be considered. In other words, our shift data rule out models which attempt to associate the low-$`T`$ macroscopic susceptibility only with a generic - therefore homogeneous- property of the frustrated lattice. Further investigations, to be detailed elsewhere , clearly confirm that the shift variation reported for this sample is an intrinsic feature of the frustrated network as it is very little dependent on the Cr/Ga substitution (at variance with the width). Next we discuss the variation of $`\mathrm{\Delta }H`$ at low $`T`$. The results, taken for various frequencies, are summarized in fig.5. At low $`T`$, the broadening scales remarkably with the applied field for both isotopes. This clearly confirms the magnetic origin of the width at low $`T`$. A Curie-like behavior is found for $`\mathrm{\Delta }H(T)`$, as shown by the solid line. We therefore plot in the inset $`\mathrm{\Delta }H/H_0`$ versus $`\chi _{macro}`$ measured for the same $`H_0`$3 Tesla field, using $`T`$ as an implicit parameter. The linearity of the relationship between $`\mathrm{\Delta }H`$ and $`\chi _{macro}`$ strongly suggests that the Curie upturn, which dominates $`\chi _{macro}`$ at low $`T`$, and the linewidth have a common origin of inhomogeneous magnetism. The deviation between $`\chi _{macro}`$ and $`K`$ below $`T=50`$ K is also straightforwardly explained by this viewpoint. In summary, our NMR results are consistent with a picture where $`\chi _{macro}`$ is a sum of two distinguished components $`\chi _{frust}`$ and $`\chi _{def}`$. $`\chi _{frust}`$ is the homogeneous susceptibility reflected in $`K`$ and representing the physics of the kagomé-based lattice. It has a Curie-Weiss like behavior at high $`T`$ and displays a crossover at $`50`$ K to a pseudogap behavior. $`\chi _{def}`$ is the inhomogeneous contribution to the susceptibility reflected in $`\mathrm{\Delta }H`$ and originating from defects of the frustrated block. This component has a pure Curie low-$`T`$ contribution and it dominates $`\chi _{macro}`$ at $`TT_g^+`$. We now turn to discuss our experimental results in light of existing theories. Susceptibility of the pure Kagomé network has been numerically simulated using various models. In many of them, such as e.g. the case of singlets formation , a gap $`\mathrm{\Delta }`$ appears in $`\chi _{frust}`$ and $`\chi _{frust}e^{\mathrm{\Delta }/T}`$ at low $`T`$. Using for $`\mathrm{\Delta }`$ the temperature $`T_{\mathrm{max}}=50`$ K where $`K`$ peaks, one would expect a much sharper decrease of $`\chi _{frust}`$ than the 20% decrease observed at $`20`$ K. The discrepancy between the theoretically expected and measured decrease of $`K`$ might be solved using a more realistic model of pyrochlore slab . In this model spins from the triangular Cr($`2a`$) layer combine with the Kagomé Cr($`12k`$) to generate a basic unit with an uncompensated moment. This moment is expected to add a $`1/T`$ homogeneous contribution to $`\chi _{frust}`$, which should weaken the drop of $`K`$ below $`\mathrm{\Delta }`$. Whether such a term somewhat counterbalances the effect of the gap on the measured $`K`$ is still speculative as, unfortunately, an experimental confirmation is prevented by the loss of NMR intensity below 15 K. Therefore, we cannot definitely conclude on the full opening of a gap at the present stage, hence the name pseudogap. Regardless of the nature of the gap, the value of $`T_{\mathrm{max}}`$ is very surprising. In most models $`T_{\mathrm{max}}<0.1J`$ where $`J`$ $``$ $`100K`$ is the exchange interaction. Here $`T_{\mathrm{max}}J/2`$, which is much bigger than expected and obviously further theories are required to explain in detail our results. To our knowledge, only a chiral model features a peak in $`\chi `$ at $`T`$ as high as $`0.4J`$ . From the absence of neutron signature, one does not expect any real magnetic order to occur around 50K. Our shift data would rather indicate an increase of the magnetic correlations peaked at the chiral wave-vector. This scenario resembles the case of the pseudogap in High $`T_c`$ cuprates. An even simpler interpretation to the high value of $`T_{\mathrm{max}}`$ relies on the fact that for any low dimensional AF correlated system, one expects the susceptibility to decrease at low $`T`$. The crossover usually occurs in non-frustrated $`2D`$ AF for $`T\theta `$, however, because of frustration, it could occur at lower temperatures, here 10 times smaller than $`\theta `$. The ratio $`\theta /T_{\mathrm{max}}`$ might, finally, prove to be a better characterization of the degree of frustration than using $`\theta `$ $`/T_g`$ since the origin of the spin glass freezing might be associated with defects, as discussed below. We now propose an interpretation for the origin of the line width in the light of the model developed in . There, the origin of the $`1/T`$ paramagnetic behavior of $`\chi _{macro}`$ is assigned to the existence of triangles of the Kagomé lattice non fully occupied by Cr<sup>3+</sup> moments. The substitution of two adjacent Cr<sup>3+</sup> sites by Ga seems necessary to generate a paramagnetic-like “defect” at low $`T`$. A priori such a paramagnetic defect could lead to a well defined feature in the spectrum and a broadening, depending on the response of the electronic spin system to this defect. The number of Ga sites directly coupled to these $`(1p)^2=1\%`$ triangles is small and the corresponding signal is thus likely unobservable. On the contrary, a staggered response (sign oscillation of the field generated by the defect as a function of distance) over few lattice constants, is expected to lead to a symmetric line broadening of the full line. This phenomenon is observed in a large number of systems such as High $`T_c`$ cuprates, or 1D spin chains and ladders, where a similar low-$`T`$ increase of the NMR linewidth is reported. Therefore, we conclude that the defects in SCGO must be coupled to the surrounding correlated spins, in agreement with the idea of Ref. . For a quantitative analysis of the low-$`T`$ contribution of Ga/Cr substitutions to $`\chi _{macro}`$, one needs to subtract the contribution from an ideal pure sample, unfortunately not stable. Nevertheless, we use a simple (and consistent) viewpoint where $`\chi _{macro}`$ is dominated by the substitution defects at low $`T`$. From the value of the low temperature Curie constant $`C_{LT}`$ deduced from $`\chi _{macro}`$, we can deduce the value of the effective moment associated with one defect, $`\mu _{eff}`$, provided the number of defects, $`N_{defect}`$ , is known. We follow and write $`N_{defect}/N_{Cr}=3/2(1p)^2`$, where $`N_{Cr}`$ is the total number of Cr. From $`C_{LT}=N_{defect}\mu _{eff}^2/3k_B=0.03`$ emu/mol , we find $`\mu _{eff}(`$defect$`)4\mu _B`$, typical of a spin $`3/2`$. This reminds a similar case for $`S=1/2`$ AF cuprates where the absence of spin in the square 2D network generates a staggered damped response of the surrounding spins with a total moment corresponding to a spin 1/2. The overall consistency with the model of Ref. is encouraging, but, of course, more NMR and susceptibility experiments are needed for other low substitution rates, in order to further check the quadratic concentration dependence of $`N_{defect}`$. In conclusion, we have demonstrated that the intrinsic Kagomé/pyrochlore slab susceptibility displays a broad maximum around $`TJ/2`$. For $`T<20`$K, our data suggest that the macroscopic susceptibility is dominated by the contribution from defects which remain coupled to the frustrated network. Finally, the occurrence of a slowing down of spin fluctuations is clearly evidenced below 15K. Our results definitely set new constraints on the theoretical models and are stimulating for other NMR studies in the broad class of frustrated systems. We acknowledge Y.J. Uemura who suggested this work and C. L’Huillier, F. Mila, H. Alloul, J. Bobroff for fruitful discussions.
warning/0004/gr-qc0004022.html	ar5iv	text	# THE RENORMALIZED STRESS TENSOR IN KERR SPACE-TIME: GENERAL RESULTS ## I Introduction One of the central quantities of physical interest in a study of quantum field theory in curved space-time is the renormalized expectation value of the stress-energy tensor (RSET), since it is this quantity which couples, via the semi-classical Einstein equations, to the background geometry and thus yields the one-loop correction dynamics of the geometry. This paper is devoted to the properties of the RSET in the states of greatest physical interest on Kerr space-times. Any assault on it by direct computation in black hole geometries is invariably a long and complex process, requiring much algebraic dexterity and ingenuity, and usually resorting to numerical analysis via computer. The aim in this paper is to present what information can be gathered from more physical principles and general considerations. The most important of these are the symmetries of the space-time and states together with the conservation equations. In addition, various restrictions on the form of the RSET follow from its behaviour at the event horizon and far from the black hole. In a subsequent paper we shall present numerical results for the RSET in the states appropriate to a Kerr black hole with and without a bounding ‘box’. The contents of this paper are as follows. In section II we briefly review the solution of the wave equation in Kerr space-time, concentrating for simplicity on the case of a conformally coupled, massless scalar field. We also introduce the standard definitions of the Boulware and Unruh vacua, and discuss the subtleties of defining the Hartle-Hawking vacuum in Kerr. In the absence of a true Hartle-Hawking state, we define two possible candidates. Next, in section IV we investigate how much information can be gathered about the stress tensor using the conservation equations, symmetries of the geometry, and regularity conditions on sections of the event horizon. This greatly reduces the number of unknown functions in the stress tensor. The analysis of this section is applicable to any quantum field, and any of the physical vacua. In section V we consider the properties of the physical vacua in the asymptotic regions, at the event horizon and at infinity, again concentrating on the massless scalar field. We calculate the differences in expectation values of the stress tensor in the Unruh vacuum and other states, which can be calculated without renormalization. These calculations are in exact agreement with our earlier analysis. We also discuss the properties of the candidate Hartle-Hawking states, in particular their symmetry and regularity on the event horizon. We follow the space-time conventions of Misner, Thorne and Wheeler and work in geometric units throughout. ## II The wave equation in Kerr space-time The Kerr line element in Boyer-Lindquist co-ordinates has the form $$ds^2=\frac{\mathrm{\Delta }}{\rho ^2}(dta\mathrm{sin}^2\theta d\varphi )^2+\frac{\mathrm{sin}^2\theta }{\rho ^2}\left((r^2+a^2)d\varphi adt\right)^2+\frac{\rho ^2}{\mathrm{\Delta }}dr^2+\rho ^2d\theta ^2$$ (1) where $`\rho ^2=r^2+a^2\mathrm{cos}^2\theta `$ and $`\mathrm{\Delta }=r^22Mr+a^2`$. Here $`M`$ is the mass of the black hole and $`a`$ its angular momentum per unit mass as viewed from infinity. The metric possesses two coordinate singularities at the roots of the equation $`\mathrm{\Delta }=0`$, which we label $`r=r_+=M+(M^2a^2)^{1/2}`$, defining the outer event horizon and $`r=r_{}=M(M^2a^2)^{1/2}`$, defining the inner Cauchy horizon. In addition, there is a curvature singularity on the ring defined by the equation $`\rho ^2=0`$ (corresponding to $`r=0`$ and $`\theta =\pi /2`$). The space-time is stationary and axisymmetric, possessing two Killing vectors, $`\zeta =/t`$ and $`\eta =/\varphi `$. The former is timelike at infinity but becomes null when $`r=r_s=M+\sqrt{M^2a^2\mathrm{cos}^2\theta }`$. This surface is known as the stationary limit surface and between it and the event horizon is a region called the ergosphere. Within the ergosphere, $`\zeta `$ is spacelike and it is impossible for observers to remain at rest with respect to infinity. The stationary limit surface is timelike except on the axis of symmetry $`\theta =0`$, where it joins the event horizon and becomes null. The Killing vector $`\zeta +\mathrm{\Omega }_+\eta `$, where $`\mathrm{\Omega }_+=a/(r_+^2+a^2)=a/2Mr_+`$ is the angular velocity of the event horizon, generates the Killing horizon at $`r=r_+`$. This Killing vector is null on the event horizon, and timelike outside it up to the velocity of light surface, at which point it becomes null again. The velocity of light surface is the surface at which an observer with angular velocity $`\mathrm{\Omega }_+`$ must move with the speed of light. It is not the same as the stationary limit surface. In addition, the space-time possesses a Killing-Yano tensor which we shall discuss later. Consider a conformally coupled massless scalar field satisfying the equation $`_\mu (g^{\frac{1}{2}}g^{\mu \nu }_\nu )\mathrm{\Phi }=0`$ (the scalar curvature $`R`$ being zero in Kerr space-time). This equation is separable in the Kerr metric and the basis functions may be taken to be $$u_{\omega lm}(x)=\frac{N_{\omega lm}}{(r^2+a^2)^{\frac{1}{2}}}e^{i\omega t+im\varphi }S_{\omega lm}(\mathrm{cos}\theta )R_{\omega lm}(r)$$ (2) where $`N_{\omega lm}`$ is a normalization constant, $`l`$ and $`m`$ are integers with $`\|m\|l`$. $`N_{\omega lm}`$ is determined so that our mode functions are orthonormal with respect to the standard inner product $$u_1,u_2=\frac{1}{2}i_\mathrm{\Sigma }\sqrt{g}(u_{2,\mu }^{}u_1u_2^{}u_{1,\mu })𝑑\mathrm{\Sigma }^\mu $$ (3) where $`\mathrm{\Sigma }`$ is any Cauchy hypersurface. $`S_{\omega lm}(\xi )`$ is a spheroidal harmonic satisfying the eigenvalue equation $$\left[\frac{d}{d\xi }(1\xi ^2)\frac{d}{d\xi }\frac{m^2}{1\xi ^2}+2ma\omega (a\omega )^2(1\xi ^2)+\lambda _{lm}(a\omega )\right]S_{\omega lm}(\xi )=0$$ (4) subject to regularity at $`\xi =\pm 1`$. The eigenvalue $`\lambda _{lm}(a\omega )`$ depends on the integers $`l`$ and $`m`$ and has the known value $`\lambda _{lm}(0)=l(l+1)`$, with $`S_{0lm}(\xi )`$ simply an associated Legendre function. We may normalize the spheroidal harmonics so that $$_1^1S_{\omega lm}(\xi )S_{\omega l^{}m}(\xi )𝑑\xi =\delta _{ll^{}}.$$ (5) The radial equation may be written in the form of a 1-dimensional time-independent Schrödinger equation $$\left[\frac{d^2}{dr_{}^2}V_{\omega lm}(r)\right]R_{\omega lm}(r)=0$$ (6) where $$V_{\omega lm}(r)=\left(\omega \frac{ma}{r^2+a^2}\right)^2+\lambda _{lm}(a\omega )\frac{\mathrm{\Delta }}{(r^2+a^2)^2}+\frac{2(Mra^2)\mathrm{\Delta }}{(r^2+a^2)^3}+\frac{3a^2\mathrm{\Delta }^2}{(r^2+a^2)^4},$$ (7) and the ‘tortoise’ co-ordinate $`r_{}`$ is defined as $$r_{}=\frac{r^2+a^2}{\mathrm{\Delta }}𝑑r=r+\frac{1}{2\kappa _+}\mathrm{log}\|rr_+\|+\frac{1}{2\kappa _{}}\mathrm{log}\|rr_{}\|,$$ (8) with $$\kappa _\pm =\frac{r_\pm r_{}}{2(r_\pm ^2+a^2)},$$ (9) being the surface gravity on the inner and outer horizons. In the asymptotic regions $`rr_+`$ ($`r_{}\mathrm{}`$) and $`r\mathrm{}`$ ($`r_{}\mathrm{}`$) the potential (7) reduces to $$V_{\omega lm}(r)\{\begin{array}{cc}(\omega m\mathrm{\Omega }_+)^2\hfill & \text{as }r_{}\mathrm{}\hfill \\ \omega ^2\hfill & \text{as }r_{}\mathrm{}.\hfill \end{array}$$ (10) We may thus choose as a basis of solutions to Eq. (6), two classes of solutions with the asymptotic forms $`R_{\omega lm}^{}(r)`$ $``$ $`\{\begin{array}{cc}e^{i\stackrel{~}{\omega }r_{}}+A_{\omega lm}^{}e^{i\stackrel{~}{\omega }r_{}}\hfill & r_{}\mathrm{}\hfill \\ B_{\omega lm}^{}e^{i\omega r_{}}\hfill & r_{}\mathrm{}\hfill \end{array}`$ (11) $`R_{\omega lm}^+(r)`$ $``$ $`\{\begin{array}{cc}B_{\omega lm}^+e^{i\stackrel{~}{\omega }r_{}}\hfill & r_{}\mathrm{}\hfill \\ e^{i\omega r_{}}+A_{\omega lm}^+e^{i\omega r_{}}\hfill & r_{}\mathrm{}\hfill \end{array}`$ (12) where $`\stackrel{~}{\omega }=\omega m\mathrm{\Omega }_+`$. In the language of the Schrödinger equation analogy it is natural to speak of $`A`$ and $`B`$ as the ‘reflection’ and ‘transmission’ coefficients, respectively. The eigenvalues $`\lambda _{lm}`$ are real and hence if $`R`$ is a solution of Eq. (6) then so too is $`R^{}`$. Using this and the constancy of the Wronskian for solutions to Eq. (6) for various combinations of the radial wavefunctions, it can be shown that following relations hold : $`1\|A_{\omega lm}^+\|^2`$ $`=`$ $`{\displaystyle \frac{\omega m\mathrm{\Omega }_+}{\omega }}\|B_{\omega lm}^+\|^2`$ (14) $`1\|A_{\omega lm}^{}\|^2`$ $`=`$ $`{\displaystyle \frac{\omega }{\omega m\mathrm{\Omega }_+}}\|B_{\omega lm}^{}\|^2`$ (15) $`\omega B_{\omega lm}^{}{}_{}{}^{}A_{\omega lm}^+`$ $`=`$ $`(\omega m\mathrm{\Omega }_+)B_{\omega lm}^+A_{\omega lm}^{}{}_{}{}^{}`$ (16) $`\omega B_{\omega lm}^{}`$ $`=`$ $`(\omega m\mathrm{\Omega }_+)B_{\omega lm}^+.`$ (17) The first two of these relations show that for $`\omega >0`$, $`\omega m\mathrm{\Omega }_+=\stackrel{~}{\omega }<0`$, both $`\|A^{}\|^2`$ and $`\|A^+\|^2`$ are greater than 1. ## III Quantum field theory in Kerr space-time ### A The mode functions We start by considering two natural complete, orthonormal sets of solutions to the Klein-Gordon equation. It is then straightforward to construct states with particular properties along a given Cauchy surface, for example $`^{}^{}`$. Later, we shall address the much more difficult question of constructing states characterized on surfaces which do not form a Cauchy surface, for example $`^{}^+`$. With the understanding that $`\omega >0`$, we take as the ‘past’ basis the following : $`u_{\omega lm}^{in}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{8\pi ^2\omega (r^2+a^2)}}}e^{i\omega t}e^{im\varphi }S_{\omega lm}(\mathrm{cos}\theta )R_{\omega lm}^+(r)\stackrel{~}{\omega }>m\mathrm{\Omega }_+`$ (19) $`u_{\omega lm}^{up}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{8\pi ^2\stackrel{~}{\omega }(r^2+a^2)}}}e^{i\omega t}e^{im\varphi }S_{\omega lm}(\mathrm{cos}\theta )R_{\omega lm}^{}(r)\stackrel{~}{\omega }>0`$ (20) $`u_{\omega lm}^{up}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{8\pi ^2(\stackrel{~}{\omega })(r^2+a^2)}}}e^{i\omega t}e^{im\varphi }S_{\omega lm}(\mathrm{cos}\theta )R_{\omega lm}^{}(r)0>\stackrel{~}{\omega }>m\mathrm{\Omega }_+`$ (21) where we have used the property $`S_{\omega lm}(\mathrm{cos}\theta )=S_{\omega lm}(\mathrm{cos}\theta )`$. These modes are orthonormal in the sense that $`(u_{\omega lm}^{in},u_{\omega ^{}l^{}m^{}}^{in})`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}\stackrel{~}{\omega }>m\mathrm{\Omega }_+[\omega >0]`$ (23) $`(u_{\omega lm}^{up},u_{\omega ^{}l^{}m^{}}^{up})`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}\stackrel{~}{\omega }>0[\omega >m\mathrm{\Omega }_+]`$ (24) $`(u_{\omega lm}^{up},u_{\omega ^{}l^{}m^{}}^{up})`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}0>\stackrel{~}{\omega }>m\mathrm{\Omega }_+[m\mathrm{\Omega }_+>\omega >0]`$ (25) with all other inner products vanishing. Our conventions here adhere to those of the ‘distant observer viewpoint’ of Frolov and Thorne which we will follow consistently throughout this series of papers. From Eq. (12), $`u_{\omega lm}^{in}`$ $``$ $`{\displaystyle \frac{S_{\omega lm}(\mathrm{cos}\theta )}{\sqrt{8\pi ^2\omega (r^2+a^2)}}}\times \{\begin{array}{cc}0\hfill & \text{at }^{}\hfill \\ \mathrm{exp}(i\omega v+im\varphi )\hfill & \text{at }^{}\hfill \\ B_{\omega lm}^+\mathrm{exp}(i\stackrel{~}{\omega }v+im\varphi _+)\hfill & \text{at }^+\hfill \\ A_{\omega lm}^+\mathrm{exp}(i\omega u+im\varphi )\hfill & \text{at }^+\hfill \end{array}\begin{array}{c}\stackrel{~}{\omega }>m\mathrm{\Omega }_+\hfill \\ \left[\omega >0\right]\hfill \end{array}`$ (29) $`u_{\omega lm}^{up}`$ $``$ $`{\displaystyle \frac{S_{\omega lm}(\mathrm{cos}\theta )}{\sqrt{8\pi ^2\stackrel{~}{\omega }(r^2+a^2)}}}\times \{\begin{array}{cc}\mathrm{exp}(i\stackrel{~}{\omega }u+im\varphi _+)\hfill & \text{at }^{}\hfill \\ 0\hfill & \text{at }^{}\hfill \\ A_{\omega lm}^{}\mathrm{exp}(i\stackrel{~}{\omega }v+im\varphi _+)\hfill & \text{at }^+\hfill \\ B_{\omega lm}^{}\mathrm{exp}(i\omega u+im\varphi )\hfill & \text{at }^+\hfill \end{array}\begin{array}{c}\stackrel{~}{\omega }>0\hfill \\ \left[\omega >m\mathrm{\Omega }_+\right]\hfill \end{array}`$ (32) $`u_{\omega lm}^{up}`$ $``$ $`{\displaystyle \frac{S_{\omega lm}(\mathrm{cos}\theta )}{\sqrt{8\pi ^2\|\stackrel{~}{\omega }\|(r^2+a^2)}}}\times \{\begin{array}{cc}\mathrm{exp}(i\|\stackrel{~}{\omega }\|uim\varphi _+)\hfill & \text{at }^{}\hfill \\ 0\hfill & \text{at }^{}\hfill \\ A_{\omega lm}^{}\mathrm{exp}(i\|\stackrel{~}{\omega }\|vim\varphi _+)\hfill & \text{at }^+\hfill \\ B_{\omega lm}^{}\mathrm{exp}(i\omega uim\varphi )\hfill & \text{at }^+\hfill \end{array}\begin{array}{c}0>\stackrel{~}{\omega }>m\mathrm{\Omega }_+\hfill \\ \left[m\mathrm{\Omega }_+>\omega >0\right]\hfill \end{array}`$ (35) where $$u=tr_{},v=t+r_{},\varphi _+=\varphi \mathrm{\Omega }_+t.$$ (37) These modes are natural to the initial surfaces $`^{}`$ and $`^{}`$ in the sense that $`u^{in}`$ describes unit flux coming in from $`^{}`$ and zero flux coming up from $`^{}`$, whereas $`u^{up}`$ describes unit flux coming up from $`^{}`$ and zero incoming flux coming in from $`^{}`$. For modes with $`\stackrel{~}{\omega }<0`$ (but $`\omega >0`$), $`\|A^{}\|^2>1`$, so that they are reflected back to $`^+`$ with an amplitude greater than that they had originally at $`^{}`$. This is the classical phenomenon of superradiance. Of course, as $`\omega >0`$ and $`\mathrm{\Omega }_+>0`$ it is only possible for $`\stackrel{~}{\omega }=\omega m\mathrm{\Omega }_+`$ to be negative if $`m>0`$, that is for co-rotating waves. Corresponding comments apply to in modes with $`\stackrel{~}{\omega }<0`$: they are reflected back to $`^+`$ with an amplitude greater than that they had originally at $`^{}`$. Another aspect of superradiance is important to our discussion. From Eq. (LABEL:eq:null\_asympupS), one sees that the up modes (21) with $`\stackrel{~}{\omega }<0`$ have a negative energy wave propagating to $`^+`$ (conservation of energy). This is a consequence of $`_t`$ not being a globally time-like Killing vector. $`_t`$ is space-like in the ergosphere, however the combination $`_t+\mathrm{\Omega }_\varphi `$, where $`\mathrm{\Omega }=g_{t\varphi }/g_{\varphi \varphi }`$, is time-like down to the horizon upon which it becomes null. Observers following integral curves of this time-like vector field are locally non-rotating observers (LNRO). A LNRO near the horizon would measure the frequency of the superradiant up modes in (21) to be $`\|\stackrel{~}{\omega }\|=\stackrel{~}{\omega }=\omega +m\mathrm{\Omega }_+`$, in particular, the LNRO would see positive frequency waves for all modes. For $`u_{\omega lm}^{in}`$ all modes are positive frequency at $`^+`$ and $`^{}`$. A LNRO near the horizon measures $`\stackrel{~}{\omega }`$ for the frequency and thus sees negative frequency modes in the superradiant regime. An up mode having positive frequency with respect to $`u`$ at $`^{}`$ will have negative frequency with respect to $`u`$ at $`^+`$ if $`\stackrel{~}{\omega }<0`$ but $`\omega >0`$. ### B The Physical Vacua We now turn to the delicate issue of defining analogs of the standard three vacuum states in Schwarzschild space-time (Boulware, Hartle-Hawking and Unruh) in Kerr space-time. (Our discussion here concerns states on the full exterior region of Kerr, in later papers we shall also talk about the case when the black hole is contained within a ‘box’.) The construction of vacuum states in Kerr is a more subtle problem than for Schwarzschild black holes, for the following reasons: 1. The existence of superradiant modes makes the definition of positive frequency more complicated. For example, in Schwarzschild, an outgoing mode which has positive frequency with respect to the retarded null co-ordinate $`u`$ at the past horizon $`^{}`$ will also have positive frequency with respect to $`u`$ at $`^+`$, so it does not matter if we define positive frequency with respect to $`u`$ at $`^{}`$ or at $`^+`$. This is no longer the case in Kerr: a superradiant mode can have positive frequency with respect to $`u`$ at $`^+`$ but negative frequency at $`^{}`$. This is why our definition of the basis of mode functions (III A) had to be so carefully done. 2. As a consequence of this, it is only straightforward to define states with particular properties along a given Cauchy surface, such as $`^{}^{}`$. By contrast, it has become conventional in Schwarzschild space-time to consider the Boulware vacuum in terms of its properties on $`^{}^+`$ and the Hartle-Hawking vacuum in terms of its properties on $`^{}^+`$. To be explicit, we may expand the scalar field $`\mathrm{\Phi }(x)`$ in terms of the mode functions we introduced above $`\mathrm{\Phi }(x)`$ $`=`$ $`{\displaystyle \underset{l,m}{}}\left({\displaystyle _0^{\mathrm{}}}𝑑\omega (a_{\omega lm}^{in}u_{\omega lm}^{in}+a_{\omega lm}^{in}u_{\omega lm}^{in})+{\displaystyle _{\omega _{min}}^{\mathrm{}}}𝑑\omega (a_{\omega lm}^{up}u_{\omega lm}^{up}+a_{\omega lm}^{up}u_{\omega lm}^{up})\right)`$ (39) $`+{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\omega _{min}}}𝑑\omega (a_{\omega lm}^{up}u_{\omega lm}^{up}+a_{\omega lm}^{up}u_{\omega lm}^{up})`$ $`=`$ $`{\displaystyle \underset{l,m}{}}\left({\displaystyle _0^{\mathrm{}}}𝑑\omega (a_{\omega lm}^{in}u_{\omega lm}^{in}+a_{\omega lm}^{in}u_{\omega lm}^{in})+{\displaystyle _0^{\mathrm{}}}𝑑\stackrel{~}{\omega }(a_{\omega lm}^{up}u_{\omega lm}^{up}+a_{\omega lm}^{up}u_{\omega lm}^{up})\right)`$ (40) where $`\omega _{min}=\mathrm{max}\{0,m\mathrm{\Omega }_+\}`$, so $`\omega _{min}=0`$ for counter-rotating waves ($`m0`$) and $`\omega _{min}=m\mathrm{\Omega }_+`$ for co-rotating waves ($`m>0`$). Given this expansion, the natural way to quantize the field is for the coefficients to become operators satisfying the commutation relations $`[\widehat{a}_{\omega lm}^{in},\widehat{a}_{\omega ^{}l^{}m^{}}^{in}]`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}\stackrel{~}{\omega }>m\mathrm{\Omega }_+`$ (42) $`[\widehat{a}_{\omega lm}^{up},\widehat{a}_{\omega ^{}l^{}m^{}}^{up}]`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}\stackrel{~}{\omega }>0`$ (43) $`[\widehat{a}_{\omega lm}^{up},\widehat{a}_{\omega ^{}l^{}m^{}}^{up}]`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}0>\stackrel{~}{\omega }>m\mathrm{\Omega }_+`$ (44) with all other commutators vanishing. From Eq.(III A), the operators $`\widehat{a}^{in}`$ and $`\widehat{a}^{up}`$ have the natural interpretation that they will, respectively, create particles incident from $`^{}`$ and $`^{}`$. With this in mind, we define a ‘past Boulware’ vacuum state by $`\widehat{a}_{\omega lm}^{in}\|B^{}`$ $`=`$ $`0\stackrel{~}{\omega }>m\mathrm{\Omega }_+`$ (46) $`\widehat{a}_{\omega lm}^{up}\|B^{}`$ $`=`$ $`0\stackrel{~}{\omega }>0`$ (47) $`\widehat{a}_{\omega lm}^{up}\|B^{}`$ $`=`$ $`00>\stackrel{~}{\omega }>m\mathrm{\Omega }_+`$ (48) corresponding to an absence of particles from $`^{}`$ and $`^{}`$. This state does not precisely correspond to the idea of a Boulware state in Schwarzschild as that state which is most empty at infinity. The state $`\|B^{}`$ contains, at $`^+`$, an outward flux of particles in the superradiant modes; this is the Unruh-Starobinskii effect . One might suppose that a more appropriate definition for the Boulware vacuum would be to define a state which is empty at $`^{}`$ and $`^+`$. However, it is straightforward to see that such a state cannot exist within conventional quantum field theory by introducing the mode functions natural for defining the ‘future Boulware’ vacuum. (We shall discuss later the non-conventional ‘$`\eta `$-formalism’ construction proposed by Frolov and Thorne.) The mode functions relevant to the ‘future Boulware’ vacuum are those representing a unit (locally-positive frequency) flux out to $`^+`$ and down $`^+`$. From the asymptotic forms for the radial functions Eq. (12), it is clear that we should take as our ‘future’ basis : $`u_{\omega lm}^{out}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{8\pi ^2\omega (r^2+a^2)}}}e^{i\omega t}e^{im\varphi }S_{\omega lm}(\mathrm{cos}\theta )R_{\omega lm}^+(r)\stackrel{~}{\omega }>m\mathrm{\Omega }_+,`$ (50) $`u_{\omega lm}^{down}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{8\pi ^2\stackrel{~}{\omega }(r^2+a^2)}}}e^{i\omega t}e^{im\varphi }S_{\omega lm}(\mathrm{cos}\theta )R_{\omega lm}^{}(r)\stackrel{~}{\omega }>0,`$ (51) $`u_{\omega lm}^{down}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{8\pi ^2\|\stackrel{~}{\omega }\|(r^2+a^2)}}}e^{i\omega t}e^{im\varphi }S_{\omega lm}(\mathrm{cos}\theta )R_{\omega lm}^{}(r)0>\stackrel{~}{\omega }>m\mathrm{\Omega }_+.`$ (52) These modes are orthonormal in the sense that $`(u_{\omega lm}^{out},u_{\omega ^{}l^{}m^{}}^{out})`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}\stackrel{~}{\omega }>m\mathrm{\Omega }_+[\omega >0]`$ (54) $`(u_{\omega lm}^{down},u_{\omega ^{}l^{}m^{}}^{down})`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}\stackrel{~}{\omega }>0[\omega >m\mathrm{\Omega }_+]`$ (55) $`(u_{\omega lm}^{down},u_{\omega ^{}l^{}m^{}}^{down})`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}0>\stackrel{~}{\omega }>m\mathrm{\Omega }_+[m\mathrm{\Omega }_+>\omega >0]`$ (56) with all other inner products vanishing. Their asymptotic properties are given by $`u_{\omega lm}^{out}`$ $``$ $`{\displaystyle \frac{S_{\omega lm}(\mathrm{cos}\theta )}{\sqrt{8\pi ^2\omega (r^2+a^2)}}}\times \{\begin{array}{cc}B_{\omega lm}^+\mathrm{exp}(i\stackrel{~}{\omega }u+im\varphi _+)\hfill & \text{at }^{}\hfill \\ A_{\omega lm}^+\mathrm{exp}(i\omega v+im\varphi )\hfill & \text{at }^{}\hfill \\ 0\hfill & \text{at }^+\hfill \\ \mathrm{exp}(i\omega u+im\varphi )\hfill & \text{at }^+\hfill \end{array}\begin{array}{c}\stackrel{~}{\omega }>m\mathrm{\Omega }_+\hfill \\ \left[\omega >0\right]\hfill \end{array}`$ (60) $`u_{\omega lm}^{down}`$ $``$ $`{\displaystyle \frac{S_{\omega lm}(\mathrm{cos}\theta )}{\sqrt{8\pi ^2\stackrel{~}{\omega }(r^2+a^2)}}}\times \{\begin{array}{cc}A_{\omega lm}^{}\mathrm{exp}(i\stackrel{~}{\omega }u+im\varphi _+)\hfill & \text{at }^{}\hfill \\ B_{\omega lm}^{}\mathrm{exp}(i\omega v+im\varphi )\hfill & \text{at }^{}\hfill \\ \mathrm{exp}(i\stackrel{~}{\omega }v+im\varphi _+)\hfill & \text{at }^+\hfill \\ 0\hfill & \text{at }^+\hfill \end{array}\begin{array}{c}\stackrel{~}{\omega }>0\hfill \\ \left[\omega >m\mathrm{\Omega }_+\right]\hfill \end{array}`$ (63) $`u_{\omega lm}^{down}`$ $``$ $`{\displaystyle \frac{S_{\omega lm}(\mathrm{cos}\theta )}{\sqrt{8\pi ^2\|\stackrel{~}{\omega }\|(r^2+a^2)}}}\times \{\begin{array}{cc}A_{\omega lm}^{}\mathrm{exp}(i\|\stackrel{~}{\omega }\|uim\varphi _+)\hfill & \text{at }^{}\hfill \\ B_{\omega lm}^{}\mathrm{exp}(i\omega vim\varphi )\hfill & \text{at }^{}\hfill \\ \mathrm{exp}(i\|\stackrel{~}{\omega }\|vim\varphi _+)\hfill & \text{at }^+\hfill \\ 0\hfill & \text{ at }^+\hfill \end{array}\begin{array}{c}0>\stackrel{~}{\omega }>m\mathrm{\Omega }_+\hfill \\ \left[m\mathrm{\Omega }_+>\omega >0\right]\hfill \end{array}`$ (66) We may expand the scalar field $`\mathrm{\Phi }(x)`$ in terms of these mode functions we introduced above $`\mathrm{\Phi }(x)`$ $`=`$ $`{\displaystyle \underset{l,m}{}}\left({\displaystyle _0^{\mathrm{}}}𝑑\omega (a_{\omega lm}^{out}u_{\omega lm}^{out}+a_{\omega lm}^{out}u_{\omega lm}^{out})+{\displaystyle _{\omega _{min}}^{\mathrm{}}}𝑑\omega (a_{\omega lm}^{down}u_{\omega lm}^{down}+a_{\omega lm}^{down}u_{\omega lm}^{down})\right)`$ (69) $`+{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\omega _{min}}}𝑑\omega (a_{\omega lm}^{down}u_{\omega lm}^{down}+a_{\omega lm}^{down}u_{\omega lm}^{down})`$ $`=`$ $`{\displaystyle \underset{l,m}{}}\left({\displaystyle _0^{\mathrm{}}}𝑑\omega (a_{\omega lm}^{out}u_{\omega lm}^{out}+a_{\omega lm}^{out}u_{\omega lm}^{out})+{\displaystyle _0^{\mathrm{}}}𝑑\stackrel{~}{\omega }(a_{\omega lm}^{down}u_{\omega lm}^{down}+a_{\omega lm}^{down}u_{\omega lm}^{down})\right)`$ (70) where $`\omega _{min}=\mathrm{max}\{0,m\mathrm{\Omega }_+\}`$, as before. Given this expansion, the natural way to quantize the field is for the coefficients become operators satisfying the commutation relations $`[\widehat{a}_{\omega lm}^{out},\widehat{a}_{\omega ^{}l^{}m^{}}^{out}]`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}\stackrel{~}{\omega }>m\mathrm{\Omega }_+`$ (72) $`[\widehat{a}_{\omega lm}^{down},\widehat{a}_{\omega ^{}l^{}m^{}}^{down}]`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}\stackrel{~}{\omega }>0`$ (73) $`[\widehat{a}_{\omega lm}^{down},\widehat{a}_{\omega ^{}l^{}m^{}}^{down}]`$ $`=`$ $`\delta (\omega \omega ^{})\delta _{ll^{}}\delta _{mm^{}}0>\stackrel{~}{\omega }>m\mathrm{\Omega }_+`$ (74) with all other commutators vanishing. From Eq.(III B), the operators $`\widehat{a}^{out}`$ and $`\widehat{a}^{down}`$ have the natural interpretation that they will, respectively, create particles incident from $`^+`$ and $`^+`$. Thus, we define the ‘future Boulware’ vacuum state by $`\widehat{a}_{\omega lm}^{out}\|B^+`$ $`=`$ $`0\stackrel{~}{\omega }>m\mathrm{\Omega }_+`$ (76) $`\widehat{a}_{\omega lm}^{down}\|B^+`$ $`=`$ $`0\stackrel{~}{\omega }>0`$ (77) $`\widehat{a}_{\omega lm}^{down}\|B^+`$ $`=`$ $`00>\stackrel{~}{\omega }>m\mathrm{\Omega }_+`$ (78) corresponding to an absence of particles from $`^+`$ and $`^+`$. In this language, the Unruh-Starobinskii effect is a statement about the behaviour of $$B^{}\|\widehat{T}_{\mu \nu }\|B^{}B^+\|\widehat{T}_{\mu \nu }\|B^+$$ (79) as $`r\mathrm{}`$. A vacuum state empty at $`^{}`$ and $`^+`$ must be constructed from modes $`u_{\omega lm}^{in}`$ and $`u_{\omega lm}^{out}`$ up to a trivial Bogoliubov transformation (i.e., one with all $`\beta `$-coefficients vanishing). However, $`u_{\omega lm}^{in}`$ and $`u_{\omega lm}^{out}`$ are not orthogonal and the fact that they cannot be made so by any trivial Bogoliubov transformation is seen most easily by writing $`u_{\omega lm}^{out}`$ in terms of the basis given by $`u_{\omega lm}^{in}`$ and $`u_{\omega lm}^{up}`$. For non-superradiant modes the transformation does correspond to a trivial Bogoliubov transformation: $`u_{\omega lm}^{out}`$ $`=`$ $`A_{\omega lm}^+u_{\omega lm}^{in}+\sqrt{{\displaystyle \frac{\stackrel{~}{\omega }}{\omega }}}B_{\omega lm}^+u_{\omega lm}^{up},\stackrel{~}{\omega }>0,`$ (81) $`u_{\omega lm}^{down}`$ $`=`$ $`\sqrt{{\displaystyle \frac{\omega }{\stackrel{~}{\omega }}}}B_{\omega lm}^{}u_{\omega lm}^{in}+A_{\omega lm}^{}u_{\omega lm}^{up},[\omega >m\mathrm{\Omega }_+],`$ (82) but for superradiant modes $`u_{\omega lm}^{out}`$ $`=`$ $`A_{\omega lm}^+u_{\omega lm}^{in}\sqrt{{\displaystyle \frac{\stackrel{~}{\omega }}{\omega }}}B_{\omega lm}^+u_{\omega lm}^{up},0>\stackrel{~}{\omega }>m\mathrm{\Omega }_+\left[m\mathrm{\Omega }_+>\omega >0\right],`$ (84) $`u_{\omega lm}^{down}`$ $`=`$ $`\sqrt{{\displaystyle \frac{\omega }{\stackrel{~}{\omega }}}}B_{\omega lm}^{}u_{\omega lm}^{in}+A_{\omega lm}^{}u_{\omega lm}^{up},0>\stackrel{~}{\omega }>m\mathrm{\Omega }_+\left[m\mathrm{\Omega }_+>\omega >0\right].`$ (85) As no trivial Bogoliubov transformation can affect the total number of ‘particles’ produced, $`_{i,r}\|\beta _{ir}\|^2`$, it is impossible to define a vacuum state empty with respect to in modes at $`^{}`$ and out modes at $`^+`$. The non-existence of a ‘true Boulware’ state is intimately linked with the non-existence a ‘true Hartle-Hawking’ state (defined as being a Hadamard state which respects the symmetries of the space-time and is regular everywhere, in particular, on both future and past event horizons) on Kerr space-time . In the former case, one wishes to define the state on $`^{}^+`$, in the latter on $`^{}^+`$. Indeed, one can make the analogy quite precise by, in the language of Frolov and Thorne, switching from a ‘distant’ to a ‘near horizon’ viewpoint. The (past) Unruh state $`\|U^{}`$ is easily defined as that state empty at $`^{}`$ but with the ‘up’ modes (natural modes on $`^{}`$) thermally populated. For a proof that this is equivalent to using modes which are positive frequency with respect to a future-increasing affine parameter on $`^{}`$ see Ref. . As before, we use the notation $`\|U^{}`$ in order to emphasize that this state is naturally defined by considerations on $`^{}^{}`$. One can, of course also define a state $`\|U^+`$ empty at $`^+`$ but with the ‘down’ modes (natural modes on $`^+`$) thermally populated. Indeed, one can also make such a distinction in the Schwarzschild case for the Unruh vacuum. However, one rarely considers $`\|U^+`$ as it is $`\|U^{}`$ that mimics the state arising at late times from the collapse of a star to a black hole. For this reason we shall usually drop the term ‘past’ but we will retain the terminology $`\|U^{}`$ to make clear that this state is naturally defined in terms of ‘in’ and ‘up’ modes. In this language, the (Kruskal space-time model of the) Hawking effect is a statement about the behaviour of $$U^{}\|\widehat{T}_{\mu \nu }\|U^{}B^+\|\widehat{T}_{\mu \nu }\|B^+$$ (86) as $`r\mathrm{}`$. With these definitions, it is straightforward to write down mode sum expressions for the two-point functions of the field in the past and future Boulware and (past) Unruh vacuum states: $`G_B^{}(x,x^{})`$ $`=`$ $`B^{}\|\widehat{\mathrm{\Phi }}(x)\widehat{\mathrm{\Phi }}(x^{})\|B^{}`$ (88) $`=`$ $`{\displaystyle \underset{l,m}{}}\left({\displaystyle _0^{\mathrm{}}}𝑑\stackrel{~}{\omega }u_{\omega lm}^{up}(x)u_{\omega lm}^{up}(x^{})+{\displaystyle _0^{\mathrm{}}}𝑑\omega u_{\omega lm}^{in}(x)u_{\omega lm}^{in}(x^{})\right)`$ (89) $`G_{B^+}(x,x^{})`$ $`=`$ $`B^+\|\widehat{\mathrm{\Phi }}(x)\widehat{\mathrm{\Phi }}(x^{})\|B^+`$ (90) $`=`$ $`{\displaystyle \underset{l,m}{}}\left({\displaystyle _0^{\mathrm{}}}𝑑\stackrel{~}{\omega }u_{\omega lm}^{down}(x)u_{\omega lm}^{down}(x^{})+{\displaystyle _0^{\mathrm{}}}𝑑\omega u_{\omega lm}^{out}(x)u_{\omega lm}^{out}(x^{})\right)`$ (91) $`G_U^{}(x,x^{})`$ $`=`$ $`U^{}\|\widehat{\mathrm{\Phi }}(x)\widehat{\mathrm{\Phi }}(x^{})\|U^{}`$ (92) $`=`$ $`{\displaystyle \underset{l,m}{}}\left({\displaystyle _0^{\mathrm{}}}𝑑\stackrel{~}{\omega }\text{coth}\left({\displaystyle \frac{\pi \stackrel{~}{\omega }}{\kappa }}\right)u_{\omega lm}^{up}(x)u_{\omega lm}^{up}(x)+{\displaystyle _0^{\mathrm{}}}𝑑\omega u_{\omega lm}^{in}(x)u_{\omega lm}^{in}(x^{})\right).`$ (93) The corresponding expressions for the unrenormalized expectation values of the stress tensor in the past and future Boulware and (past) Unruh vacuum states are: $`B^{}\|\widehat{T}_{\mu \nu }\|B^{}`$ $`=`$ $`{\displaystyle \underset{l,m}{}}\left({\displaystyle _0^{\mathrm{}}}𝑑\stackrel{~}{\omega }T_{\mu \nu }[u_{\omega lm}^{up},u_{\omega lm}^{up}]+{\displaystyle _0^{\mathrm{}}}𝑑\omega T_{\mu \nu }[u_{\omega lm}^{in},u_{\omega lm}^{in}]\right)`$ (95) $`B^+\|\widehat{T}_{\mu \nu }\|B^+`$ $`=`$ $`{\displaystyle \underset{l,m}{}}\left({\displaystyle _0^{\mathrm{}}}𝑑\stackrel{~}{\omega }T_{\mu \nu }[u_{\omega lm}^{down},u_{\omega lm}^{down}]+{\displaystyle _0^{\mathrm{}}}𝑑\omega T_{\mu \nu }[u_{\omega lm}^{out},u_{\omega lm}^{out}]\right)`$ (96) $`U^{}\|\widehat{T}_{\mu \nu }\|U^{}`$ $`=`$ $`{\displaystyle \underset{l,m}{}}\left({\displaystyle _0^{\mathrm{}}}𝑑\stackrel{~}{\omega }\text{coth}\left({\displaystyle \frac{\pi \stackrel{~}{\omega }}{\kappa }}\right)T_{\mu \nu }[u_{\omega lm}^{up},u_{\omega lm}^{up}]+{\displaystyle _0^{\mathrm{}}}𝑑\omega T_{\mu \nu }[u_{\omega lm}^{in},u_{\omega lm}^{in}]\right)`$ (97) where the contribution to the stress-energy tensor, for a massless scalar field mode in Ricci-flat Kerr space-time, assuming conformal coupling, is $$T_{\mu \nu }[u,u^{}]=\frac{1}{3}(u_{;\mu }u_{;\nu }^{}+u_{;\mu }^{}u_{;\nu })\frac{1}{6}(u_{;\mu \nu }u^{}+u_{;\mu \nu }^{}u)\frac{1}{6}g_{\mu \nu }u_{;\tau }u^{;\tau }.$$ (98) Kay and Wald have shown that there does not exist a Hadamard state which respects the symmetries of the space-time and is regular everywhere in Kerr space-time. In the absence of such a ‘true Hartle-Hawking’ vacuum we consider the following states, which are attempts in the literature to define a thermal state with most (but not all) of the properties of the Hartle-Hawking state. The first state is that introduced by Candelas, Chrzanowski and Howard , which is constructed by thermalizing the ‘in’ and ‘up’ modes with respect to their natural energy, so $`G_{CCH}(x,x^{})`$ $`=`$ $`CCH\|\widehat{\mathrm{\Phi }}(x)\widehat{\mathrm{\Phi }}(x^{})\|CCH`$ (99) $`=`$ $`{\displaystyle \underset{l,m}{}}({\displaystyle _0^{\mathrm{}}}d\stackrel{~}{\omega }\mathrm{coth}\left({\displaystyle \frac{\pi \stackrel{~}{\omega }}{\kappa }}\right)u_{\omega lm}^{up}(x)u_{\omega lm}^{up}(x^{})`$ (101) $`+{\displaystyle _0^{\mathrm{}}}d\omega \mathrm{coth}\left({\displaystyle \frac{\pi \omega }{\kappa }}\right)u_{\omega lm}^{in}(x)u_{\omega lm}^{in}(x^{})).`$ and $`CCH\|\widehat{T}_{\mu \nu }\|CCH`$ $`=`$ $`{\displaystyle \underset{l,m}{}}({\displaystyle _0^{\mathrm{}}}d\stackrel{~}{\omega }\mathrm{coth}\left({\displaystyle \frac{\pi \stackrel{~}{\omega }}{\kappa }}\right)T_{\mu \nu }[u_{\omega lm}^{up},u_{\omega lm}^{up}]`$ (103) $`+{\displaystyle _0^{\mathrm{}}}d\omega \mathrm{coth}\left({\displaystyle \frac{\pi \omega }{\kappa }}\right)T_{\mu \nu }[u_{\omega lm}^{in},u_{\omega lm}^{in}]).`$ As such, it might naturally, be described as the ‘past Hartle-Hawking’ vacuum, however, given the discussion above it is not surprising that as we shall show in detail below, this definition gives a state which does not respect the simultaneous $`t`$-$`\varphi `$ reversal invariance of Kerr space-time. The second state we shall consider is that introduced by Frolov and Thorne who used the ‘$`\eta `$ formalism’ to treat the quantization of the superradiant modes. They derived the following expressions in the state, denoted here by $`\|FT`$, which they claim defined the Hartle-Hawking vacuum (at least close to the horizon): $`G_{FT}(x,x^{})`$ $`=`$ $`FT\|\eta \widehat{\mathrm{\Phi }}(x)\eta \widehat{\mathrm{\Phi }}(x^{})\eta \|FT`$ (104) $`=`$ $`{\displaystyle \underset{l,m}{}}({\displaystyle _0^{\mathrm{}}}d\stackrel{~}{\omega }\mathrm{coth}\left({\displaystyle \frac{\pi \stackrel{~}{\omega }}{\kappa }}\right)u_{\omega lm}^{up}(x)u_{\omega lm}^{up}(x^{})`$ (106) $`+{\displaystyle _0^{\mathrm{}}}d\omega \mathrm{coth}\left({\displaystyle \frac{\pi \stackrel{~}{\omega }}{\kappa }}\right)u_{\omega lm}^{in}(x)u_{\omega lm}^{in}(x^{}))`$ and $`FT\|\widehat{T}_{\mu \nu }\|FT`$ $`=`$ $`{\displaystyle \underset{l,m}{}}({\displaystyle _0^{\mathrm{}}}d\stackrel{~}{\omega }\mathrm{coth}\left({\displaystyle \frac{\pi \stackrel{~}{\omega }}{\kappa }}\right)T_{\mu \nu }[u_{\omega lm}^{up},u_{\omega lm}^{up}]`$ (108) $`+{\displaystyle _0^{\mathrm{}}}d\omega \mathrm{coth}\left({\displaystyle \frac{\pi \stackrel{~}{\omega }}{\kappa }}\right)T_{\mu \nu }[u_{\omega lm}^{in},u_{\omega lm}^{in}]).`$ Thus, the Frolov-Thorne state differs in its choice of the appropriate ‘energy’ for the thermal factor corresponding to the ‘in’ modes. This state is formally invariant under simultaneous $`t`$-$`\varphi `$ reversal. Frolov and Thorne claim that the state defined by Eq. (108) is regular out to the speed-of-light surface and is ill-defined outside. However, Kay and Wald’s result is essentially local and the Frolov-Thorne state appears to violate the spirit if not the letter of the result proved by Kay and Wald. Below and in subsequent papers in this series where we address the issues numerically, we shall show that the Frolov-Thorne state is fundamentally flawed while the Candelas-Chrzanowski-Howard state is workable but cannot claim to represent an equilibrium state. ## IV Constraints on the stress tensor We now investigate how much information can be gathered about the stress-energy tensor in Kerr space-time from general physical principles. We shall have in mind the physical vacua which have been defined in the previous section. ### A Solution of the conservation equations In this section, we consider the solution of the conservation equations $`_\nu T_\mu {}_{}{}^{\nu }=0`$. To avoid the calculation of Christoffel symbols, since $`T_{\mu \nu }`$ is a symmetric tensor, the conservation equations may be written in the alternative form $$_\nu (T_\mu {}_{}{}^{\nu }\sqrt{g})=\frac{1}{2}\sqrt{g}(_\mu g_{\alpha \beta })T^{\alpha \beta }$$ (109) where $`g`$ is the determinant of the matrix of metric coefficients given by $`g=\rho ^4\mathrm{sin}^2\theta `$. Since we are interested in the renormalized stress tensor for states which respect the symmetries of the space-time, we assume that the stress-energy tensor, like the metric, is independent of $`t`$ and $`\varphi `$. The $`\mu =t`$ and $`\mu =\varphi `$ equations then become, respectively, $`_r(\rho ^2\mathrm{sin}\theta T_t{}_{}{}^{r})+_\theta (\rho ^2\mathrm{sin}\theta T_t{}_{}{}^{\theta })`$ $`=`$ $`0`$ (110) $`_r(\rho ^2\mathrm{sin}\theta T_\varphi {}_{}{}^{r})+_\theta (\rho ^2\mathrm{sin}\theta T_\varphi {}_{}{}^{\theta })`$ $`=`$ $`0.`$ (111) These may be integrated immediately over $`r`$ to yield $`T_{tr}`$ $`=`$ $`{\displaystyle \frac{K(\theta )}{\mathrm{\Delta }}}{\displaystyle \frac{1}{\mathrm{\Delta }\mathrm{sin}\theta }}_\theta \left(\mathrm{sin}\theta {\displaystyle _{r_+}^r}T_{t\theta }𝑑r^{}\right)`$ (112) $`T_{\varphi r}`$ $`=`$ $`{\displaystyle \frac{L(\theta )}{\mathrm{\Delta }}}{\displaystyle \frac{1}{\mathrm{\Delta }\mathrm{sin}\theta }}_\theta \left(\mathrm{sin}\theta {\displaystyle _{r_+}^r}T_{\varphi \theta }𝑑r^{}\right)`$ (113) where $`K(\theta )`$ and $`L(\theta )`$ are arbitrary functions of $`\theta `$ alone. The $`\mu =r`$ and $`\mu =\theta `$ equations are, respectively, $`F(r,\theta )`$ $`=`$ $`_r(\rho ^2T_r{}_{}{}^{r})+\mathrm{\Delta }^1\mathrm{csc}\theta _\theta (\rho ^2\mathrm{sin}\theta T_\theta {}_{}{}^{r})rT_\theta ^\theta `$ (115) $`\mathrm{\Delta }^1(ra^2\mathrm{sin}\theta \mathrm{\Lambda })T_r^r`$ $`G(r,\theta )`$ $`=`$ $`_r(\rho ^2T_\theta {}_{}{}^{r})+\mathrm{csc}\theta _\theta (\rho ^2\mathrm{sin}\theta T_\theta {}_{}{}^{\theta })`$ (117) $`+a^2\mathrm{sin}\theta \mathrm{cos}\theta T_r{}_{}{}^{r}+a^2\mathrm{sin}\theta \mathrm{cos}\theta T_\theta ^\theta `$ where $`F(r,\theta )`$ $`=`$ $`\rho ^2[\mathrm{\Lambda }T^{tt}+2a\mathrm{\Lambda }\mathrm{sin}^2\theta T^{t\varphi }+\mathrm{sin}^2\theta (\mathrm{\Lambda }a^2\mathrm{sin}^2\theta +r\rho ^4)T^{\varphi \varphi }]`$ (118) $`G(r,\theta )`$ $`=`$ $`{\displaystyle \frac{a^2(r^2+a^2\mathrm{\Delta })}{\rho ^2\mathrm{\Delta }(r^2+a^2)}}\mathrm{sin}\theta \mathrm{cos}\theta [(r^2+a^2)^2T_{tt}+2a(r^2+a^2)T_{t\varphi }`$ (120) $`+a^2T_{\varphi \varphi }]+{\displaystyle \frac{\rho ^2\mathrm{cos}\theta }{(r^2+a^2)\mathrm{sin}^3\theta }}T_{\varphi \varphi }`$ with $`\mathrm{\Lambda }=M(r^2a^2\mathrm{cos}^2\theta )`$. Here we have two equations in six unknowns each of which is a function of two variables $`r`$ and $`\theta `$. One other symmetry immediately apparent from the form of the metric is invariance under the transformation $$\theta \stackrel{~}{\theta }=\pi \theta .$$ (121) The components of the stress-energy tensor will also possess this symmetry, so in particular $$_\theta (T_{\mu \nu })=0\mathrm{when}\theta =\pi /2.$$ (122) This does not imply that any components of $`T_{\mu \nu }`$ vanish, so $`T_{r\theta }`$ is non-zero in general. However, from the conservation equations (117), it follows that $$T_{r\theta }=0\mathrm{when}\theta =\pi /2.$$ (123) The other symmetry of the geometry which should be mentioned here is invariance under simultaneous $`t`$-$`\varphi `$ reversal, that is, $`tt`$ and $`\varphi \varphi `$. The stress tensor for a state satisfying this invariance must have $`T_{tr}=T_{t\theta }=T_{\varphi r}=T_{\varphi \theta }=0`$ and correspondingly $`K(\theta )=L(\theta )=0`$. It might be thought that this simple symmetry of the space-time should be mirrored by the stress tensor for the physical vacua in which we are interested. However, as discussed above this is not the case, because of the superradiant modes. Neither the Boulware vacuum $`\|B^{}`$ nor the Unruh vacuum $`\|U^{}`$ defined in section III B is invariant under simultaneous $`t`$-$`\varphi `$ reversal. This in contrast to the situation for Schwarzschild black holes, where the Boulware vacuum is time-reversal invariant, although the Unruh vacuum is not, due to the Hawking flux. In Schwarzschild space-time, the Hartle-Hawking state is also time-reversal invariant. Of the two Hartle-Hawking-like states, $`\|CCH`$ is not invariant under simultaneous $`t`$-$`\varphi `$ reversal but $`\|FT`$ is. In section V we shall consider further the symmetry and other properties of these states. ### B The trace anomaly As is well-known, conformally invariant field theories on a curved background $`g_{\mu \nu }`$ possess a conformal anomaly which means that the renormalized stress tensor has a trace even though the classical stress tensor must be trace-free. As it arises from the renormalization procedure, the trace anomaly is a geometrical scalar, depending only on the geometry and the nature of the quantum field under consideration, not on the actual quantum state. All methods of regularization agree that it has the form $$\widehat{T}_\alpha ^\alpha _{ren}=k_1C_{\alpha \beta \gamma \delta }C^{\alpha \beta \gamma \delta }+k_2(R_{\alpha \beta }R^{\alpha \beta }\frac{1}{3}R^2)+k_3_\alpha ^\alpha R$$ (124) in four dimensions. Here $`k_1`$, $`k_2`$, $`k_3`$ are constants which are independent of the space-time geometry and depend only on the quantum field. For example, for a massless scalar field, $`k_1=k_2=k_3=(2880\pi ^2)^1`$. Although all methods of regularization agree on the values of $`k_1`$, $`k_2`$, $`k_3`$ for scalar and neutrino fields, and on $`k_1`$ and $`k_2`$ for the electromagnetic field, there is disagreement on the value of $`k_3`$. Dimensional regularization gives $`k_3=0`$ whilst both point separation and $`\zeta `$-function renormalization give $`k_3=(96\pi ^2)^1`$. This discrepancy is unimportant for us as $`R=0`$ for a Kerr black hole. For a Kerr black hole of mass $`M`$ and angular momentum $`Ma`$, $`C_{\alpha \beta \gamma \delta }C^{\alpha \beta \gamma \delta }`$ $`=`$ $`48\rho ^{12}\{M^2r^815M^2r^4a^2\mathrm{cos}^2\theta `$ (126) $`+15M^2r^2a^4\mathrm{cos}^4\theta M^2a^6\mathrm{cos}^6\theta \}`$ where, as before, $`\rho ^2=r^2+a^2\mathrm{cos}^2\theta `$. The trace anomaly is, of course, finite except at a curvature singularity of the space-time. We may now replace one of the stress tensor components by the trace. Hence we may substitute $$T_\theta {}_{}{}^{\theta }=T_\alpha {}_{}{}^{\alpha }T_t{}_{}{}^{t}T_r{}_{}{}^{r}T_\varphi ^\varphi $$ (127) to yield $`\stackrel{~}{F}(r,\theta )`$ $`=`$ $`_r(\rho ^2T_r{}_{}{}^{r})+\mathrm{\Delta }^1\mathrm{csc}\theta _\theta (\rho ^2\mathrm{sin}\theta T_\theta {}_{}{}^{r})+T_r^r`$ (129) $`\mathrm{\Delta }^1(ra^2\mathrm{sin}^2\theta \mathrm{\Lambda })T_r^r`$ $`\stackrel{~}{G}(r,\theta )`$ $`=`$ $`_r(\rho ^2T_\theta {}_{}{}^{r})\mathrm{csc}\theta _\theta (\rho ^2\mathrm{sin}\theta T_r{}_{}{}^{r})`$ (130) where $`\stackrel{~}{F}(r,\theta )`$ $`=`$ $`F(r,\theta )+rT_\alpha {}_{}{}^{\alpha }rT_t{}_{}{}^{t}rT_\varphi ^\varphi `$ (132) $`=`$ $`rT_\alpha {}_{}{}^{\alpha }+{\displaystyle \frac{(Mr)}{\mathrm{\Delta }^2}}[(r^2+a^2)^2T_{tt}+2a(r^2+a^2)T_{t\varphi }+a^2T_{\varphi \varphi }]`$ (134) $`+{\displaystyle \frac{2r}{\mathrm{\Delta }}}\left[(r^2+a^2)T_{tt}+aT_{t\varphi }\right]`$ $`\stackrel{~}{G}(r,\theta )`$ $`=`$ $`G(r,\theta )a^2\mathrm{sin}\theta \mathrm{cos}\theta (T_\alpha {}_{}{}^{\alpha }T_t{}_{}{}^{t}T_\varphi {}_{}{}^{\varphi })`$ (136) $`\mathrm{csc}\theta _\theta \left(\rho ^2\mathrm{sin}\theta [T_\alpha {}_{}{}^{\alpha }T_t{}_{}{}^{t}T_\varphi {}_{}{}^{\varphi }]\right)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }\mathrm{sin}\theta }}_\theta \left(\mathrm{sin}\theta \left[(r^2+a^2)^2T_{tt}+2a(r^2+a^2)T_{t\varphi }+a^2T_{\varphi \varphi }\right]\right)`$ (139) $`+2a\mathrm{cot}\theta (a\mathrm{sin}^2\theta T_{tt}+T_{t\varphi })+a^2\mathrm{sin}\theta _\theta T_{tt}+2a_\theta T_{t\varphi }`$ $`+\mathrm{csc}^2\theta _\theta T_{\varphi \varphi }\rho ^2_\theta T_\alpha {}_{}{}^{\alpha }+\mathrm{cos}\theta (a^2\mathrm{sin}\theta \rho )^2T_\alpha {}_{}{}^{\alpha }.`$ Equations (130) can be written in the alternative form: $`_r(\mathrm{\Delta }^{\frac{1}{2}}\rho ^2\mathrm{sin}\theta T_r{}_{}{}^{r})+_\theta (\mathrm{\Delta }^{\frac{1}{2}}\rho ^2\mathrm{sin}\theta T_\theta {}_{}{}^{r})`$ $`=`$ $`\mathrm{\Delta }^{\frac{1}{2}}\stackrel{~}{F}(r,\theta )\mathrm{sin}\theta `$ (140) $`_r(\rho ^2\mathrm{sin}\theta T_\theta {}_{}{}^{r})_\theta (\rho ^2\mathrm{sin}\theta T_r{}_{}{}^{r})`$ $`=`$ $`\stackrel{~}{G}(r,\theta )\mathrm{sin}\theta .`$ (141) These equations can now be integrated over $`r`$ to give $`T_r^r`$ $`=`$ $`{\displaystyle \frac{R(\theta )}{\mathrm{\Delta }^{\frac{1}{2}}\rho ^2}}+{\displaystyle \frac{1}{\mathrm{\Delta }^{\frac{1}{2}}\rho ^2\mathrm{sin}\theta }}{\displaystyle _{r_+}^r}(\mathrm{\Delta }^{\frac{1}{2}}\stackrel{~}{F}(r^{},\theta )\mathrm{sin}\theta \mathrm{\Delta }^{\frac{1}{2}}_\theta (\rho ^2\mathrm{sin}\theta T_\theta {}_{}{}^{r}))dr^{}`$ (142) $`T_\theta ^r`$ $`=`$ $`{\displaystyle \frac{S(\theta )}{\rho ^2}}+{\displaystyle \frac{1}{\rho ^2\mathrm{sin}\theta }}{\displaystyle _{r_+}^r}(\stackrel{~}{G}(r^{},\theta )\mathrm{sin}\theta +_\theta (\rho ^2\mathrm{sin}\theta T_r{}_{}{}^{r}))dr^{}`$ (143) where $`R(\theta )`$, $`S(\theta )`$ are arbitrary functions of $`\theta `$ alone. Choice of a particular vacuum state will place restrictions on the four arbitrary functions $`K(\theta )`$, $`L(\theta )`$, $`R(\theta )`$ and $`S(\theta )`$ and also on $`\stackrel{~}{F}(r,\theta )`$ and $`\stackrel{~}{G}(r,\theta )`$, which depend on three unknown stress tensor components, $`T_{tt}`$, $`T_{t\varphi }`$ and $`T_{\varphi \varphi }`$. The solutions (143) are particularly useful for finding the behaviour of the stress tensor close to the event horizon, but we still have the coupling between $`T_r^r`$ and $`T_\theta ^r`$. Uncoupled equations for $`T_r^r`$ and $`T_\theta ^r`$, can be obtained from (141) in the form: $`\mathrm{\Delta }^{\frac{1}{2}}_r\left[\mathrm{\Delta }^{\frac{1}{2}}_r𝒯_1\right]+_\theta ^2𝒯_1`$ $`=`$ $`\mathrm{\Delta }^{\frac{1}{2}}_r\left(\right)_\theta \left(𝒢\right)`$ (144) $`\mathrm{\Delta }^{\frac{1}{2}}_r\left[\mathrm{\Delta }^{\frac{1}{2}}_r\left(𝒯_2\right)\right]+_\theta ^2\left(𝒯_2\right)`$ $`=`$ $`\mathrm{\Delta }^{\frac{1}{2}}_r\left(𝒢\right)+_\theta \left(\right)`$ (145) where $$𝒯_1=T_r{}_{}{}^{r}\mathrm{\Delta }_{}^{\frac{1}{2}}\rho ^2\mathrm{sin}\theta 𝒯_2=T_\theta {}_{}{}^{r}\rho _{}^{2}\mathrm{sin}\theta =\mathrm{\Delta }\stackrel{~}{F}\mathrm{sin}\theta 𝒢=\mathrm{\Delta }^{\frac{1}{2}}\stackrel{~}{G}\mathrm{sin}\theta .$$ (146) We now define a new variable $`x`$ by: $$x=2\mathrm{\Delta }^{\frac{1}{2}}+2r2M,$$ (147) in terms of which the equations (145) now have the usual polar form of the Laplacian: $`x_x\left[x_x𝒯_1\right]+_\theta ^2𝒯_1`$ $`=`$ $`x_x_\theta 𝒢`$ (149) $`x_x\left[x_x𝒯_2\right]+_\theta ^2𝒯_2`$ $`=`$ $`x_x𝒢+_\theta .`$ (150) The domain of these equations is $`x((r^2a^2)^{\frac{1}{2}},\mathrm{})`$, $`\theta (0,\pi )`$, that is the punctured half-plane. By constructing a Green’s function for this domain, a unique solution for $`𝒯_1`$ and $`𝒯_2`$ can be found if they are specified on the boundary, provided we know $``$ and $`𝒢`$ throughout the region. Therefore, we need to know $`T_r^r`$ and $`T_\theta ^r`$ on the event horizon (where $`x=(r^2a^2)^{\frac{1}{2}}`$), and the three components of the stress tensor, $`T_{tt}`$, $`T_{t\varphi }`$ and $`T_{\varphi \varphi }`$ everywhere outside the event horizon. From Eq. (146), it can be seen that $`𝒯_1`$ and $`𝒯_2`$ must vanish on the axis $`\theta =0,\pi `$ provided that $`T_r^r`$ and $`T_\theta ^r`$ are well-defined there. Therefore this reduces the number of boundary functions which are unknown. Although it looks like $`𝒯_1`$ vanishes on the event horizon, the analysis of subsection IV D will show that even for a quantum state which is regular on the event horizon, $`T_r^r`$ diverges as $`\mathrm{\Delta }^1`$ as $`rr_+`$, giving a divergent value for $`𝒯_1`$ on the horizon. This means that the Green’s function method is not directly applicable to Eq. (149). However, the second equation can be solved uniquely using a Green’s function, and the solution then fed into Eq. (143) to give the behaviour of $`T_r^r`$. Note that our calculations in section V confirm that, for the Unruh and (past and future) Boulware vacua, the function $`𝒯_2`$ vanishes sufficiently quickly at infinity that the Green’s function method gives a unique solution. ### C The Killing-Yano Tensor So far in our analysis we have exploited the Killing vector symmetries of the Kerr geometry to assume that the stress tensor is a function only of $`r`$ and $`\theta `$. The Kerr geometry also possesses a Killing-Yano tensor , which is a skew-symmetric tensor $`f_{\mu \nu }`$ satisfying $$^{(\mu }f^{\nu )\lambda }=0.$$ (151) We shall now show that the consequence of the existence of the Killing-Yano tensor is that $`T_{x\theta }=0`$, when $`x=t`$ or $`x=\varphi `$, for the quantum states we are interested in. For any quantum state, the renormalized expectation value of the quantum stress tensor can be calculated using the technique of point splitting: $$T_{\mu \nu }_{\mathrm{ren}}=\underset{xx^{}}{lim}\left[T_{\mu \nu }(x,x^{})T_{\mu \nu }^{\mathrm{div}}(x,x^{})\right]$$ (152) where $`T_{\mu \nu }(x,x^{})`$ is the point-separated stress tensor for our particular quantum state and $`T_{\mu \nu }^{\mathrm{div}}(x,x^{})`$ are the divergent subtraction terms. The unrenormalized stress tensor components for the quantum states in which we are interested are given as mode sums (95108), the mode sum contribution to $`T_{A\theta }`$ for $`A=t`$ or $`A=\varphi `$ being $$T_{A\theta }[u,u^{}]=\mathrm{}e\left[\frac{1}{3}\left(u_{;A}u_{;\theta }^{}+u_{;A}^{}u_{;\theta }\right)\frac{1}{6}\left(u_{;A\theta }u^{}+u_{;A\theta }^{}u\right)\right].$$ (153) The existence of the Killing-Yano tensor has the result that the wave equation for a massless scalar field on the Kerr geometry is separable , with the mode solutions given by (2). In addition, we have $$u_{;A\theta }=u_{,A\theta }\mathrm{\Gamma }_{A\theta }^tu_{,t}\mathrm{\Gamma }_{A\theta }^\varphi u_{,\varphi }.$$ (154) From the mode functions, $`u_{,A}=iku`$ where $`k=\omega `$ if $`A=t`$ and $`k=m`$ if $`A=\varphi `$; also $$u_{,\theta }\left(r^2+a^2\right)^{\frac{1}{2}}e^{i\omega t+im\varphi }R_{\omega lm}(r)S_{\omega lm}^{}(\mathrm{cos}\theta )\mathrm{sin}\theta ;$$ (155) and $`u_{,x\theta }=iku_{,\theta }`$. Since the spheroidal harmonics $`S_{\omega lm}`$ are real, the quantities appearing in (153) are all purely imaginary and hence $`T_{x\theta }[u,u^{}]=0`$. Therefore the point-separated stress tensor components $`T_{A\theta }(x,x^{})`$ vanish for all of the states under consideration. In addition, it is shown in that the subtraction terms $`T_{A\theta }^{\mathrm{div}}(x,x^{})`$ are also zero, so that $`T_{x\theta }_{\mathrm{ren}}`$ vanishes for all the states we are considering here. This property was proved by Frolov and Thorne for $`\|FT`$ but we have shown here that this is a quite general result. ### D Behaviour on the event and Cauchy horizons Next we shall investigate the behaviour of the stress tensor at the future and past event and Cauchy horizons. It is convenient to introduce co-ordinate systems that are regular at the horizons. We first introduce two double-null co-ordinate systems $`u`$, $`v`$, $`\theta `$, $`\varphi _\pm `$ by $$u=tr_{},v=t+r_{},\varphi _\pm =\varphi \frac{a}{r_\pm ^2+a^2}t=\varphi \mathrm{\Omega }_\pm t,$$ (156) where the last equation defines $`\mathrm{\Omega }_+`$ and $`\mathrm{\Omega }_{}`$ which are the angular velocity of the event and Cauchy horizons respectively, and $`r_{}`$ is the ‘tortoise’ co-ordinate given by Eq. (8). The two sets of Kruskal co-ordinates $`U_\pm `$, $`V_\pm `$ are then defined by $$U_\pm =e^{\kappa _\pm u},V_\pm =e^{\kappa _\pm v}.$$ (157) From the definition of $`U_\pm `$, $`V_\pm `$ and $`r_{}`$, $$U_\pm V_\pm =e^{2\kappa _\pm r_{}}=e^{2\kappa _\pm r}(rr_\pm )\|rr_{}\|^{\kappa _\pm /\kappa _{}}.$$ (158) The exterior region corresponds to $`U_+<0`$, $`V_+>0`$ with the past event horizon at $`V_+=0`$ and the future event horizon at $`U_+=0`$. These coordinates may be extended to cover the event horizons in a regular fashion but are singular at the Cauchy horizons. Correspondingly, the coordinates $`U_{}`$ and $`V_{}`$ may be extended to cover the Cauchy horizons ($`U_{}=0`$ and $`V_{}=0`$) in a regular fashion but are singular at the event horizons. The stress tensor components in these Kruskal co-ordinate systems are $`T_{U_\pm U_\pm }`$ $`=`$ $`\kappa _\pm ^2U_\pm ^2[{\displaystyle \frac{1}{4}}T_{tt}+{\displaystyle \frac{1}{2}}\mathrm{\Omega }_\pm T_{t\varphi }+{\displaystyle \frac{1}{4}}\mathrm{\Omega }_\pm ^2T_{\varphi \varphi }{\displaystyle \frac{\mathrm{\Delta }}{2(r^2+a^2)}}T_{rt}`$ (161) $`+{\displaystyle \frac{\mathrm{\Delta }^2}{4(r^2+a^2)^2}}T_{rr}{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\pm }{2(r^2+a^2)}}T_{r\varphi }]`$ $`T_{U_\pm V_\pm }`$ $`=`$ $`\kappa _\pm ^2U_\pm ^1V_\pm ^1\left[{\displaystyle \frac{1}{4}}T_{tt}+{\displaystyle \frac{1}{2}}\mathrm{\Omega }_\pm T_{t\varphi }+{\displaystyle \frac{1}{4}}\mathrm{\Omega }_\pm ^2T_{\varphi \varphi }{\displaystyle \frac{\mathrm{\Delta }^2}{4(r^2+a^2)}}T_{rr}\right]`$ (162) $`T_{V_\pm V_\pm }`$ $`=`$ $`\kappa _\pm ^2V_\pm ^2[{\displaystyle \frac{1}{4}}T_{tt}+{\displaystyle \frac{1}{2}}\mathrm{\Omega }_\pm T_{t\varphi }+{\displaystyle \frac{1}{4}}\mathrm{\Omega }_\pm ^2T_{\varphi \varphi }+{\displaystyle \frac{\mathrm{\Delta }}{2(r^2+a^2)}}T_{rt}`$ (164) $`+{\displaystyle \frac{\mathrm{\Delta }^2}{4(r^2+a^2)^2}}T_{rr}+{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\pm }{2(r^2+a^2)}}T_{r\varphi }]`$ $`T_{U_\pm \theta }`$ $`=`$ $`\kappa _\pm ^1U_\pm ^1{\displaystyle \frac{\mathrm{\Delta }}{2(r^2+a^2)}}T_{r\theta }`$ (165) $`T_{V_\pm \theta }`$ $`=`$ $`\kappa _\pm ^1V_\pm ^1{\displaystyle \frac{\mathrm{\Delta }}{2(r^2+a^2)}}T_{r\theta }`$ (166) $`T_{U_\pm \varphi _\pm }`$ $`=`$ $`\kappa _\pm ^1U_\pm ^1\left[{\displaystyle \frac{1}{2}}T_{t\varphi }+{\displaystyle \frac{1}{2}}\mathrm{\Omega }_\pm T_{\varphi \varphi }{\displaystyle \frac{\mathrm{\Delta }}{2(r^2+a^2)}}T_{r\varphi }\right]`$ (167) $`T_{V_\pm \varphi _\pm }`$ $`=`$ $`\kappa _\pm ^1V_\pm ^1\left[{\displaystyle \frac{1}{2}}T_{t\varphi }+{\displaystyle \frac{1}{2}}\mathrm{\Omega }_\pm T_{\varphi \varphi }+{\displaystyle \frac{\mathrm{\Delta }}{2(r^2+a^2)}}T_{r\varphi }\right]`$ (168) with $`T_{\theta \theta }=T_{\theta \theta }`$, $`T_{\theta \varphi _\pm }=0`$ and $`T_{\varphi _\pm \varphi _\pm }=T_{\varphi \varphi }`$, where we have set $`T_{t\theta }=T_{\theta \varphi }=0`$. It follows immediately that regularity of the stress tensor on any horizon requires that $`T_{\theta \theta }`$, $`T_{\varphi \varphi }`$ and $`T_{r\theta }`$ be finite as the horizon is approached. For a general stress tensor with $`T_{t\theta }=T_{\theta \varphi }=0`$, we have by Eqs. (113) $$T_{tr}=\frac{K(\theta )}{\mathrm{\Delta }},T_{\varphi r}=\frac{L(\theta )}{\mathrm{\Delta }}.$$ (169) In this case, consideration of the $`T_{U_\pm \varphi _\pm }`$ and $`T_{V_\pm \varphi _\pm }`$ components shows that regularity requires $$T_{t\varphi }(r,\theta )=\pm \frac{L(\theta )}{r_\pm ^2+a^2}\mathrm{\Omega }_\pm T_{\varphi \varphi }(r_\pm ,\theta )+O(rr_\pm )\mathrm{as}rr_\pm ,$$ (170) where the positive sign is taken for regularity on the future horizon ($`U_\pm =0`$) and the negative sign on the past horizon ($`V_\pm =0`$). Note that if $`L(\theta )`$ is non-zero, only one of these conditions can be met on either the future or past event horizon. Regularity of the $`T_{U_\pm U_\pm }`$, $`T_{V_\pm V_\pm }`$ and $`T_{U_\pm V_\pm }`$ components implies that $`T_{tt}`$ $`=`$ $`\pm {\displaystyle \frac{K(\theta )\mathrm{\Omega }_\pm L(\theta )}{r^2+a^2}}+\mathrm{\Omega }_\pm ^2T_{\varphi \varphi }+O(rr_\pm )`$ (172) $`T_{rr}`$ $`=`$ $`\pm \left[K(\theta )+\mathrm{\Omega }_\pm L(\theta )\right]{\displaystyle \frac{r^2+a^2}{\mathrm{\Delta }^2}}+O(rr_\pm )^1`$ (173) as $`rr_\pm `$, with the positive sign for regularity on the future horizon, and the negative sign for the past horizon, as before. Finiteness of $`T_{r\theta }`$ as the horizon is approached implies that the function $`S(\theta )`$ in (143) vanishes, whilst the form (173) of $`T_{rr}`$ near the horizon tells us that $`R(\theta )`$ in (143) is also identically zero. It should be stressed that the forms (170-173) are compatible with the solution of the conservation equations (143) with $`R`$ and $`S`$ identically zero. We note that our analysis is in agreement with that of , in that unless both $`K(\theta )`$ and $`L(\theta )`$ vanish identically, the stress tensor must diverge at one of the event horizons, and at least one of the Cauchy horizons. The past and future Boulware vacua are not expected to be regular on either event horizon. For the Unruh vacuum state, it is expected that the divergences occur on the past event horizon and future Cauchy horizon . For $`\|FT`$ simultaneous $`t`$-$`\varphi `$ invariance required that $`K(\theta )`$ and $`L(\theta )`$ vanish consistent with regularity. On the other hand, for $`\|CCH`$ there was no requirement that $`K(\theta )`$ and $`L(\theta )`$ vanish and so one expects that there will be divergences on the past event horizon in line with the Unruh vacuum. We shall return to this issue in section V. At this stage, we need to step back and see how much information about the stress tensor we have managed to obtain from our approach. We began with ten stress tensor components, each a function of the two variables $`r`$ and $`\theta `$. The Killing-Yano symmetry revealed that two of these components $`T_{t\theta }`$ and $`T_{\theta \varphi }`$ vanished identically, whilst another component could be eliminated by using the known trace anomaly, leaving seven unknown functions of $`r`$ and $`\theta `$. Using the conservation equations, we need to know three functions of $`r`$ and $`\theta `$ (corresponding to $`T_{tt}`$, $`T_{t\varphi }`$ and $`T_{\varphi \varphi }`$), and four functions of $`\theta `$ ($`K`$, $`L`$, $`T_{rr}`$ and $`T_{r\theta }`$ on the event horizon). Finally, for a state which is regular on one of the event horizons, this reduces to three functions of $`\theta `$ since the behaviour of $`T_{rr}`$ is given in terms of $`K`$ and $`L`$. In addition, we know the behaviour of the three unknown components, $`T_{tt}`$, $`T_{t\varphi }`$ and $`T_{\varphi \varphi }`$ on the event horizon, in terms of $`K`$, $`L`$ and $`T_{\varphi \varphi }`$. Thus our analysis has significantly reduced the number of degrees of freedom of the stress tensor in Kerr space-time. Of course, this reduction is rather less significant than the corresponding analysis for Schwarzschild black holes , but this was to be expected due to the fact that Kerr has fewer symmetries than Schwarzschild. ## V Asymptotic behaviour of the physical vacua In this section we shall consider the asymptotic behaviour of the physical states of interest near the event horizon and at infinity. This will provide a consistency check on the analysis of the previous section. We shall also use the properties of the Unruh and Boulware vacua (whose asymptotic behaviour are well understood) to reveal information about the states $`\|CCH`$ and $`\|FT`$. It is known that the divergent terms which have to be subtracted from the unrenormalized expectation value of the stress tensor are independent of the quantum state under consideration. Therefore we shall consider the differences in expectation values of the stress tensor in two different states, since these can be calculated without renormalization. Such differences in expectation values will be traceless tensors since the trace anomaly is the same for all quantum states. We shall begin by concentrating on the Unruh vacuum, since its stress tensor has been calculated in the asymptotic regimes by Punsley using an equivalence principle approach. This will provide a useful check of our calculations. Firstly, we consider the behaviour at infinity, and calculate $`U^{}\|\widehat{T}_{\mu \nu }\|U^{}_{ren}B^{}\|\widehat{T}_{\mu \nu }\|B^{}_{ren}`$ $`=`$ $`U^{}\|\widehat{T}_{\mu \nu }\|U^{}B^{}\|\widehat{T}_{\mu \nu }\|B^{}`$ (174) $`=`$ $`{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{2d\stackrel{~}{\omega }}{e^{2\pi \stackrel{~}{\omega }/\kappa }1}}T_{\mu \nu }[u_{\omega lm}^{up},u_{\omega lm}^{up}].`$ (175) Using the asymptotic form of the mode functions (12), we have, as $`r\mathrm{}`$, $`U^{}\|\widehat{T}_\mu ^\nu \|U^{}_{ren}B^{}\|\widehat{T}_\mu ^\nu \|B^{}_{ren}`$ (176) $``$ $`{\displaystyle \frac{1}{4\pi ^2r^2}}{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\omega \mathrm{d}\stackrel{~}{\omega }}{\stackrel{~}{\omega }(e^{2\pi \stackrel{~}{\omega }/\kappa }1)}}\|B_{\omega lm}^{}\|^2\|S_{\omega lm}(\mathrm{cos}\theta )\|^2\left(\begin{array}{cccc}\omega & \omega & 0& m\\ \omega & \omega & 0& m\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right).`$ (181) In order to obtain the behaviour of the Unruh vacuum at future null infinity, we need to consider the ‘past’ Boulware vacuum at infinity. The ‘past’ Boulware vacuum contains at future null infinity an outward flux of particles due to the Unruh-Starobinskii effect , so that, as we approach $`^+`$, $`B^{}\|\widehat{T}_\mu ^\nu \|B^{}_{ren}`$ $``$ $`B^{}\|\widehat{T}_\mu ^\nu \|B^{}B^+\|\widehat{T}_\mu ^\nu \|B^+`$ (182) $``$ $`{\displaystyle \frac{1}{4\pi ^2r^2}}{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\omega _{min}}}{\displaystyle \frac{\omega d\omega }{\stackrel{~}{\omega }(e^{2\pi \stackrel{~}{\omega }/\kappa }1)}}\|B_{\omega lm}^{}\|^2\|S_{\omega lm}(\mathrm{cos}\theta )\|^2\left(\begin{array}{cccc}\omega & \omega & 0& m\\ \omega & \omega & 0& m\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right).`$ (187) Adding these two tensors gives the asymptotic behaviour of the Unruh vacuum at future null infinity as: $`U\|\widehat{T}_\mu ^\nu \|U_{ren}`$ $``$ $`{\displaystyle \frac{1}{4\pi ^2r^2}}{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\omega d\omega }{\stackrel{~}{\omega }(e^{2\pi \stackrel{~}{\omega }/\kappa }1)}}\|B_{\omega lm}^{}\|^2\|S_{\omega lm}(\mathrm{cos}\theta )\|^2\left(\begin{array}{cccc}\omega & \omega & 0& m\\ \omega & \omega & 0& m\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right).`$ (192) This is in agreement with the form obtained in , and represents the expected thermal flux at infinity. It should be noted that, despite initial appearances, the integrands are regular when $`\stackrel{~}{\omega }=0`$ due to the Wronskian relations (17) which ensure that $`\|B_{\omega lm}^{}\|^2=O(\stackrel{~}{\omega }^2)`$ as $`\stackrel{~}{\omega }0`$. From Eq. (192) we can read off the forms of the functions $`K`$ and $`L`$ (169) for the Unruh vacuum: $`K_U^{}(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\omega ^2d\omega }{\stackrel{~}{\omega }(e^{2\pi \stackrel{~}{\omega }/\kappa }1)}}\|B_{\omega lm}^{}\|^2\|S_{\omega lm}(\mathrm{cos}\theta )\|^2`$ (193) $`L_U^{}(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{m\omega d\omega }{\stackrel{~}{\omega }(e^{2\pi \stackrel{~}{\omega }/\kappa }1)}}\|B_{\omega lm}^{}\|^2\|S_{\omega lm}(\mathrm{cos}\theta )\|^2.`$ (194) We now turn to the behaviour of the Unruh vacuum at the event horizon. In Schwarzschild, the Hartle-Hawking state is regular on both event horizons, and so the behaviour of the Unruh vacuum as $`rr_+`$ is found from: $`U^{}\|\widehat{T}_{\mu \nu }\|U^{}_{ren}`$ $``$ $`U^{}\|\widehat{T}_{\mu \nu }\|U^{}_{ren}H\|\widehat{T}_{\mu \nu }\|H_{ren}`$ (195) $`=`$ $`U^{}\|\widehat{T}_{\mu \nu }\|U^{}H\|\widehat{T}_{\mu \nu }\|H.`$ (196) In the absence of a Hartle-Hawking state for Kerr, we shall instead consider the differences of the stress tensors in the Unruh vacuum and the states $`\|FT`$ and $`\|CCH`$. These are given by: $`U^{}\|\widehat{T}_{\mu \nu }\|U^{}FT\|\widehat{T}_{\mu \nu }\|FT`$ $`=`$ $`{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{2d\omega }{e^{2\pi \stackrel{~}{\omega }/\kappa }1}}T_{\mu \nu }[u_{\omega lm}^{in},u_{\omega lm}^{in}],`$ (198) $`U^{}\|\widehat{T}_{\mu \nu }\|U^{}CCH\|\widehat{T}_{\mu \nu }\|CCH`$ $`=`$ $`{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{2d\omega }{e^{2\pi \omega /\kappa }1}}T_{\mu \nu }[u_{\omega lm}^{in},u_{\omega lm}^{in}].`$ (199) As $`rr_+`$, one finds $`U^{}\|\widehat{T}_\mu ^\nu \|U^{}FT\|\widehat{T}_\mu ^\nu \|FT`$ (200) $`{\displaystyle \frac{1}{4\pi ^2\rho ^2}}{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega }{\omega (e^{2\pi \stackrel{~}{\omega }/\kappa }1)}}\|B_{\omega lm}^+\|^2\|S_{\omega lm}(\mathrm{cos}\theta )\|^2`$ (201) $`\times \left(\begin{array}{cccc}\mathrm{\Delta }^1(r_+^2+a^2)\omega \stackrel{~}{\omega }& \omega \stackrel{~}{\omega }& 0& \mathrm{\Delta }^1a\omega \stackrel{~}{\omega }\\ \mathrm{\Delta }^2(r_+^2+a^2)^2\stackrel{~}{\omega }^2& \mathrm{\Delta }^1(r_+^2+a^2)^2\stackrel{~}{\omega }^2& O(1)& \mathrm{\Delta }^2a(r_+^2+a^2)\stackrel{~}{\omega }^2\\ 0& O(\mathrm{\Delta })& O(1)& 0\\ \mathrm{\Delta }^1(r_+^2+a^2)m\stackrel{~}{\omega }& m\stackrel{~}{\omega }& 0& \mathrm{\Delta }^1am\stackrel{~}{\omega }\end{array}\right).`$ (206) The expression for $`U^{}\|\widehat{T}_\mu ^\nu \|U^{}CCH\|\widehat{T}_\mu ^\nu \|CCH`$ is identical to Eq. (206), with the denominator $`e^{2\pi \stackrel{~}{\omega }/\kappa }1`$ replaced by $`e^{2\pi \omega /\kappa }1`$. In both cases the integrand is regular for all values of $`\omega `$, by virtue of the Wronskian relations (17). The difference in expectation values of the stress tensor in the Unruh and Frolov-Thorne states (206) agrees with the stress tensor for the Unruh vacuum found in , whereas when we have the state $`\|CCH`$ instead of $`\|FT`$ the thermal terms in the denominator do not agree. Furthermore, the tensor (206) is regular on the future event horizon but not on the past event horizon, the same behaviour that we would expect for the Unruh vacuum. Therefore we can compare the tensor (206) with the behaviour near the event horizon derived in section IV D. There is exact agreement, using the functions $`K_U^{}(\theta )`$ and $`L_U^{}(\theta )`$ found from the expectation value of the stress tensor at infinity in the Unruh vacuum (194), and the Wronskian relations. From the regularity of the tensor (206) on the future event horizon, we can conclude that the expectation value of the stress tensor in the state $`\|FT`$ is regular on at least one event horizon (and, since it is invariant under simultaneous $`t`$, $`\varphi `$ reversal, it will be regular on both event horizons). Thus, it may appear that the state $`\|FT`$ in fact has the properties that we require of the Hartle-Hawking state. However, whilst the expectation value of the stress tensor in the state $`\|FT`$ is regular on the event horizon, the expectation value of $`\widehat{\mathrm{\Phi }}^2`$ is not. We calculate, as $`rr_+`$, $`U^{}\|\widehat{\mathrm{\Phi }}^2\|U^{}FT\|\widehat{\mathrm{\Phi }}^2\|FT`$ (207) $``$ $`{\displaystyle \frac{1}{4\pi ^2(r_+^2+a^2)}}{\displaystyle \underset{l,m}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{2d\omega }{\omega (e^{2\pi \stackrel{~}{\omega }/\kappa }1)}}\|B_{\omega lm}^+\|^2\|S_{\omega lm}(\mathrm{cos}\theta )\|^2.`$ (208) The integrand in the above expression is regular at $`\omega =0`$ because of the Wronskian relations (17), but has a pole at $`\stackrel{~}{\omega }=0`$, giving a divergent integral. If we attempt to calculate the difference in expectation values (198) anywhere outside the event horizon, then the integral over $`\omega `$ also has a pole at $`\stackrel{~}{\omega }=0`$, leading to a divergent result. Therefore it seems that the regularity of the difference in expectation values of the stress tensor (206) at the event horizon does not reflect the true nature of the state $`\|FT`$, and that this state in fact fails to be regular almost everywhere, both on or outside the event horizon, although it formally has attractive symmetry properties. There is one exception to the regularity of the state $`\|FT`$ which is that on the axis the terms with $`m0`$ (and, in particular, all superradiant modes) do not contribute. Thus, if one point is on the axis the $`\|FT`$ and $`\|CCH`$ two-point functions agree: $`G^{FT/CCH}(t,r,\theta ,\varphi ;t^{},r^{},0,\varphi ^{})`$ $`=`$ $`{\displaystyle \underset{l}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega \mathrm{coth}(\pi \omega /\kappa )}{\omega \sqrt{(r^2+a^2)(r_{}^{}{}_{}{}^{2}+a^2)}}}`$ (210) $`\times \left[R_{\omega l0}^+(r)R_{\omega l0}^+(r^{})+R_{\omega l0}^{}(r)R_{\omega l0}^{}(r^{})\right]S_{\omega l0}(\mathrm{cos}\theta )S_{\omega l0}(1).`$ In the asymptotic regions, the integrals are dominated by the contribution from near $`\omega =0`$. In this limit the spheroidal functions reduce to Legendre polynomials $$S_{0lm}=\frac{1}{\sqrt{4\pi }}P_l(\mathrm{cos}\theta ),\lambda (0)=l(l+1).$$ (211) In addition, $`T_{lm}(r)=R_{0lm}(r)/\sqrt{r^2+a^2}`$ satisfies the equation $$\frac{\mathrm{d}}{\mathrm{d}\eta }(\eta ^21)\frac{\mathrm{d}T_{lm}}{\mathrm{d}\eta }\left[l(l+1)+\frac{m^2a^2}{(M^2a^2)(\eta ^21)}\right]T_{lm}=0$$ (212) where $$\eta =\frac{2r(r_++r_{})}{(r_+r_{})}=\frac{rM}{\sqrt{M^2a^2}},$$ (213) with solutions $`P_l^{ma/\sqrt{M^2a^2}}(\eta )`$ and $`Q_l^{ma/\sqrt{M^2a^2}}(\eta )`$. In particular, a steepest descent analysis of Eq. (210) as $`r^{}r_+`$ yields $`G^{FT/CCH}(t,r,\theta ,\varphi ;t^{},r_+,0,\varphi ^{})`$ $`=`$ $`{\displaystyle \frac{\kappa _+}{16\pi ^2\sqrt{M^2a^2}}}{\displaystyle \underset{l}{}}(2l+1)Q_l\left({\displaystyle \frac{rM}{\sqrt{M^2a^2}}}\right)P_l(\mathrm{cos}\theta )`$ (214) $`=`$ $`{\displaystyle \frac{\kappa _+}{8\pi ^2}}{\displaystyle \frac{1}{rM\sqrt{M^2a^2}\mathrm{cos}\theta }},`$ (215) where the second line follows from Heine’s formula. This result was first given by Frolov and enabled him to calculate the renormalized value of the expectation value of $`\widehat{\mathrm{\Phi }}^2`$ on the pole of the event horizon. Later with Zel’nikov he extended this calculation to calculate the renormalized value of the expectation value of $`\widehat{T}_{\mu \nu }`$ on the pole of the event horizon. Our point is that, unfortunately, these calculations were only possible because the troublesome superradiant modes do not contribute on the axis and have actually led to a false confidence concerning the Hartle-Hawking vacuum. Finally we return the properties of the state $`\|CCH`$. This has a different thermal factor from $`\|FT`$ (103) which means that the difference in expectation values of the stress tensor in $`\|U^{}`$ and $`\|CCH`$ at the event horizon is rather different from simply the stress tensor in the state $`\|U^{}`$. The difference in thermal factors also means that the state $`\|CCH`$ is not invariant under simultaneous $`t`$-$`\varphi `$ reversal. However, the quantity $`U^{}\|\widehat{T}_\mu ^\nu \|U^{}CCH\|\widehat{T}_\mu ^\nu \|CCH`$ is regular on the future event horizon (but not on the past), so, using the expected regularity of the Unruh vacuum, we can conclude that $`CCH\|\widehat{T}_\mu ^\nu \|CCH`$ is also regular on the future event horizon (but not on the past). If we consider the difference in expectation values of $`\widehat{\mathrm{\Phi }}^2`$ at the event horizon, the answer is the same as (208), but with $`e^{2\pi \stackrel{~}{\omega }/\kappa }`$ replaced by $`e^{2\pi \omega /\kappa }`$. Using the Wronskian relations (17), this gives a finite answer, further strengthening our argument that $`\|CCH`$ is a regular state on the future event horizon. ## VI Conclusions In this paper we have considered the renormalized stress energy tensor on Kerr space-time, and used the anticipated physical properties of this tensor (symmetry, conservation equations, and regularity conditions) in order to derive as much information as possible. As expected, the analysis is considerably more complex than the corresponding problem in Schwarzschild , and the solution gives us less information, although we are able to reduce the number of unknowns to three functions of $`r`$ and $`\theta `$ and three functions of $`\theta `$. Our results are in agreement with the known form of the Unruh vacuum at the event horizon and at infinity. We also considered two candidates for the state analogous to the Hartle-Hawking state in Schwarzschild. From the Kay-Wald theorem , we know that there is no state in Kerr which is regular at the event horizon and everywhere outside, invariant under simultaneous $`t`$, $`\varphi `$ reversal and thermal in nature. Of our two candidate states, one is invariant under $`t`$, $`\varphi `$ reversal, but fails to be regular on the event horizon, whilst the other is regular on the event horizon but not invariant under simultaneous $`t`$, $`\varphi `$ reversal. We should add that are conclusions are based on a mode by mode analysis and it is possible, though in our opinion unlikely, that subtle cancellations could rescue the Frolov-Thorne state. A detailed numerical investigation would be necessary to elucidate further details of the properties of these states outside the event horizon. This paper has laid the foundation for such an investigation which we will present in following papers in this series. It is possible to draw some conclusions on the basis of our analysis without resorting to a numerical investigation. For example, one can show that any state which is isotropic in a tetrad which co-rotates with the event horizon must become divergent on the velocity of light surface . This implies that even if we could construct a state which is regular on the event horizon and has the desired thermal properties, then that state may well turn out not to be regular on the velocity of light surface, in agreement with the Kay-Wald theorem that the state must fail to be regular somewhere. This paper has shown that whilst quantum field theory in Kerr space-time is more complex than in Schwarzschild, application of the same physical principles which have proved to be so valuable in Schwarzschild also makes the picture much clearer and more simple in Kerr. ###### Acknowledgements. E.W. thanks Oriel College, Oxford, for a fellowship supporting this research, and the Department of Physics, University of Newcastle, Newcastle-upon-Tyne, for hospitality during the completion of this work.
warning/0004/math-ph0004032.html	ar5iv	text	# Untitled Document $`Z_3`$-GRADED EXTERIOR DIFFERENTIAL CALCULUS AND GAUGE THEORIES OF HIGHER ORDER by Richard Kerner Laboratoire de Gravitation et Cosmologie Relativistes Université Pierre-et-Marie-Curie, CNRS - D0 769 Tour 22, 4-ème étage, Boîte 142 4, Place Jussieu, 75005 Paris FRANCE Summary. We present a possible generalization of the exterior differential calculus, based on the operator $`d`$ such that $`d^3=0`$, but $`d^20`$. The entities $`dx^i`$ and $`d^2x^k`$ generate an associative algebra; we shall suppose that the products $`dx^idx^k`$ are independent of $`dx^kdx^i`$, while the ternary products will satisfy the relation: $`dx^idx^kdx^m=jdx^kdx^mdx^i=j^2dx^mdx^idx^k`$ , complemented by the relation $`dx^id^2x^k=jd^2x^kdx^i`$, with $`j:=e^{\frac{2\pi i}{3}}`$. We shall attribute grade 1 to the differentials $`dx^i`$ and grade 2 to the ”second differentials” $`d^2x^k`$ ; under the associative multiplication law the grades add up modulo 3. We show how the notion of covariant derivation can be generalized with a 1-form $`A`$ so that $`D\mathrm{\Phi }:=d\mathrm{\Phi }+A\mathrm{\Phi }`$ , and we give the expression in local coordinates of the curvature 3-form defined as $`\mathrm{\Omega }:=d^2A+d(A^2)+AdA+A^3`$. Finally, the introduction of notions of a scalar product and integration of the $`Z_3`$-graded exterior forms enables us to define variational principle and to derive the differential equations satisfied by the 3-form $`\mathrm{\Omega }`$. The lagrangian obtained in this way contains the invariants of the ordinary gauge field tensor $`F_{ik}`$ and its covariant derivatives $`D_iF_{km}`$. 1. INTRODUCTION The models of fundamental interactions and field theories based on the non- commutative geometry have been the object of intense studies in past few years (cf. refs. to ). Most of the effort has been concentrated on reproducing the Weinberg-Salam unified theory of electroweak interactions by means of generalized gauge theories developed in non-commutative geometries instead of fibre bundles, and with finite algebra of matrices replacing the infinite algebra of functions on the manifold. In most developed variants of this approach, the $`Z_2`$-graded matrix algebras have been used, whose $`Z_2`$-graded internal differential (defined as a graded commutator with a matrix whose square was equal to 1) was combined with the usual exterior differential forming a $`Z_2`$-graded Grassmann algebra. In the tensor products of these two algebras, i.e. in the algebra of matrix-valued exterior forms, the $`Z_2`$-grades of the matrices were added (modulo $`2`$) to the $`Z_2`$-grades of the exterior forms, defining the $`Z_2`$\- grade of the whole object (cf. refs. and ). Here we would like to generalize this scheme to the case of $`Z_3`$-grading. The resulting matrix algebra can be represented in the simplest case as the algebra of $`3\times 3`$ matrices with entries from an associative algebra over complex numbers, on which a $`Z_3`$-graded commutator replaces the usual commutation and anti-commutational rules for the $`Z_2`$-graded $`2\times 2`$-matrices. This leads naturally to the differential $`d`$ whose square $`d^2`$ is different from $`0`$, but whose cube does vanish identically, $`d^3=0`$. We show how such matrix algebra appears naturally as the algebra of linear transformations of a $`Z_3`$-graded generalization of the Grassmann algebra. We also introduce a $`Z_3`$-graded generalization of exterior algebra of differential forms and show how the notions of connection and curvature can be generalized, too, and how the lagrangians containing the invariants of the curvature and its derivatives can be constructed. Finally, we briefly discuss the resulting variational principle for the gauge fields, leading to invariant equations of fourth order. Similar theories with even higher order derivatives of the curvature and lagrangians depending on their invariants can be obtained with a similar scheme based on $`Z_N`$-graded differentials satisfying $`d^N=0`$. 2. $`Z_3`$-GRADED ANALOG OF GRASSMANN ALGEBRA. The cyclic group $`Z_3`$ can be represented in the complex plane by means of the cubic roots of 1 : let $`j:=e^{\frac{2\pi i}{3}}`$; then one has $`j^3=1`$ and $`j+j^2+1=0`$; obviously, $`j^n=j^{(n+3)}`$. By analogy with the $`Z_2`$-graded Grassmann algebras spanned by the set of anti-commuting generators, we may introduce an associative algebra spanned by $`N`$ generators $`\theta ^A`$, $`A,B=1,2\mathrm{}N`$, whose binary products $`\theta ^A\theta ^B`$ will be considered as $`N^2`$ independent quantities, whereas we shall impose a ternary analog of the anti-commutation relations: $$\theta ^A\theta ^B\theta ^C=j\theta ^B\theta ^C\theta ^A=j^2\theta ^C\theta ^A\theta ^B$$ (1) A more precise formulation is to say that the algebra in question is the universal algebra defined by the above relations. Corollary : The cube of any generator must vanish (because in this case the relation (1) amounts to $`(\theta ^A)^3=j(\theta ^A)^3=0`$ ; all the monomials of order $`4`$ or higher are identically null (the proof that follows makes use of the associativity of the postulated product and of the relation $`1`$): (the low braces are there just to indicate to which triple of $`\theta `$’s the circular permutation is being applied) $$\stackrel{C}{\underset{}{\theta ^A\theta ^B\theta }}\theta ^D=j\theta ^B\stackrel{D}{\underset{}{\theta ^C\theta ^A\theta }}=j^2\stackrel{D}{\underset{}{\theta ^B\theta ^A\theta }}\theta ^C=\theta ^A\stackrel{C}{\underset{}{\theta ^D\theta ^B\theta }}=j\theta ^A\theta ^B\theta ^C\theta ^D;$$ (2) therefore, as $`1j0`$, one has $`\theta ^A\theta ^B\theta ^C\theta ^D=0`$. The dimension of this $`Z_3`$-graded generalization of Grassmann algebra is equal to $`N+N^2+(N^3N)/3`$; we may also add a ”neutral” element denoted by 1 and commuting with all other generators. One can note a dissymetry between the components of this algebra with the grades 1 et 2 : as a matter of fact, there are $`N`$ elements of grade 1, (the $`\theta `$’s ) and $`N^2`$ elements of grade 2 ($`\theta \theta `$). A natural way to re-establish the symmetry is to introduce the set of $`N`$conjugate” generators , $`\overline{\theta }^A`$, of grade 2, that would satisfy conjugate ternary relations (in which $`j`$ is replaced by $`j^2`$): $$\overline{\theta }^A\overline{\theta }^B\overline{\theta }^C=j^2\overline{\theta }^B\overline{\theta }^C\overline{\theta }^A$$ (3) The ternary relation between the $`\theta ^A`$’s can be interpreted as follows: $`\theta ^A\stackrel{C}{\underset{}{\theta ^B\theta }}=j\stackrel{C}{\underset{}{\theta ^B\theta }}\theta ^A`$, which suggests the following relations between the generators $`\theta ^A`$ and $`\overline{\theta }^B`$ : $$\theta ^A\overline{\theta }^B=j\overline{\theta }^B\theta ^A,\overline{\theta }^B\theta ^A=j^2\theta ^A\overline{\theta }^B.$$ (4) The $`Z_3`$-graded algebra so defined can be naturally divided in three parts, of grade 0, 1 et 2 respectively, with the dimensions of the sub-spaces of grades 1 and 2 being equal: one can write symbolically $`A=A_0+A_1+A_2`$, where $`A_0`$ contains: 1, $`\theta ^A\overline{\theta }^B,\theta ^A\theta ^B\theta ^C,\overline{\theta }^A\overline{\theta }^B\overline{\theta }^C,\theta ^A\theta ^B\overline{\theta }^C\overline{\theta }^D\mathrm{and}\theta ^A\theta ^B\theta ^C\overline{\theta }^D\overline{\theta }^E\overline{\theta }^F`$ ; $`A_1`$ contains: $`\theta ^A,\overline{\theta }^B\overline{\theta }^C,\theta ^A\theta ^B\overline{\theta }^C,\theta ^A\overline{\theta }^A\overline{\theta }^B\overline{\theta }^C`$, and $`A_2`$ contains: $`\overline{\theta }^A,\theta ^A\theta ^B,\theta ^A\overline{\theta }^B\overline{\theta }^C,\theta ^A\theta ^B\theta ^C\overline{\theta }^D`$ . In the case of usual $`Z_2`$-graded Grassmann algebras the anti-commutation between the generators of the algebra and the assumed associativity imply automatically the fact that all grade $`0`$ elements commute with the rest of the algebra, while any two elements of grade $`1`$ anti-commute. In the case of the $`Z_3`$-graded generalization such an extension of ternary and binary relations does not follow automatically, and must be imposed explicitly. If we decide to extend these relations to all elements of the algebra having a well-defined grade (i.e. the monomials in $`\theta `$’s and $`\overline{\theta }`$’s , then many additional expressions must vanish, e.g.: $`\theta ^A\stackrel{C}{\underset{}{\theta ^B\overline{\theta }}}=\stackrel{C}{\underset{}{\theta ^B\overline{\theta }}}\theta ^A=\theta ^B\stackrel{A}{\underset{}{\overline{\theta }^C\theta }}=\overline{\theta }^C\theta ^A\theta ^B=0`$ ; because on the one side, $`\theta ^A\overline{\theta }^C`$ is of grade 0 and commutes with all other elements; at the same time, commuting $`\overline{\theta }^C`$ with $`\theta ^A\theta ^B`$ one gets twice the factor $`j^2`$, which leads to the overall factor $`j\overline{\theta }^C\theta ^A\theta ^B`$; this produces a contradiction which can be solved only by supposing that $`\theta ^A\theta ^B\overline{\theta }^C=0`$. The resulting $`Z_3`$-graded algebra contains only the following products of generators: $$A_1=\theta ,\{\overline{\theta }\overline{\theta }\};A_2=\overline{\theta },\{\theta \theta \};A_0=\{\theta \overline{\theta }\},\{\theta \theta \theta \},\{\overline{\theta }\overline{\theta }\overline{\theta }\}$$ (5) Let us note that the set of grade $`0`$ (which obviously forms a sub-algebra of the $`Z_3`$-graded Grassmann algebra) contains the products which could symbolize the only observable combinations of quark fields in quantum chromodynamics based on the $`SU(3)`$-symmetry. If we align the basis of our algebra, with all the elements of grade $`0`$ first, next all the elements of grade $`1`$ and finally the elements of grade $`2`$ in a one-column vector, a general linear transformation that would leave these entries in the same order can be symbolized by a matrix whose entries have a definite $`Z_3`$-grade placed as follows: $$\left(\begin{array}{ccc}0& 2& 1\\ 1& 0& 2\\ 2& 1& 0\end{array}\right)\left(\begin{array}{c}0\\ 1\\ 2\end{array}\right)=\left(\begin{array}{c}0\\ 1\\ 2\end{array}\right)$$ (6) Note that the position of the three grades did not change in the resulting column; we shall call such an operator a grade 0 matrix. We can introduce two other kinds of matrices that raise all the grades by $`1`$ (resp. by $`2`$), and call them respectively grade 1 and grade 2 matrices, as follows: $$\left(\begin{array}{ccc}1& 0& 2\\ 2& 1& 0\\ 0& 2& 1\end{array}\right)\left(\begin{array}{c}0\\ 1\\ 2\end{array}\right)=\left(\begin{array}{c}1\\ 2\\ 0\end{array}\right),\mathrm{and}\left(\begin{array}{ccc}2& 1& 0\\ 0& 2& 1\\ 1& 0& 2\end{array}\right)\left(\begin{array}{c}0\\ 1\\ 2\end{array}\right)=\left(\begin{array}{c}2\\ 0\\ 1\end{array}\right)$$ (7) (the numbers $`0,1,2`$ symbolize the grades of the respective entries in the matrices). If we restrict the character of the matrices, admitting only complex-valued matrix elements, then the grades $`0,1`$ and $`2`$ will reduce themselves to the following three types of $`3\times 3`$-block matrices: $$\left(\begin{array}{ccc}a& 0& 0\\ 0& b& 0\\ 0& 0& c\end{array}\right),\left(\begin{array}{ccc}0& \alpha & 0\\ 0& 0& \beta \\ \gamma & 0& 0\end{array}\right),\left(\begin{array}{ccc}0& 0& \gamma \\ \alpha & 0& 0\\ 0& \beta & 0\end{array}\right)$$ (8) representing arbitrary matrices with respective grade $`0,1`$ and $`2`$. It is easy to check that these grades add up modulo $`3`$ under the associative matrix multiplication law. Let $`B,C`$ denote two matrices whose grades are $`grad(A)=a`$ and $`grad(B)=b`$ , respectively. We can define the $`Z_3`$-graded commutator $`[B,C]`$ as follows: $$[B,C]_{Z_3}:=BCj^{bc}CB,$$ (9) (Note that this $`Z_3`$-graded commutator does not satisfy the Jacobi identity). Let $`\eta `$ be a matrix of grade $`0`$; we can choose for the sake of simplicity $$\eta =\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 1\\ 1& 0& 0\end{array}\right)$$ (10) With the help of the matrix $`\eta `$ we can define a formal ”differential” on the $`Z_3`$-graded algebra of $`3`$ matrices as follows: $$dB:=[\eta ,B]_{Z_3}=\eta Bj^bB\eta $$ (11) It is easy to show that $`d(BC)=(dB)C+j^bB(dC)`$ and that $`d^3=0`$. The first identity is trivial, whereas the last one follows from the fact that $`\eta ^3=Id`$ does commute with all the elements of the algebra. A formal algebraic analogue of connection and curvature forms have been discussed elsewhere (cf. ref. and ). Here we shall show how such a $`Z_3`$-graded exterior differential may be realized on a differential manifold. Some interesting ideas concerning similar techniques may be found in ref.. 3. $`Z_3`$-GRADED EXTERIOR DIFFERENTIAL. Let $`M_n`$ be a differentiable manifold of dimension $`N`$, with local coordinates $`x^k`$. The variables $`x^k`$ commute and with respect to the $`Z_3`$-grading are of grade $`0`$. Our aim now is to define such a linear operation acting on functions of the coordinates $`x^k`$, and by extension, on a larger algebra of forms that still remains to be defined, that would reproduce the properties of the ”algebraic differential” introduced above. To do this, we have to define a linear operation $`d`$ which acts on the functions of $`x^k`$; for the definition, it is enough to define its action on the coordinates $`x^k`$ and on their products. We postulate that the linear operator $`d`$ applied to $`x^k`$ produces a 1-form whose $`Z_3`$-grade is $`1`$ by definition; when applied two times by iteration, it will produce a new entity, which we shall call a 1-form of grade $`2`$, denoted by $`d^2x^k`$. Finally, we require that $`d^3=0`$. Let F(M) denote the algebra of functions $`C^{\mathrm{}}(M)`$, over which the $`Z_3`$-graded algebra generated by the forms $`dx^i`$ and $`d^2x^k`$ behaves as a left module. In other words, we shall be able to multiply the forms $`dx^i`$ , $`d^2x^k`$, $`dx^idx^k`$ , etc. by the functions on the left only; right multiplication will just not be considered here. That is why we will write by definition, e.g. $$d(x^ix^k):=x^idx^k+x^kdx^i$$ (12) We shall also assume the following Leibniz rule for the operator $`d`$ with respect to the multiplication of the $`Z_3`$-graded forms: when $`d`$ crosses a form of grade $`p`$, and of arbitrary rank, the factor $`j^p`$ appears as follows: $$d(\omega \varphi )=(d\omega )\varphi +j^p\omega (d\varphi )$$ (13) Let us note that in contrast with the $`Z_2`$-graded case, the forms are treated as one whole, even when multiplied from the left by an arbitrary function; that means that we can not identify e.g. $`(\omega _idx^i)(\varphi _kdx^k)`$ with $`(\omega _i\varphi _k)dx^idx^k`$ This is equivalent with saying that the products of functions by the forms are to be understood in the sense of tensor products, which is associative, but non-commutative. Nevertheless, we shall show later that such an identification can be done for the forms of maximal degree (i.e. $`3`$), which contain the products of the type $`dx^idx^kdx^m`$ or $`dx^id^2x^m`$, whose exterior differentials vanish irrespective of the order of the multiplication; as a matter of fact, it will be easy to show that with the rules of differentiation that we expose below, $$d((\alpha _idx^i)(\beta _kdx^k)(\gamma _mdx^m))=d((\alpha _i\beta _k\gamma _m)dx^idx^kdx^m))=0.$$ (14) With the so established $`Z_3`$-graded Leibniz rule, the postulate $`d^3=0`$ suggests in an almost unique way the ternary and binary commutation rules for the differentials $`dx^i`$ and $`d^2x^k`$. To begin with, consider the differentials of a function of the coordinates $`x^k`$, with the ”first differential” $`df`$ coinciding with the usual one: $$df:=(_if)dx^i;d^2f:=(_k_if)dx^kdx^i+(_if)d^2x^i;$$ (15) $$d^3f=(_m_k_if)dx^mdx^kdx^i+(_k_if)d^2x^kdx^i+j(_k_if)dx^id^2x^k+(_k_if)dx^kd^2x^i;$$ (16) (we remind that the last part of the differential, $`(_if)d^3x^i`$, vanishes by virtue of the postulate $`d^3x^i=0`$). Supposing that the partial derivatives commute, exchanging the summation indices $`i`$ et $`k`$ in the last expression and replacing $`1+j`$ by $`j^2`$, we arrive at the following two conditions that lead to the vanishing of $`d^3f`$ : $$dx^mdx^kdx^i+dx^kdx^idx^m+dx^idx^mdx^k=0;d^2x^kdx^ij^2dx^id^2x^k=0.$$ (17) which lead in turn to the following choice of relations: $$dx^idx^kdx^m=jdx^kdx^mdx^i,\mathrm{and}dx^id^2x^k=jd^2x^kdx^i.$$ (18) By extending these rules to all the expressions with a well-defined grade, and applying the associativity of the $`Z_3`$-exterior product, we see that all the expressions of the type $`dx^idx^kdx^mdx^n`$ and $`dx^idx^kd^2x^m`$ must vanish, and along with them, also the monomials of higher order that would contain them as factors. Still, this is not sufficient in order to satisfy the rule $`d^3=0`$ on all the forms spanned by the generators $`dx^1`$ and $`d^2x^k`$. It can be proved without much pain that the expressions containing $`d^2x^id^2x^k`$ must vanish, too. For example, if we take the particular 1-form $`x^idx^k`$ and and apply to it the operator $`d`$, we get $$d(x^idx^k)=dx^idx^k+x^id^2x^k;$$ (19) $$d^2(x^idx^k)=d^2x^idx^k+(1+j)dx^id^2x^k=d^2x^idx^kd^2x^kdx^i;$$ (20) which leads to $`d^3(x^idx^k)=d^2x^id^2x^kd^2x^kd^2x^i`$; then, if we want to keep both the associativity of the ”exterior product” and the ternary rule for the entities of grade $`2`$, i.e. $`d^2x^id^2x^kd^2x^m=j^2d^2x^kd^2x^md^2x^i`$, then the only solution is to impose $`d^2x^id^2x^k=0`$ and to set forward the additional rule declaring that any expression containing four or more operators $`d`$ must vanish identically. With this set of rules we can check that $`d^3=0`$ on all the forms, whatever their grade or degree. Let us show how such calculus works on the example of a 1-form $`\omega =\omega _kdx^k`$: $$d(\omega _kdx^k)=(_i\omega _k)dx^idx^k+\omega _kd^2x^k;$$ (21) $$d^2(\omega _kdx^k)=(_m_i\omega _k)dx^mdx^idx^k+(_i\omega _k)(d^2x^idx^k+jdx^id^2x^k)+_i\omega _kdx^id^2x^k;$$ (22) after exchanging the summation indices $`i`$ and $`k`$ in two last terms and using the fact that $`j+1=j^2`$ and the commutation relations between $`dx^k`$ and $`d^2x^i`$, we can write $$d^2(\omega _kdx^k)=(_m_i\omega _k)dx^mdx^idx^k+(_i\omega _k_k\omega _i)d^2x^idx^k.$$ (23) where it is interesting to note how the usual anti-symmetric exterior differential appears as a part of the whole expression. It is also easy to check that $$Im(d)Ker(d^2),\mathrm{and}Im(d^2)Ker(d)$$ (24) 4. COVARIANT DIFFERENTIAL, CONNECTION AND CURVATURE. Let A be an associative algebra with unit element, and let H be a free left module over this algebra. Let $`A`$ be a A-valued 1-form defined on a differential manifold $`M`$, and let $`\mathrm{\Phi }`$ be a function on the manifold $`M`$ with values in the module H. We shall introduce the covariant differential as usual: $$D\mathrm{\Phi }:=d\mathrm{\Phi }+A\mathrm{\Phi };$$ (25) If the module is a free one, any of its elements $`\mathrm{\Phi }`$ can be represented by an appropriate element of the algebra acting on a fixed element of H, so that one can always write $`\mathrm{\Phi }=B\mathrm{\Phi }_o`$; then the action of the group of automorphisms of H can be translated as the action of the same group on the algebra A. Let $`U`$ be a function defined on $`M`$ with its values in the group of the automorphisms of H. The definition of a covariant differential is equivalent with the requirement $`DU^1B=U^1DB`$; as in the usual case, this leads to the following well-known transformation for the connection 1-form $`A`$ : $$AU^1AU+U^1dU;$$ (26) But here, unlike in the usual theory, the second covariant differential $`D^2\mathrm{\Phi }`$ is not an automorphism: as a matter of fact, we have: $$D^2\mathrm{\Phi }=d(d\mathrm{\Phi }+A\mathrm{\Phi })+A(d\mathrm{\Phi }+A\mathrm{\Phi })=d^2\mathrm{\Phi }+dA\mathrm{\Phi }+jAd\mathrm{\Phi }+Ad\mathrm{\Phi }+A^2\mathrm{\Phi };$$ (27) the expression containing $`d^2\mathrm{\Phi }`$ and $`d\mathrm{\Phi }`$ ; whereas $`D^3\mathrm{\Phi }`$ is an automorphism indeed, because it contains only $`\mathrm{\Phi }`$ multiplied on the left by an algebra-valued 3-form: $$D^3\mathrm{\Phi }=d(D^2\mathrm{\Phi })+A(D^2\mathrm{\Phi }),$$ (28) which gives explicitly: $$d(d^2\mathrm{\Phi }+dA\mathrm{\Phi }+jAd\mathrm{\Phi }+A^2\mathrm{\Phi })+A(d^2\mathrm{\Phi }+dA\mathrm{\Phi }+jAd\mathrm{\Phi }+Ad\mathrm{\Phi }+A^2\mathrm{\Phi })$$ (29) With a direct calculus one observes that all the terms containing $`d\mathrm{\Phi }`$ or $`d^2\mathrm{\Phi }`$ simplify because of the identity $`1+j+j^2=0`$, leaving only $$D^3\mathrm{\Phi }=(d^2A+d(A^2)+AdA+A^3)\mathrm{\Phi }=(D^2A)\mathrm{\Phi }:=\mathrm{\Omega }\mathrm{\Phi };$$ (30) Obviously, because $`D(U^1\mathrm{\Phi })=U^1(D\mathrm{\Phi })`$, one also has: $`D^3(U^1\mathrm{\Phi })=U^1(D^3\mathrm{\Phi })=U^1\mathrm{\Omega }\mathrm{\Phi }=U^1\mathrm{\Omega }UU^1\mathrm{\Phi }`$, which proves that the 3-form $`\mathrm{\Omega }`$ transforms as usual, $`\mathrm{\Omega }U^1\mathrm{\Omega }U`$ when the connection 1-form transforms according to the law: $`AU^1AU+U^1dU`$. It can be also proved by a direct calculus that the curvature 3-form $`\mathrm{\Omega }`$ does vanish identically for $`A=U^1dU`$. This computation illustrates very well the technique of the $`Z_3`$-graded exterior differential calculus introduced above: as a matter of fact, one has $$d(U^1dU)=dU^1dU+U^1d^2U,$$ (31) so that the term corresponding to $`d^2A`$ gives: $$d^2(U^1dU)=d^2U^1dU+jdU^1d^2U+dU^1d^2U;$$ (32) next, the term corresponding to $`d(A^2)=d(U^1dUU^1dU)`$ gives $$dU^1dUU^1dU+U^1d^2UU^1dU+jU^1dUdU^1dU+jU^1dUU^1d^2U;$$ (33) whereas $$AdA=U^1dUdU^1dU+U^1dUU^1d^2U;$$ (34) and finally, the term corresponding to $`A^3=U^1dUU^1dUU^1dU`$ can be written as $`dU^1dUU^1dU`$ by virtue of the identity $`dUU^1=UdU^1`$ which follows from the Leibniz rule applied to $`UU^1=\mathrm{𝟏}`$, i.e. $`d(UU^1)=dUU^1+UdU^1=0.`$ Using this identity whenever possible, and replacing $`1+j`$ by $`j^2`$, we can reduce the whole expression to the following sum of three terms $$d^2U^1dU+U^1d^2UU^1dUj^2U^1dUdU^1dU$$ (35) whose vanishing does not at all seem obvious. However, it is not very difficult to prove that this expression is identically null. First of all, it is enough to prove the vanishing of the expression $`d^2U^1+U^1d^2UU^1j^2U^1dUdU^1`$ , because all the three terms contain the same factor $`dU`$ on the right; then, by multiplying on the left by $`U`$, we get $$Ud^2U^1+d^2UU^1j^2dUdU^1$$ (36) At this point let us note that $`d^2(UU^1)=d^2(Id)=0`$, but then, according to our $`Z_3`$-graded Leibniz rule, $$d^2(UU^1)=d(dUU^1+UdU^1)=d^2UU^1+jdUdU^1+dUdU^1+Ud^2U^1;$$ (37) so that $`Ud^2U^1+d^2UU^1=dUdU^1jdUdU^1`$, and the result can be written as $$dUdU^1jdUdU^1j^2dUdU^1=0$$ (38) because here again the common factor $`dUdU^1`$ is multiplied by $`(1+j+j^2)=0`$ , which completes the proof. 5. EXPRESSIONS IN LOCAL COORDINATES. The curvature 3-form $`\mathrm{\Omega }=d^2A+d(A^2)+AdA+A^3`$ is of grade $`0`$; therefore it must be decomposed along the elements $`dx^idx^kdx^m`$ and $`d^2x^idx^k`$. Here is how we can compute its components in a local coordinate system. By definition, $`A=A_idx^i`$, so we have: $$dA=_iA_kdx^idx^k+A_kd^2x^k;$$ (39) $$d^2A=_m_iA_kdx^mdx^idx^k+_iA_kd^2x^idx^k+j_iA_kdx^id^2x^k+_iA_kdx^id^2x^k;$$ (40) After replacing $`1+j`$ by $`j^2`$, and taking into account the relation $`dx^kd^2x^i=jd^2x^idx^k`$, we get: $$d^2A=(_m_iA_k)dx^mdx^idx^k+(_iA_k_kA_i)d^2x^idx^k;$$ (41) $$\mathrm{Then},d(A^2)+AdA=dAA+jAdA+AdA=dAAj^2AdA,$$ (42) which leads easily to $$(_iA_kA_m)dx^idx^kdx^mj^2(A_m_iA_k)dx^mdx^idx^k+A_kA_md^2x^kdx^mj^2A_mA_kdx^md^2x^k;$$ (43) and due to the relations $`dx^md^2x^k=jd^2x^kdx^m`$ et $`dx^mdx^idx^k=jdx^idx^kdx^m`$, $$d(A^2)+AdA=(A_m_iA_k_iA_kA_m)dx^mdx^idx^k+(A_kA_mA_mA_k)d^2x^kdx^m.$$ (44) Finally, as $`A^3=A_iA_kA_mdx^idx^kdx^m`$, the curvature 3-form can be written in local coordinates as follows: $$\mathrm{\Omega }=d^2A+d(A^2)+AdA+A^3=\mathrm{\Omega }_{ikm}dx^idx^kdx^m+F_{ik}d^2x^idx^k$$ (45) $$\mathrm{where}\mathrm{\Omega }_{ikm}:=_i_kA_m+A_i_kA_m_kA_mA_i+A_iA_kA_m,$$ (46) $$\mathrm{and}F_{ik}:=_iA_k_kA_i+A_iA_kA_kA_i;$$ (47) In $`F_{ik}`$ one can easily recognize the 2-form of curvature of the usual gauge theories. We know that the expression $`F_{ik}`$ is covariant with respect to the gauge transformations; on the other hand, the 3-form $`\mathrm{\Omega }`$ is also covariant; therefore, the local expression $`\mathrm{\Omega }_{ijk}`$ must be covariant, too. As a matter of fact, it can be expressed as a combination of covariant derivatives of the 2-form $`F_{ik}`$. In order to find the covariant expression of $`\mathrm{\Omega }_{ikm}`$, it suffices to recall that due to the particular symmetry of the ternary exterior product $`dx^idx^kdx^m`$, we can replace $`\mathrm{\Omega }_{ikm}`$ by $`\frac{1}{3}(\mathrm{\Omega }_{ikm}+j^2\mathrm{\Omega }_{kmi}+j\mathrm{\Omega }_{mik})`$ and analyze the abelian case, when this expression reduces itself to $`\mathrm{\Omega }_{ikm}=_i_kA_m`$. Substituting for $`_i_kA_m`$ the equivalent expression $`\frac{1}{3}(_i_kA_m+j^2_k_mA_i+j_m_iA_k)`$ and then $`\frac{1}{3}(j(_k_mA_i_i_kA_m)+j(_m_iA_k_i_kA_m))`$ , because $`1=jj^2`$, we can easily recognize $`\frac{1}{3}(j_i[_mA_k_kA_m]+j^2_k[_mA_i_iA_k])`$ ; which in a general non-abelian case must lead to the following expression: $$\mathrm{\Omega }_{ikm}=\frac{1}{3}[jD_iF_{mk}+j^2D_kF_{mi}],$$ (48) or, equivalently, $$\mathrm{\Omega }_{ikm}=\frac{1}{6}[D_iF_{mk}+D_kF_{mi}]+\frac{i\sqrt{3}}{6}[D_iF_{mk}D_kF_{mi}]$$ (49) 6. DUALITY, INTEGRATION, VARIATIONAL PRINCIPLE. The natural symmetry between $`j`$ et $`j^2`$ , which leads to the possibility of choosing one of these two complex numbers as the generator of the group $`Z_3`$ , and simultaneous interchanging the rôles between the grades $`1`$ and $`2`$, suggests that we could extend the notion of complex conjugation $`j(j)^{}:=j^2`$, with $`((j)^{})^{}=j`$, to the algebra of $`Z_3`$\- graded exterior forms and the operator $`d`$ itself. It does not seem reasonable to use the ”second differentials” $`d^2x^i`$ as the objects conjugate to the ”first differentials” $`dx^i`$, because the rules of $`Z_3`$-graded exterior differentiation we have imposed break the symmetry between these two kinds of differentials: remember that the products $`dx^idx^k`$, and $`dx^idx^kdx^m`$ are admitted, while we require that $`d^2x^id^2x^k`$ and $`d^2x^id^2x^kd^2x^m`$ must vanish. This suggests the introduction of a ”conjugate” differential $`\delta `$ of grade $`2`$, the image of the differential $`d`$ under the conjugation $``$, satisfying the following conjugate relations: $$\delta x^i\delta x^k\delta x^m=j^2\delta x^k\delta x^m\delta x^i,\delta x^i\delta ^2x^k=j^2\delta ^2x^k\delta x^i.$$ (50) One notes that $`\delta ^2x^k`$ is of grade 1 ($`2+2=4=1(mod3)`$). All the relations existing between the operator $`d`$ and the exterior forms generated by $`dx^i`$ and $`d^2x^k`$ are faithfully reproduced under the conjugation $``$ if we consider the $`Z_3`$-graded algebra generated by the entities $`\delta x^i`$ and $`\delta ^2x^k`$ as a right module over the algebra of functions F(M) , with the operator $`\delta `$ acting on the right on this module. The rules $`d^3=0`$ and $`\delta ^3=0`$ suggest their natural extension: $$d\delta =\delta d=0$$ (51) We would like to be able to form quadratic expressions that could define a scalar product; to do this, we should postulate that the algebra generated by the elements $`dx^i,d^2x^k`$ and its conjugate algebra generated by the elements $`\delta x^i,\delta x^k`$ commute with each other. Then, we can define scalar products for the forms of maximal degree $`3`$: $`<\omega \varphi >:=\omega \varphi `$ , and integrating this result with respect to the usual volume element defined on the manifold $`M`$, which gives explicitly: $$\overline{\omega }_{ikm}\varphi _{prs}<\delta x^i\delta x^k\delta x^mdx^pdx^rdx^s>,\overline{\psi }_{ik}\chi _{mn}<\delta x^i\delta ^2x^kd^2x^mdx^n>$$ (52) What remains now is to determine the scalar products of the basis of forms; in order to assure the hermiticity of the product, one can always choose an ”orthonormal” basis in which we should have: $$<\delta x^i\delta x^k\delta x^mdx^pdx^rdx^s>=\delta _s^i\delta _r^k\delta _p^m\mathrm{et}<\delta x^i\delta ^2x^kd^2x^mdx^n>=\mu \delta _n^i\delta _m^k.$$ (53) Here the first scalar product is normed to $`1`$, and $`\mu `$ is the ratio between the two types of ”elementary volume”. We shall consider the two types of forms of degree $`3`$, $`dx^idx^kdx^m`$ and $`d^2x^pdx^r`$, as being mutually orthogonal. Because in both cases the differentials of the volume forms are identically null, we can formally apply Stokes’ formula with a vanishing contribution of the ”surface term”, which enables us to derive the classical Euler-Lagrange equations. With the Lagrangean density defined on the manifold as $`<\mathrm{\Omega }\mathrm{\Omega }>`$ we get: $$L=[\frac{4}{3}D_iF_{mk}D^iF^{mk}\frac{2}{3}D_kF_{mi}D^iF^{mk}]+\mu [4F_{ik}F^{ik}]$$ (54) (We suppose here that the manifold $`M`$ is flat, and we raise and lower the indeces by means of the metric tensor $`\delta _{ik}`$ and its inverse $`\delta ^{kl}`$). The corresponding Euler-Lagrange are then easily deduced: $$D^mD^iD_iF_{mk}D^iD^kD_mF_{ik}+\frac{3\mu }{4}(D^iF_{ik})=0$$ (55) In the abelian case, and with the usual choice of gauge ($`^iA_i=0`$), these equations reduce themselves to: $$A_k+\frac{3\mu }{4}A_k=0$$ (56) It is worthwhile to note that this scheme can be readily generalized to a $`Z_N`$-graded case, with an exterior differential operator satisfying $`d^N=0`$, leading to the invariants built with the higher-order covariant derivatives of the usual curvature tensor, leading to the equations containing higher powers of the Laplace-Beltrami operator. The equations of this type have been studied before (\[L.Vainerman, (1974)\]) In quantum field theory, they appear in the perturbative developments of non-linear effective Lagrangeans (the so-called Schwinger terms). In such expansions, the consecutive powers of the Laplacian (or the d’Alembertian) operator appear, with the coeficients depending on the particular choice of Lagrangean. It would be interesting to compare the coefficients appearing in these theories with the coefficients depending on the choice of normalization of our various ”volume elements” of higher orders, which appear in the above model. Nevertheless, from the mathematical point of view it seems that a theory based on the $`Z_3`$-grading is particularly interesting because of its exceptional character related to the fact that the group $`S_3`$ is the last of the permutation groups that has a faithful representation in complex plane, which could lead to some special applications in physics, e.g. a better description of the quark fields ( cf.\[R.Kerner, 1992\], \[B. Le Roy, 1995\], and \[ V.Abramov, R.Kerner, B.Le Roy (1995)\].). Another possible application of this formalism may be a new description of differential operators on Riemannian manifolds. Consider a linear operation $`d`$ that produces a symmetric $`2`$-form from a given $`1`$-form, $`(d\xi )_{ij}=_i\xi _j+_j\xi _i`$ The complex thus includes $`1`$-forms and symmetric $`2`$-forms (metrics). Now, let $`d`$ of a metric be defined as its Christoffel connection, and $`d`$ of the latter as Riemannian curvature. Then we have $`d^20`$, but $`d^3=0`$. The applications to gauge theories with well defined non-abelian gauge groups will be the object of our forthcoming articles. Références 1. M. Dubois-Violette, R. Kerner, J. Madore, Journ.Math.Phys.31 p.316, ibid, p. 323 (1990) 2. M. Dubois-Violette, J. Madore, R. Kerner, Class. and Quantum Gravity, 8, p. 1077 - 1089 (1991) 3. A. Connes, J. Lott, Nucl. Phys. B (Proc.Suppl.) 18 , p.29 (1990) 4. R. Coquereaux, G. Esposito-Farese, P. Vaillant, Nucl. Phys. B, 353 p. 689 (1991) 5. R. Kerner, Journal of Geometry and Physics, 11 (1 - 4) p.325 (1993) 6. A.H. Chamseddine, G. Felder, J. Frohlich, Phys. Lett. B 296, p. 109 (1992) 7. V. Abramov, R. Kerner, B. Le Roy, to appear (1995) 8. R. Kerner, C.R.Acad.Sci., Paris, 312 (Ser. II) , p.191-196, (1991) 9. R. Kerner, Journal of Math. Physics, 33 (1), p. 411 - 416, (1992) 10. R. Kerner, in ”Generalized Symmetries in Physics”, H.-D. Doebner, V.K. Dobrev and A.G. Ushveridze ed., World Scientific, p.375 - 394, (1994) 11.V.Takhtajan, Comm.Math.Phys., 160, p. 295 (1994) 12. B. Le Roy, C.R.Acad.Sci., to appear (1995) 13. L. Vainerman, Soviet.Math.Dokl, (1974) Acknowledgements . Numerous enlightening discussions with V. Abramov, B. Le Roy, L. Vainerman, M. Dubois-Violette and J. Madore are gratefully acknowledged. Laboratoire GCR, Université Pierre-et-Marie-Curie - CNRS URA 769, Tour 22, Boíte 142, 4, Place Jussieu, 75005 Paris. e-mail address : rk@ccr.jussieu.fr FAX : (33 1) 44 27 72 87 .
warning/0004/physics0004016.html	ar5iv	text	# Fast Algorithm for Finding the Eigenvalue Distribution of Very Large Matrices ## I Introduction The calculation of the distribution of eigenvalues of very large matrices is a central problem in quantum physics. This distribution determines the thermodynamic properties of the system (see below). It is directly related to the single-particle density of states (DOS) or Green’s function. In a one-particle (e.g., one-electron) description knowledge of the DOS suffices to compute the transport properties . The most direct method to compute the DOS, i.e. all the eigenvalues, is to diagonalize the matrix $`H`$ representing the Hamiltonian of the system. This approach has two obvious limitations: The number of operations increases as the third power of the dimension $`D`$ of $`H`$ and, perhaps most importantly, the amount of memory required by state-of-the-art algorithms grows as $`D^2`$ . This scaling behavior limits the application of this approach to matrices of dimension $`D=𝒪(10000)`$, which is too small for many problems of interest. What is needed are methods that scale linearly with $`D`$. There has been considerable interest in developing “fast” (i.e. $`𝒪(D)`$) algorithms to compute the DOS and other similar quantities. One such algorithm and an application of it to electron motion in disordered alloy models was given by Alben et al. . In this approach the DOS is obtained by solving the time-dependent Schrödinger equation (TDSE) of a particle moving on a lattice, followed by a Fourier transform of the retarded Green’s function . Using the unconditionally stable split-step Fast Fourier Transform (FFT) method to solve the TDSE, it was shown that the eigenvalue spectrum of a particle moving in continuum space can be computed in the same manner . Fast algorithms of this kind proved useful to study various aspects of localization of waves and other one-particle problems . Application of these ideas to quantum many-body systems triggered further development of flexible and efficient methods to solve the TDSE. Based on Suzuki’s product formula approach, an unconditionally stable algorithm was developed and used to compute the time-evolution of two-dimensional S=1/2 Heisenberg-like models . Results for the DOS of matrices of dimension $`D1000000`$ where reported . A potentially interesting feature of these fast algorithms is that they may run very efficiently on a quantum computer . A common feature of these fast algorithms is that they solve the TDSE for a sample of randomly chosen initial states. The efficiency of this approach as a whole relies on the hypothesis (suggested by the central limit theorem) that satisfactory accuracy can be achieved by using a small sample of initial states. Experience not only shows that this hypothesis is correct, it strongly suggests that for a fixed sample size the statistical error on physical quantities such as the energy and specific heat decreases with the dimension $`D`$ of the Hilbert space . In view of the general applicability of these fast algorithms to a wide variety of quantum problems it seems warranted to analyze in detail their properties and the peculiar $`D`$ dependence in particular. In Sections II and III we recapitulate the essence of the approach. We present a rigorous estimate for the mean square error (variance) on the trace of a matrix. In Section IV we describe the imaginary-time version of the method. The statistical analysis of the numerical data is discussed in Section V. Section VI describes the model systems that are used in our numerical experiments. The algorithm used to solve the TDSE is reviewed in Section VII. In Section VIII we derive rigorous bounds on the accuracy with which all eigenvalues can be determined and demonstrate that this accuracy decreases linearly with the time over which the TDSE is solved. The results of our numerical calculations are presented in Section IX and our conclusions are given in Section X. ## II Theory The trace of a matrix $`A`$ acting on a $`D`$-dimensional Hilbert-space spanned by an orthonormal set of states $`\{\|\varphi _n\}`$ is given by $`TrA={\displaystyle \underset{n=1}{\overset{D}{}}}\varphi _n\|A\varphi _n.`$ (1) Note that according to (1) we have $`Tr\mathrm{\hspace{0.17em}1}=D`$. If $`D`$ is very large one might think of approximating Eq. (1) by sampling over a subset of $`K`$ ($`KD`$)“important” basis vectors. The problem with this approach is that the notion “important” may be very model-dependent. Therefore it is better to sample in a different manner. We construct a random vector $`\|\psi `$ by choosing $`D`$ complex random numbers, $`c_nf_n+ig_n`$, with mean $`0`$, for $`n=1\mathrm{}D`$, so $`\|\psi ={\displaystyle \underset{n=1}{\overset{D}{}}}c_n\|\varphi _n,`$ (2) and calculate $`\psi \|A\psi ={\displaystyle \underset{n,m=1}{\overset{D}{}}}c_m^{}c_n^{}\varphi _m\|A\varphi _n.`$ (3) If we now sample over $`S`$ realizations of the random vectors $`\{\psi \}`$ and calculate the average, we obtain $`{\displaystyle \frac{1}{S}}{\displaystyle \underset{p=1}{\overset{S}{}}}\psi _p\|A\psi _p={\displaystyle \frac{1}{S}}{\displaystyle \underset{p=1}{\overset{S}{}}}{\displaystyle \underset{n,m=1}{\overset{D}{}}}c_{m,p}^{}c_{n,p}^{}\varphi _m\|A\varphi _n.`$ (4) Assuming that there is no correlation between the random numbers in different realizations and that the random numbers $`f_{n,p}`$ and $`g_{n,p}`$ are drawn from an even and symmetric (both with respect to each variable) probability distribution (see Appendix A for more details), we have $`\underset{S\mathrm{}}{lim}{\displaystyle \frac{1}{S}}{\displaystyle \underset{p=1}{\overset{S}{}}}c_{m,p}^{}c_{n,p}^{}=E\left(\|c\|^2\right)\delta _{m,n},`$ (5) where $`E(.)`$ denotes the expectation value with respect to the probability distribution used to generate the $`c_{n,p}`$’s. In the r.h.s of (5) the subscripts of $`c_{n,p}`$ have been dropped to indicate that the expectation value does not depend on $`n`$ or $`p`$. It follows immediately that $`\underset{S\mathrm{}}{lim}{\displaystyle \frac{1}{S}}{\displaystyle \underset{p=1}{\overset{S}{}}}\psi _p\|A\psi _p=E\left(\|c\|^2\right)TrA=E\left(\|c\|^2\right){\displaystyle \underset{n=1}{\overset{D}{}}}\varphi _n\|A\varphi _n,`$ (6) showing that we can compute the trace of $`A`$ by sampling over random states $`\{\psi _p\}`$, provided there is an efficient algorithm to calculate $`\psi _p\|A\psi _p`$ (see Section VII). According to the central limit theorem, for a large but finite $`S`$ we have $`{\displaystyle \frac{1}{S}}{\displaystyle \underset{p=1}{\overset{S}{}}}c_{m,p}^{}c_{n,p}^{}=E\left(\|c\|^2\right)\delta _{m,n}+𝒪\left({\displaystyle \frac{1}{\sqrt{S}}}\right),`$ (7) meaning that the statistical error on the trace vanishes like $`1/\sqrt{S}`$, which is not surprising. What is surprising is that one can prove a much stronger result as follows. Let us first normalize the $`c_{n,p}`$’s so that, for all $`p`$, $`{\displaystyle \underset{n=1}{\overset{D}{}}}\|c_{n,p}\|^2=1.`$ (8) This innocent looking step has far reaching consequences. First we note that the normalization renders the method exact in (the rather trivial) case that the matrix $`A`$ is proportional to the unit matrix. The price we pay for this is that for fixed $`p`$, the $`c_{n,p}`$ are now correlated but that does not cause problems (see Appendix A). Second it follows that $`E\left(\|c\|^2\right)=1/D`$. Obviously the error can be written as $`TrA{\displaystyle \frac{D}{S}}{\displaystyle \underset{p=1}{\overset{S}{}}}\psi _p\|A\psi _p=TrRA,`$ (9) where $`R_{m,n}\delta _{m,n}{\displaystyle \frac{D}{S}}{\displaystyle \underset{p=1}{\overset{S}{}}}c_{m,p}^{}c_{n,p}^{},`$ (10) is a traceless (due to Eq. (8)) Hermitian matrix of random numbers. We put $`X=TrRA`$ and compute $`E(\|X\|^2)`$. The result for the general case can be found in Appendix A. For a uniform distribution of the $`c_{n,p}`$’s on the hyper-sphere defined by $`_{n=1}^D\|c_{n,p}\|^2=1`$ the expression simplifies considerably and we find $`E\left(\left\|TrRA\right\|^2\right)`$ $`={\displaystyle \frac{DTrA^{}A\|TrA\|^2}{S(D+1)}},`$ (11) an exact expression for the variance in terms of the sample size $`S`$, the dimension $`D`$ of the matrix $`A`$ and the (unknown) constants $`TrA^{}A`$ and $`\|TrA\|`$. Invoking a generalization of Markov’s inequality $`𝐏(\|X\|^2a){\displaystyle \frac{E(\|X\|^2)}{a}};a>0,`$ (12) where $`𝐏(Q)`$ denotes the probability for the statement $`Q`$ to be true. We find that the probability that $`\|TrRA\|^2`$ exceeds a fraction $`a`$ of $`\|TrA\|^2`$ is bounded by $`𝐏({\displaystyle \frac{\|TrRA\|^2}{\|TrA\|^2}}a){\displaystyle \frac{1}{aS(D+1)}}{\displaystyle \frac{DTrA^{}A\|TrA\|^2}{\|TrA\|^2}};a>0,`$ (13) or, in other words, the relative statistical error $`e_A`$ on the estimator of the trace of $`A`$ is given by $`e_A\sqrt{{\displaystyle \frac{DTrA^{}A\|TrA\|^2}{S(D+1)\|TrA\|^2}}},`$ (14) if $`\|TrA\|>0`$. We see that $`e_A=0`$ if $`A`$ is proportional to a unit matrix. From (14) it follows that, in general, we may expect $`e_A`$ to vanish with the square root of $`SD`$. The prefactor is a measure for the relative spread of the eigenvalues of $`A`$ and is obviously model dependent. The dependence of $`e_A`$ on $`S`$, $`D`$ and the spectrum of $`A`$ is corroborated by the numerical results presented below. It is also of interest to examine the effect of not normalizing the $`c_{n,p}`$’s. A calculation similar to the one that lead to the above results yields $`e_A=\sqrt{{\displaystyle \frac{TrA^{}A}{S\|TrA\|^2}}}.`$ (15) Clearly this bound is less sharp and does not vanish if $`A`$ is proportional to a unit matrix. ## III Real-time method The distribution of eigenvalues or density of states (DOS) of a quantum system is defined as $$𝒟(ϵ)=\underset{n=1}{\overset{D}{}}\delta (ϵE_n)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}e^{itϵ}Tre^{itH}dt,$$ (16) where $`H`$ is the Hamiltonian of the system and $`n`$ runs over all the eigenvalues of $`H`$. The DOS contains all the physical information about the equilibrium properties of the system. For instance the partition function, the energy, and the heat capacity are given by $`Z`$ $`={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑ϵ𝒟(ϵ)e^{\beta ϵ},`$ (17) $`E`$ $`={\displaystyle \frac{1}{Z}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑ϵϵ𝒟(ϵ)e^{\beta ϵ},`$ (18) $`C`$ $`=\beta ^2\left({\displaystyle \frac{1}{Z}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑ϵϵ^2𝒟(ϵ)e^{\beta ϵ}E^2\right),`$ (19) respectively. Here $`\beta =1/k_BT`$ and $`k_B`$ is Boltzmann’s constant (we put $`k_B=1`$ and $`\mathrm{}=1`$ from now on). As explained above the trace in the integral (16) can be estimated by sampling over random vectors. For the statistical error analysis discussed below it is convenient to define a DOS-per-sample by $$d_p(ϵ)\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}e^{itϵ}\psi _p\|e^{itH}\psi _p𝑑t,$$ (20) where the subscript $`p`$ labels the particular realization of the random state $`\|\psi _p`$. The DOS is then given by $$𝒟(ϵ)=\underset{S\mathrm{}}{lim}\frac{1}{S}\underset{p=1}{\overset{S}{}}d_p(ϵ).$$ (21) Schematically the algorithm to compute $`d_p(ϵ)`$ consists of the following steps: 1. Generate a random state $`\|\psi _p(0)`$, set $`t=0`$. 2. Copy this state to $`\|\psi _p(t)`$. 3. Calculate $`\psi _p(0)\|\psi _p(t)`$ and store the result. 4. Solve the TDSE for a small time step $`\tau `$, replacing $`\|\psi _p(t)`$ by $`\|\psi _p(t+\tau )`$ (see Section VII for model specific details). 5. Repeat $`N`$ times from Step 3. 6. Perform a Fourier transform on the tabulated result and store $`d_p(ϵ)`$. In practice the Fourier transform in Eq. (16) is performed by the Fast Fourier Transform (FFT). We use a Gaussian window to account for the finite time $`\tau N`$ used in the numerical time-integration of the TDSE. The number of time step $`N`$ determines the accuracy with which the eigenvalues can be computed. In Section VIII we prove that this systematic error in the eigenvalues vanishes as $`1/\tau N`$. Since for any reasonable physical system (or finite matrix) the smallest eigenvalue $`E_0`$ is finite, for all practical purposes $`d_p(ϵ)=0`$ for $`ϵ<ϵ_0<E_0`$. The value of $`ϵ_0`$ is easily determined by examination of the bottom of spectrum. To compute $`Z`$, $`E`$, or $`C`$ we simply replace the interval $`[\mathrm{},+\mathrm{}]`$ by $`[ϵ_0,+\mathrm{}]`$. ## IV Imaginary-time Method The real-time approach has the advantage that it yields information on all eigenvalues and can be used to compute both dynamic and static properties without suffering from numerical instabilities. However for the computation of the thermodynamic properties, the imaginary-time version is more efficient. We will use the imaginary-time method as an independent check on the results obtained by the real-time algorithm. Repeating the steps that lead to Eq. (17) we find $`Z`$ $`=Tr\mathrm{exp}(\beta H)`$ (22) $`=\underset{S\mathrm{}}{lim}{\displaystyle \frac{1}{S}}{\displaystyle \underset{p=1}{\overset{S}{}}}\psi _p\|\mathrm{exp}(\beta H)\psi _p,`$ (23) with similar expressions for $`E`$ and $`C`$. Furthermore we have $`\psi _p\|H^ne^{\beta H}\psi _p=e^{\beta H/2}\psi _p\|H^ne^{\beta H/2}\psi _p,`$ (24) assuming $`H`$ is Hermitian. Therefore we only need to propagate the random state for an imaginary time $`\beta /2`$ instead of $`\beta `$. Furthermore we do not need to perform an FFT. Disregarding these minor differences, the algorithm is the same as in the real-time case with $`\tau `$ replaced by $`i\tau `$. ## V Error Analysis Estimating the statistical error on the partition function $`Z`$ is easy because it depends linearly on the trace of the (imaginary) time evolution operator. However the error on $`E`$ and $`C`$ depends on this trace in a more complicated manner and this fact has to be taken into account. First we define $`z_p`$ $`{\displaystyle _{ϵ_0}^{\mathrm{}}}𝑑ϵd_p(ϵ)e^{\beta ϵ},`$ (25) $`h_p`$ $`{\displaystyle _{ϵ_0}^{\mathrm{}}}𝑑ϵd_p(ϵ)ϵe^{\beta ϵ},`$ (26) $`w_p`$ $`{\displaystyle _{ϵ_0}^{\mathrm{}}}𝑑ϵd_p(ϵ)ϵ^2e^{\beta ϵ},`$ (27) for the real-time method and $`z_p`$ $`\psi _p\|e^{\beta H}\psi _p,`$ (28) $`h_p`$ $`\psi _p\|He^{\beta H}\psi _p,`$ (29) $`w_p`$ $`\psi _p\|H^2e^{\beta H}\psi _p,`$ (30) for the imaginary-time method. For each value of $`\beta `$ we generate the data $`\{z_p\}`$, $`\{h_p\}`$, and $`\{w_p\}`$, for $`p=1,\mathrm{},S`$. For both cases we have $`Z`$ $`=\underset{S\mathrm{}}{lim}\overline{z},`$ (31) $`E`$ $`=\underset{S\mathrm{}}{lim}{\displaystyle \frac{\overline{h}}{\overline{z}}},`$ (32) $`C`$ $`=\underset{S\mathrm{}}{lim}\beta ^2\left({\displaystyle \frac{\overline{w}}{\overline{z}}}{\displaystyle \frac{\overline{h}^2}{\overline{z}^2}}\right),`$ (33) where $`\overline{x}S^1_{p=1}^Sx_p`$. The standard deviations on $`\overline{z}`$, $`\overline{h}`$, and $`\overline{w}`$ are given by $`\delta z`$ $`=\sqrt{{\displaystyle \frac{var(z)}{S1}}},`$ (34) $`\delta h`$ $`=\sqrt{{\displaystyle \frac{var(h)}{S1}}},`$ (35) $`\delta w`$ $`=\sqrt{{\displaystyle \frac{var(w)}{S1}}},`$ (36) where $`var(x)\overline{x^2}\overline{x}^2`$ denotes the variance on the data $`\{x_p\}`$. However the sets of data $`\{z_p\}`$, $`\{h_p\}`$ and $`\{w_p\}`$ are correlated since they are calculated from the same set $`\{\|\psi _p\}`$. These correlations in the data are accounted for by calculating the covariance matrix $`M_{k,l}`$ ($`k,l=1,\mathrm{},3`$) the elements of which are given by $`\overline{x_kx_l}\overline{x_k}\overline{x_l}`$, where $`\{x_1\}`$, $`\{x_2\}`$, and $`\{x_3\}`$ are a shorthand for $`\{z_p\}`$, $`\{h_p\}`$, and $`\{w_p\}`$ respectively. The estimates for the errors in $`Z`$, $`E`$ and $`C`$ are given by $`\delta Z^2`$ $`={\displaystyle \frac{1}{S1}}\delta z^2,`$ (37) $`\delta E^2`$ $`={\displaystyle \frac{1}{S1}}{\displaystyle \underset{k,l=1}{\overset{3}{}}}M_{k,l}{\displaystyle \frac{d\overline{E}}{d\overline{x_k}}}{\displaystyle \frac{d\overline{E}}{d\overline{x_l}}},`$ (38) $`\delta C^2`$ $`={\displaystyle \frac{1}{S1}}{\displaystyle \underset{k,l=1}{\overset{3}{}}}M_{k,l}{\displaystyle \frac{d\overline{C}}{d\overline{x_k}}}{\displaystyle \frac{d\overline{C}}{d\overline{x_l}}},`$ (39) where $`\overline{E}=\overline{h}/\overline{z}`$ and and $`\overline{C}=\beta ^2(\overline{w}/\overline{z}\overline{h}^2/\overline{z}^2)`$. ## VI Exactly Solvable Spin $`1/2`$ Models The most direct way to assess the validity of the approach described above is to carry out numerical experiments on exactly solvable models. In this paper we consider three different exactly solvable models, two spin-1/2 chains and a mean-field spin-1/2 model. The former have a complicated spectrum, the latter has a highly degenerate eigenvalue distribution. These spin models differ from those studied elsewhere in that they belong to the class of integrable systems. ### A Spin chains Open spin chains of $`L`$ sites described by the Hamiltonian $`H=J{\displaystyle \underset{i=1}{\overset{L1}{}}}(\sigma _i^x\sigma _{i+1}^x+\mathrm{\Delta }\sigma _i^y\sigma _{i+1}^y)h{\displaystyle \underset{i=1}{\overset{L}{}}}\sigma _i^z,`$ (40) where $`\sigma _i^x`$, $`\sigma _i^y`$, and $`\sigma _i^z`$ denote the Pauli matrices and $`J`$, $`\mathrm{\Delta }`$ and $`h`$ are model parameters, can be solved exactly. They can be reduced to diagonal form by means of the Jordan-Wigner transformation . We have $`H={\displaystyle \underset{i,j=1}{\overset{L}{}}}\left[c_i^+A_{i,j}c_j^{}+{\displaystyle \frac{1}{2}}\left(c_i^+B_{i,j}c_j^++c_j^{}B_{j,i}^{}c_i^{}\right)\right]+hL,`$ (41) where $`c_i^+`$ and $`c_i^{}`$ are spin-less fermion operators and $`A_{i,j}`$ $`=J(1+\mathrm{\Delta })(\delta _{i,j1}+\delta _{i1,j})2h\delta _{i,j},`$ (42) $`B_{i,j}`$ $`=J(1\mathrm{\Delta })(\delta _{i,j1}\delta _{i1,j}),`$ (43) are $`L\times L`$ matrices. By further canonical transformation this Hamiltonian can be written as $`H={\displaystyle \underset{k=1}{\overset{L}{}}}\mathrm{\Lambda }_k\left(n_k{\displaystyle \frac{1}{2}}\right)+{\displaystyle \frac{1}{2}}TrA+hL,`$ (44) where $`n_k`$ is the number operator of state $`k`$ and the $`\mathrm{\Lambda }_k`$’s are given by the solution of the eigenvalue equation $`(AB)(A+B)\varphi _k=\mathrm{\Lambda }_k^2\varphi _k.`$ (45) In the general case this eigenvalue problem of the $`L\times L`$ Hermitian matrix $`(AB)(A+B)`$ is most easily solved numerically. In the present paper we confine ourselves to two limiting cases: The XY model ($`\mathrm{\Delta }=1`$) and the Ising model in a transverse field ($`\mathrm{\Delta }=0`$). ### B Mean field model The Hamiltonian of the mean-field model reads $`H={\displaystyle \frac{J}{L}}{\displaystyle \underset{i>j=1}{\overset{L}{}}}\stackrel{}{\sigma }_i\stackrel{}{\sigma }_jh{\displaystyle \underset{i=1}{\overset{L}{}}}\sigma _i^z,`$ (46) and can be rewritten as $`H=2{\displaystyle \frac{J}{L}}\stackrel{}{S}\stackrel{}{S}2hS^z+{\displaystyle \frac{3}{2}}J,`$ (47) with $`\stackrel{}{S}`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{L}{}}}\stackrel{}{\sigma }_i.`$ (48) The single spin-$`L/2`$ Hamiltonian has eigenvalues $`E_{l,m}`$ $`=2Jl(l+1)/L2hm+{\displaystyle \frac{3}{2}}J,`$ (49) with degeneracy $`n_{l,m}`$ $`={\displaystyle \frac{2l+1}{L/2+l+1}}\left(\begin{array}{c}L\\ L/2l\end{array}\right).`$ (52) This rather trivial model serves as a test for the case of highly degenerate eigenvalues. ## VII Time Evolution For the approach outlined in Sections III and IV to be of practical use it is necessary that the matrix elements of the exponential of $`H`$ can be calculated efficiently. The purpose of this section is to describe how this can be done. The general form of the Hamiltonians of the models we study is $$H=\underset{i,j=1}{\overset{L}{}}\underset{\alpha =x,y,z}{}J_{i,j}^\alpha \sigma _i^\alpha \sigma _j^\alpha \underset{i=1}{\overset{L}{}}\underset{\alpha =x,y,z}{}h_i^\alpha \sigma _i^\alpha ,$$ (53) where the first sum runs over all pairs $`P`$ of spins, $`\sigma _i^\alpha `$ ($`\alpha =x,y,z`$) denotes the $`\alpha `$-th component of the spin-1/2 operator representing the $`i`$-th spin. For both methods, we have to calculate the evolution of a random state, i.e. $`U(\tau )\|\psi \mathrm{exp}(i\tau H)\|\psi `$ or $`U(\tau )\|\psi \mathrm{exp}(\tau H)\|\psi `$ for the real and imaginary time method respectively. We will discuss the real-time case only, the imaginary-time problem can be solved in the same manner. Using the semi-group property $`U(t_1)U(t_2)=U(t_1+t_2)`$ we can write $`U(t)=U(\tau )^m`$ where $`t=m\tau `$. Then the main step is to replace $`U(\tau )`$ by a symmetrized product-formula approximation . For the case at hand it is expedient to take $`U(\tau )\stackrel{~}{U}(\tau )`$ $`e^{i\tau H_z/2}e^{i\tau H_y/2}e^{i\tau H_x}e^{i\tau H_y/2}e^{i\tau H_z/2},`$ (54) where $$H_\alpha =\underset{i,j=1}{\overset{L}{}}J_{i,j}^\alpha \sigma _i^\alpha \sigma _j^\alpha \underset{i=1}{\overset{L}{}}h_i^\alpha \sigma _i^\alpha ;\alpha =x,y,z.$$ (55) Other decompositions work equally well but are somewhat less efficient for the cases at hand. In the real-time approach $`\stackrel{~}{U}(\tau )`$ is unitary and hence the method is unconditionally stable (also the imaginary-time method can be made unconditionally stable). It can be shown that $`U(\tau )\stackrel{~}{U}(\tau )s\tau ^3`$ ($`s>0`$ a constant) , implying that the algorithm is correct to second order in the time step $`\tau `$ . Usually it is not difficult to choose $`\tau `$ so small that for all practical purposes the results obtained can be considered as being “exact”. Moreover, if necessary, $`\stackrel{~}{U}(\tau )`$ can be used as a building block to construct higher-order algorithms . In Appendix B we will derive bounds on the error in the eigenvalues when they are calculated using a symmetric product formula. As basis states $`\{\|\varphi _n\}`$ we take the direct product of the eigenvectors of the $`S_i^z`$ (i.e. spin-up $`\|_i`$ and spin-down $`\|_i`$). In this basis, $`e^{i\tau H_z/2}`$ changes the input state by altering the phase of each of the basis vectors. As $`H_z`$ is a sum of pair interactions it is trivial to rewrite this operation as a direct product of 4x4 diagonal matrices (containing the interaction-controlled phase shifts) and 4x4 unit matrices. Still working in the same representation, the action of $`e^{i\tau H_y/2}`$ can be written in a similar manner but the matrices that contain the interaction-controlled phase-shift have to be replaced by non-diagonal matrices. Although this does not present a real problem it is more efficient and systematic to proceed as follows. Let us denote by $`X`$($`Y`$) the rotation by $`\pi /2`$ of each spin about the $`x`$($`y`$)-axis. As $$e^{i\tau H_y/2}=XX^{}e^{i\tau H_y/2}XX^{}=Xe^{i\tau H_z^{}/2}X^{},$$ (56) it is clear that the action of $`e^{i\tau H_y/2}`$ can be computed by applying to each spin, the inverse of $`X`$ followed by an interaction-controlled phase-shift and $`X`$ itself. The prime in (56) indicates that $`J_{i,j}^z`$ and $`h_i^z`$ in $`H_z`$ have to be replaced by $`J_{i,j}^y`$ and $`h_i^y`$ respectively. A similar procedure is used to compute the action of $`e^{i\tau H_x}`$. We only have to replace $`X`$ by $`Y`$. ## VIII Accuracy of the computed eigenvalues First we consider the problem of how to choose the number of time steps $`N`$ to obtain the DOS with acceptable accuracy. According to the Nyquist sampling theorem employing a sampling interval $`\mathrm{\Delta }t=\pi /\mathrm{max}_i\|E_i\|`$ is sufficient to cover the full range of eigenvalues. On the other hand the time step also determines the accuracy of the approximation $`\stackrel{~}{U}(\tau )`$. Let us call the maximum value of $`\tau `$ which gives satisfactory accuracy $`\tau _0`$ (for the imaginary-time method, this is the only parameter). For the examples treated here $`\tau _0<\mathrm{\Delta }t)`$, implying that we have to use more steps to solve the TDSE than we actually use to compute the FFT. Eigenvalues that differ less than $`\mathrm{\Delta }ϵ=\pi /N\mathrm{\Delta }t`$ cannot be identified properly. However since $`\mathrm{\Delta }ϵN^1`$ we only have to extend the length of the calculation by a factor of two to increase the resolution by the same factor. At first glance the above reasoning may seem to be a little optimistic. It apparently overlooks the fact that if we integrate the TDSE over longer and longer times the error on the wave function also increases (although it remains bounded because of the unconditional stability of the product formula algorithm). In fact it has been shown that in general $$e^{itH}\|\psi (0)\stackrel{~}{U}^m(\tau )\|\psi (0)c\tau ^2t,$$ (57) where $`t=m\tau `$, suggesting that the loss in accuracy on the wave function may well compensate for the gain in resolution that we get by using more data in the Fourier transform. Fortunately this argument does not apply when we want to determine the eigenvalues as we now show. As before we will discuss the real-time algorithm only because the same reasoning (but different mathematical proofs) holds for the imaginary-time case. Consider the time-step operator (52). Using the fact that any unitary matrix can be written as the matrix exponential of a Hermitian matrix we can write $$\stackrel{~}{U}(\tau )=e^{i\tau H_z/2}e^{i\tau H_y/2}e^{i\tau H_x}e^{i\tau H_y/2}e^{i\tau H_z/2}e^{i\tau \stackrel{~}{H}(\tau )}.$$ (58) It is clear that in practice the real-time method yields the spectrum of $`\stackrel{~}{H}(\tau )`$, not the one of $`H`$. Therefore the relevant question is: How much do the spectra of $`\stackrel{~}{H}(\tau )`$ and $`H`$ differ? In Appendix B we give a rigorous proof that the difference between the eigenvalues of $`\stackrel{~}{H}(\tau )`$ and $`H`$ vanishes as $`\tau ^2`$. In other words the value of $`m`$ (or $`t=m\tau `$) has no effect whatsoever on the accuracy with which the spectrum can be determined. Therefore the final conclusion is that the error in the eigenvalues vanishes as $`\tau ^2/N`$ where $`N`$ is the number of data points used in the Fourier transform of $`Tre^{it\stackrel{~}{H}(\tau )}`$. ## IX Results In all our calculations we take $`J=1`$ and $`h=0`$, except for the Ising model in a transverse field, where we take $`h=0.75`$. The random numbers $`c_{n,p}`$ are generated such that the conditions Eqs. (A5) and (A8) are satisfied. We use two different techniques to generate these random numbers: 1. A uniform random number generator produces $`\{f_{n,p}\}`$ and $`\{g_{n,p}\}`$ with $`1f_{n,p},g_{n,p}1`$. We then normalize the vector (see Eq. (8)). 2. The $`c_{n,p}`$’s are obtained from a two-variable (real and imaginary part) Gaussian random number generator and the resulting vector is normalized. Both methods satisfy the basic requirements Eqs. (A5) and (A8) but because the first samples points out of a $`2D`$-dimensional hypercube and subsequently projects the vector onto a sphere, the points are not distributed uniformly over the surface of the unit hyper-sphere. The second method is known to generate numbers which are distributed uniformly over the hyper-surface. Although the first method does not satisfy all the mathematical conditions that lead to the error (14), our numerical experiments with both generators give identical results, within statistical errors of course. Also, within the statistical errors, the results from the imaginary and real-time algorithm are the same. Therefore we only show some representative results as obtained from the real-time algorithm. In Fig. 1 we show a typical result for the DOS $`D(ϵ)`$ of the XY model, the Ising model in a transverse field and the mean-field model, all with $`L=15`$ spins and using $`S=20`$ samples. Because of the very high degeneracy we plotted the DOS for the mean-field model on a logarithmic scale. In Fig. 2 we show the relative error $`\delta Z/Z`$ based on Eq. (37) for the three models of various size, as obtained from the simulation (symbols). For these figures we used the imaginary-time algorithm, because then the statistical error can be related to $`e_A`$ directly (see Eq. (14) with $`A=\mathrm{exp}(\beta H)`$). The theoretical results (lines) for the error estimate, obtained by a direct exact numerical evaluation of (14) are shown too. We conclude that for all systems, lattice sizes and temperatures there is very good agreement between numerical experiment and theory. Results for the energy $`E`$ and specific heat $`C`$ are presented in Fig. 3 (XY model), 4 (Ising model in a transverse field), and 5 (mean-field model). The solid lines represent the exact result for the case shown. Simulation data as obtained from $`S=5`$ and $`S=20`$ samples are represented by symbols, the estimates of the statistical error by error bars. We see that the data are in excellent agreement with the exact results and equally important, the estimate for the error captures the deviation from the exact result very well. We also see that in general the error decreases with the system size. Both the imaginary and real-time method seem to work very well, yielding accurate results for the energy and specific heat of quantum spin systems with modest amounts of computational effort. ## X Conclusions The theoretical analysis presented in this paper gives a solid justification of the remarkable efficiency of the real-time equation-of-motion method for computing the distribution of all eigenvalues of very large matrices. The real-time method can be used whenever the more conventional, Lanczos-like, sparse-matrix techniques can be applied: Memory and CPU requirements for each iteration (time-step) are roughly the same (depending on the actual implementation) for both approaches. We do not recommend using the real-time method if one is interested in the smallest (or largest) eigenvalue only. Then the Lanczos method is computationally more efficient because it needs less iterations (time-steps) than the real-time approach. However if one needs information about all eigenvalues and direct diagonalization is not possible (because of memory/CPU-time) there is as yet no alternative to the real-time method. The matrices used in this example (up to $`32768\times 32768`$) are not representative in this respect: The real-time method has been used to compute the distribution of eigenvalues for matrices of dimension $`16777216\times 16777216`$ . Once the eigenvalue distribution is known the thermodynamic quantities directly follow. However if one is interested in the accurate determination of the temperature dependence of thermodynamic (and static correlation functions) properties but not in the eigenvalue distribution itself, the imaginary-time method is by far the most efficient method to compute these quantities. For instance the calculation of the thermodynamic properties for $`\beta J=0,\mathrm{},10`$ of a 15-site spin-1/2 system (i.e. implicitly solving the full $`32768\times 32768`$ eigenvalue problem) takes 1410 seconds per sample on a Mobile Pentium III 500 Mhz system. Finally we remark that although we used quantum-spin models to illustrate various aspects, there is nothing in the real or imaginary-time method that is specific to the models used. The only requirement for these methods to be useful in practice is that the matrix is sparse and (very) large. ## Acknowledgments Support from the Dutch “Stichting Nationale Computer Faciliteiten (NCF)” and the Dutch “Stichting voor Fundamenteel Onderzoek der Materie (FOM)” is gratefully acknowledged. ## A Expectation value calculation In this appendix we calculate the expectation value of the error squared, as defined in Section II. By definition we have $`E\left(\left\|TrRA\right\|^2\right)`$ $`=E\left(\left\|{\displaystyle \frac{1}{S}}{\displaystyle \underset{p=1}{\overset{S}{}}}{\displaystyle \underset{m,n=1}{\overset{D}{}}}\left(\delta _{m,n}Dc_{m,p}^{}c_{n,p}^{}\right)A_{m,n}\right\|^2\right)`$ (A1) $`={\displaystyle \frac{1}{S^2}}{\displaystyle \underset{p,p^{}=1}{\overset{S}{}}}{\displaystyle \underset{k,l,m,n=1}{\overset{D}{}}}(\delta _{k,l}\delta _{m,n}D\delta _{k,l}E(c_{m,p}^{}c_{n,p}^{})`$ (A2) $`D\delta _{m,n}E(c_{k,p^{}}^{}c_{l,p^{}}^{})+D^2E(c_{m,p}^{}c_{n,p}^{}c_{k,p^{}}^{}c_{l,p^{}}^{}))A^{}_{k,l}A_{m,n},`$ (A3) where $`p`$ and $`p^{}`$ label the realization of the random numbers $`c_{n,p}f_{n,p}+ig_{n,p}`$. First we assume that different realizations $`pp^{}`$ are independent implying that $`E(c_{m,p}^{}c_{n,p}c_{k,p^{}}c_{l,p^{}}^{})_{pp^{}}=E(c_{m,p}^{}c_{n,p})E(c_{k,p^{}}c_{l,p^{}}^{}).`$ (A4) Second we assume that the random numbers are drawn from a probability distribution that is an even function of each variable $`P(f_{1,p},g_{1,p},f_{2,p},g_{2,p},\mathrm{},f_{k,p},g_{k,p},\mathrm{},f_{D,p},g_{D,p})`$ (A5) $`=P(f_{1,p},g_{1,p},f_{2,p},g_{2,p},\mathrm{},f_{k,p},g_{k,p},\mathrm{},f_{D,p},g_{D,p})`$ (A6) $`=P(f_{1,p},g_{1,p},f_{2,p},g_{2,p},\mathrm{},f_{k,p},g_{k,p},\mathrm{},f_{D,p},g_{D,p}),`$ (A7) and symmetric under interchange of any two variables $`P(f_{1,p},g_{1,p},\mathrm{},f_{i,p},g_{i,p},\mathrm{},f_{j,p},g_{j,p},\mathrm{},f_{D,p},g_{D,p})`$ (A8) $`=P(f_{1,p},g_{1,p},\mathrm{},f_{j,p},g_{i,p},\mathrm{},f_{i,p},g_{j,p},\mathrm{},f_{D,p},g_{D,p})`$ (A9) $`=P(f_{1,p},g_{1,p},\mathrm{},g_{i,p},f_{i,p},\mathrm{},f_{j,p},g_{j,p},\mathrm{},f_{D,p},g_{D,p}),`$ (A10) for all $`i,j,k=1,\mathrm{},D`$. This is most easily done by drawing individual numbers from the same even probability distribution i.e. $`P(f_{1,p},g_{1,p},\mathrm{},f_{j,p},g_{i,p},\mathrm{},f_{i,p},g_{j,p},\mathrm{},f_{D,p},g_{D,p})`$ $`={\displaystyle \underset{n,m=1}{\overset{D}{}}}P(f_{n,p})P(g_{n,p}),`$ (A11) where $`P(x)=P(x)`$. Normalizing the vector $`(f_{1,p},g_{1,p},\mathrm{},f_{D,p},g_{D,p})`$ such that $`_{i=1}^D\|c_{n,p}\|^2=1`$ (for $`p=1,\mathrm{},S`$) does not affect the basic requirements (A5) and (A8). Making use of the above properties of $`P(f_1,g_1,\mathrm{},f_D,g_D)`$ we find that $`E(c_{m,p}^{}c_{n,p}^{})=\delta _{m,n}E(\|c_{m,p}\|^2)=\delta _{m,n}E(\|c\|^2),`$ (A12) where in the last equality we omitted the subscripts of $`c_{m,p}`$ to indicate that the expectation value does not depend on $`m`$ or $`p`$. An expectation value of a product of two $`c^{}`$’s and two $`c`$’s can be written as $`E(c_{m,p}^{}c_{n,p}^{}c_{k,p^{}}^{}c_{l,p^{}}^{})=`$ $`(1\delta _{p,p^{}})\delta _{m,n}\delta _{k,l}E(\|c_{m,p}\|^2)E(\|c_{m,p^{}}\|^2)`$ (A13) $`+\delta _{p,p^{}}\delta _{m,n}\delta _{k,l}(1\delta _{mk})E(c_m^{}c_m^{}c_k^{}c_k^{})`$ (A14) $`+\delta _{p,p^{}}\delta _{m,k}\delta _{n,l}(1\delta _{m,n})E(c_m^{}c_n^{}c_m^{}c_n^{})`$ (A15) $`+\delta _{p,p^{}}\delta _{m,l}\delta _{n,k}(1\delta _{m,n})E(c_m^{}c_n^{}c_n^{}c_m^{})`$ (A16) $`+\delta _{p,p^{}}\delta _{m,l}\delta _{n,k}\delta _{m,n}E(c_m^{}c_m^{}c_m^{}c_m^{})`$ (A17) $`=`$ $`(1\delta _{p,p^{}})\delta _{m,n}\delta _{k,l}E(\|c\|^2)^2`$ (A18) $`+\delta _{p,p^{}}\delta _{m,n}\delta _{k,l}(1\delta _{m,k})E(\|c_{m,p}\|^2\|c_{k,p}\|^2)`$ (A19) $`+\delta _{p,p^{}}\delta _{m,k}\delta _{n,l}(1\delta _{m,n})E(\|c_{m,p}\|^2\|c_{n,p}\|^2)`$ (A20) $`+\delta _{p,p^{}}\delta _{m,l}\delta _{n,k}(1\delta _{m,n})E(c_{m,p}^{}c_{n,p}^{}c_{n,p}^{}c_{m,p}^{})`$ (A21) $`+\delta _{p,p^{}}\delta _{m,l}\delta _{n,k}\delta _{m,n}E(\|c_{m,p}\|^4).`$ (A22) Furthermore for $`mn`$ we have $`E(c_{m,p}^{}c_{n,p}^{}c_{n,p}^{}c_{m,p}^{})=`$ $`E((f_{m,p}^22if_{m,p}g_{m,p}g_{m,p}^2)(f_{n,p}^2+2if_{n,p}g_{n,p}g_{n,p}^2))`$ (A23) $`=`$ $`E(f_{m,p}^2f_{n,p}^2)+2iE(f_{m,p}^2f_{n,p}g_{n,p})E(f_{m,p}^2g_{n,p}^2)`$ (A24) $`2iE(f_{m,p}g_{m,p}f_{n,p}^2)+4E(f_{m,p}f_{n,p}g_{m,p}g_{n,p})+2iE(f_{m,p}g_{m,p}g_{n,p}^2)`$ (A25) $`E(g_{m,p}^2f_{n,p}^2)2iE(g_{m,p}^2f_{n,p}g_{n,p})+E(g_{m,p}^2g_{n,p}^2)`$ (A26) $`=`$ $`E(f_{m,p}^2f_{n,p}^2)E(g_{m,p}^2f_{n,p}^2)E(f_{m,p}^2g_{n,p}^2)+E(g_{m,p}^2g_{n,p}^2)`$ (A27) $`=`$ $`0.`$ (A28) By symmetry $`E(\|c_{m,p}\|^2\|c_{n,p}\|^2)`$ does not depend on $`m`$, $`n`$ or $`p`$ and the same holds for $`E(\|c_{m,p}\|^4)`$. The fact that the vector $`(c_{1,p},\mathrm{},c_{D,p})`$ is normalized yields the identities $`E\left({\displaystyle \underset{n=1}{\overset{D}{}}}\|c_{n,p}\|^2\right)={\displaystyle \underset{n=1}{\overset{D}{}}}E(\|c_{n,p}\|^2)=DE(\|c\|^2)=E(1)=1,`$ (A29) and $`E\left(\left({\displaystyle \underset{n=1}{\overset{D}{}}}\|c_{n,p}\|^2\right)^2\right)`$ $`={\displaystyle \underset{m,n=1}{\overset{D}{}}}E(\|c_{n,p}\|^2\|c_{m,p}\|^2)`$ (A30) $`={\displaystyle \underset{n=1}{\overset{D}{}}}E(\|c_{n,p}\|^4)+{\displaystyle \underset{m,n=1}{\overset{D}{}}}(1\delta _{m,n})E(\|c_n\|^2\|c_m\|^2)`$ (A31) $`=DE(\|c\|^4)+D(D1)E(\|c\|^2\|c^{}\|^2)=E(1)=1,`$ (A32) where $`c`$ and $`c^{}`$ refer to two different complex random variables. Therefore we have $`E(\|c\|^2)`$ $`=1/D,`$ (A33) and $`E(\|c\|^2\|c^{}\|^2)`$ $`={\displaystyle \frac{1DE(\|c\|^4)}{D(D1)}}.`$ (A34) Substitution into (A13) yields $`E(c_{m,p}^{}c_{n,p}c_{k,p^{}}c_{l,p^{}}^{})=`$ $`(1\delta _{p,p^{}})\delta _{m,n}\delta _{k,l}D^2`$ (A35) $`+\delta _{p,p^{}}{\displaystyle \frac{1DE(\|c\|^4)}{D(D1)}}\left(\delta _{m,n}\delta _{k,l}(1\delta _{m,k})+\delta _{m,k}\delta _{n,l}(1\delta _{m,n})\right)`$ (A36) $`+\delta _{p,p^{}}\delta _{m,l}\delta _{n,k}\delta _{m,n}E(\|c\|^4).`$ (A37) and the final result for the variance reads $`E\left(\left\|TrRA\right\|^2\right)=`$ $`{\displaystyle \frac{1}{S}}({\displaystyle \frac{DD^2E(\|c\|^4)}{D1}}TrA^{}A+{\displaystyle \frac{1D^2E(\|c\|^4)}{D1}}\|TrA\|^2`$ (A38) $`+{\displaystyle \frac{(D+1)D^2E(\|c\|^4)2D}{D1}}{\displaystyle \underset{n=1}{\overset{D}{}}}\|A_{n,n}\|^2).`$ (A39) An expression for the fourth moment $`E(\|c\|^4)`$ cannot be derived from general properties of the probability distribution or normalization of random vector. We can only make progress by specifying the former explicitly. As an example we take a probability distribution such that for each realization $`p`$ the random numbers $`f_{n,p}`$ and $`g_{n,p}`$ are distributed uniformly over the surface of a $`2D`$-dimensional sphere of radius $`1`$. This probability distribution can be written as $`P(f_1,g_1,f_2,g_2,\mathrm{},f_D,g_D)\delta (f_1^2+g_1^2+f_2^2+g_2^2+\mathrm{}+f_D^2+g_D^21),`$ (A40) where we omitted the subscript $`p`$ because it is irrelevant for what follows. The even moments of $`\|c_n\|=(f_n^2+g_n^2)^{1/2}`$ are defined by $`E(\|c\|^{2M})={\displaystyle \frac{_{\mathrm{}}^{\mathrm{}}(f_1^2+g_1^2)^M\delta (f_1^2+g_1^2+f_2^2+g_2^2+\mathrm{}+f_D^2+g_D^21)𝑑f_1𝑑g_1𝑑f_2𝑑g_2\mathrm{}𝑑f_D𝑑g_D}{_{\mathrm{}}^{\mathrm{}}\delta (f_1^2+g_1^2+\mathrm{}+f_D^2+g_D^21)𝑑f_1𝑑g_1\mathrm{}𝑑f_D𝑑g_D}}.`$ (A41) It is expedient to introduce an auxiliary integration variable $`X`$ by $`E(\|c\|^{2M})={\displaystyle \frac{_{\mathrm{}}^{\mathrm{}}X^M\delta (f_1^2+g_1^2X)\delta (f_2^2+g_2^2+\mathrm{}+f_D^2+g_D^2(1X))𝑑X𝑑f_1𝑑g_1𝑑f_2𝑑g_2\mathrm{}𝑑f_D𝑑g_D}{_{\mathrm{}}^{\mathrm{}}\delta (f_1^2+g_1^2+\mathrm{}+f_D^2+g_D^21)𝑑f_1𝑑g_1\mathrm{}𝑑f_D𝑑g_D}}.`$ (A42) We can perform the integration over $`X`$ last and regard (A42) as the $`M`$-th moment of the variable $`X`$ with respect to the probability distribution $`P(X)={\displaystyle \frac{_{\mathrm{}}^{\mathrm{}}\delta (f_1^2+g_1^2X)\delta (f_2^2+g_2^2+\mathrm{}+f_D^2+g_D^2(1X))𝑑f_1𝑑g_1𝑑f_2𝑑g_2\mathrm{}𝑑f_D𝑑g_D}{_{\mathrm{}}^{\mathrm{}}\delta (f_1^2+g_1^2+\mathrm{}+f_D^2+g_D^21)𝑑f_1𝑑g_1\mathrm{}𝑑f_D𝑑g_D}}.`$ (A43) The calculation of $`P(X)`$ amounts to computing integrals of the form $`I_N(X)={\displaystyle _{\mathrm{}}^{\mathrm{}}}\delta \left({\displaystyle \underset{n=1}{\overset{N}{}}}x_n^2X\right)𝑑x_1𝑑x_2\mathrm{}𝑑x_N.`$ (A44) Changing to spherical coordinates we have $`I_N(X)`$ $`={\displaystyle \frac{2\pi ^{N/2}}{\mathrm{\Gamma }(N/2)}}{\displaystyle _0^{\mathrm{}}}r^{N1}\delta (r^2X)𝑑r`$ (A45) $`={\displaystyle \frac{\pi ^{N/2}}{\mathrm{\Gamma }(N/2)}}X^{N/21}\theta (X),`$ (A46) yielding $`P(X)`$ $`={\displaystyle \frac{I_2(X)I_{2D2}(1X)}{I_{2D}(1)}}`$ (A47) $`=(D1)(1X)^{D2}\theta (X)\theta (1X).`$ (A48) The moments $`E(\|c\|^{2M})`$ are given by $`E(\|c\|^{2M})`$ $`={\displaystyle _{\mathrm{}}^{\mathrm{}}}X^MP(X)𝑑X`$ (A49) $`=(D1){\displaystyle _0^1}X^M(1X)^{D2}𝑑X`$ (A50) $`={\displaystyle \frac{\mathrm{\Gamma }(D)\mathrm{\Gamma }(1+M)}{\mathrm{\Gamma }(D+M)}},`$ (A51) and the values of interest to us are $`E(\|c\|^0)=1,E(\|c\|^2)={\displaystyle \frac{1}{D}},E(\|c\|^4)={\displaystyle \frac{2}{D(D+1)}},`$ (A52) where the first two results provide some check on the procedure used. Substituting (A52) into (A38) yields $`E\left(\left\|TrRA\right\|^2\right)`$ $`={\displaystyle \frac{DTrA^{}A\|TrA\|^2}{S(D+1)}}.`$ (A53) ## B Error Bounds Here we prove that the difference between the eigenvalues of the Hermitian matrix $`A+B`$ and those obtained from the approximate time-evolution $`\mathrm{exp}(zA/2)\mathrm{exp}(zB)\mathrm{exp}(zA/2)`$ ($`z=i\tau ,\tau `$) is bounded by $`\tau ^2`$. In the following we assume $`A`$ and $`B`$ are Hermitian matrices and take $`\tau `$ a real, non-negative number. We start with the imaginary-time case. We define the difference $`R(\tau )`$ by $`R(\tau )`$ $`e^{\tau (A+B)}e^{\tau A/2}e^{\tau B}e^{\tau A/2}`$ (B1) $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }d\lambda {\displaystyle _0^\lambda }d\mu {\displaystyle _0^\mu }d\nu e^{\lambda A/2}e^{\lambda B}\{e^{\nu B}[2B,[A,B]]e^{\nu B}`$ (B2) $`+e^{\nu A/2}[A,[A,B]]e^{\nu A/2}\}e^{\lambda A/2}e^{(\tau \lambda )(A+B)},`$ (B3) a well-known result . We have $`R(\tau )`$ $`{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }𝑑\lambda {\displaystyle _0^\lambda }𝑑\mu {\displaystyle _0^\mu }𝑑\nu e^{\lambda A/2}e^{(\lambda \nu )B}[2B,[A,B]]e^{\nu B}e^{\lambda A/2}e^{(\tau \lambda )(A+B)}`$ (B4) $`+{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }𝑑\lambda {\displaystyle _0^\lambda }𝑑\mu {\displaystyle _0^\mu }𝑑\nu e^{\lambda A/2}e^{\lambda B}e^{\nu A/2}[A,[A,B]]e^{(\lambda \nu )A/2}e^{(\tau \lambda )(A+B)}`$ (B5) $``$ $`{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }𝑑\lambda {\displaystyle _0^\lambda }𝑑\mu {\displaystyle _0^\mu }𝑑\nu e^{\lambda A/2}e^{(\lambda \nu )B}[2B,[A,B]]e^{\nu B}e^{\lambda A/2}e^{(\tau \lambda )(A+B)}`$ (B6) $`+{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }𝑑\lambda {\displaystyle _0^\lambda }𝑑\mu {\displaystyle _0^\mu }𝑑\nu e^{\lambda A/2}e^{\lambda B}e^{\nu A/2}[A,[A,B]]e^{(\lambda \nu )A/2}e^{(\tau \lambda )(A+B)}`$ (B7) $`={\displaystyle \frac{1}{24}}\tau ^3e^{\tau (A+B)}\left([A,[A,B]]+[2B,[A,B]]\right),`$ (B8) and $`R(\tau )`$ $`{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }𝑑\lambda {\displaystyle _0^\lambda }𝑑\mu {\displaystyle _0^\mu }𝑑\nu e^{\lambda A/2}e^{(\lambda \nu )B}[2B,[A,B]]e^{\nu B}e^{\lambda A/2}e^{(\tau \lambda )(A+B)}`$ (B9) $`+{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }𝑑\lambda {\displaystyle _0^\lambda }𝑑\mu {\displaystyle _0^\mu }𝑑\nu e^{\lambda A/2}e^{\lambda B}e^{\nu A/2}[A,[A,B]]e^{(\lambda \nu )A/2}e^{(\tau \lambda )(A+B)}`$ (B10) $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }𝑑\lambda {\displaystyle _0^\lambda }𝑑\mu {\displaystyle _0^\mu }𝑑\nu e^{\lambda A/2}e^{(\lambda +\nu )B}[2B,[A,B]]e^{\nu B}e^{\lambda A/2}e^{(\tau +\lambda )(A+B)}`$ (B11) $`+{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }𝑑\lambda {\displaystyle _0^\lambda }𝑑\mu {\displaystyle _0^\mu }𝑑\nu e^{\lambda A/2}e^{\lambda B}e^{\nu A/2}[A,[A,B]]e^{(\lambda +\nu )A/2}e^{(\tau +\lambda )(A+B)}`$ (B12) $``$ $`{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }𝑑\lambda {\displaystyle _0^\lambda }𝑑\mu {\displaystyle _0^\mu }𝑑\nu e^{\lambda A/2}e^{(\lambda \nu )B}[2B,[A,B]]e^{\nu B}e^{\lambda A/2}e^{(\tau \lambda )(A+B)}`$ (B13) $`+{\displaystyle \frac{1}{4}}{\displaystyle _0^\tau }𝑑\lambda {\displaystyle _0^\lambda }𝑑\mu {\displaystyle _0^\mu }𝑑\nu e^{\lambda A/2}e^{\lambda B}e^{\nu A/2}[A,[A,B]]e^{(\lambda \nu )A/2}e^{(\tau \lambda )(A+B)}`$ (B14) $`=`$ $`{\displaystyle \frac{1}{24}}\tau ^3e^{\tau (A+B)}\left([A,[A,B]]+[2B,[A,B]]\right).`$ (B15) Hence the bound in $`R(\tau )`$ does not depend on the sign of $`\tau `$ so that we can write $`R(\tau )`$ $`s\|\tau \|^3e^{\|\tau \|(A+B)},`$ (B16) where $`s{\displaystyle \frac{1}{24}}[A,[A,B]]+[2B,[A,B]].`$ (B17) For real $`\tau `$ we have $`e^{\tau A/2}e^{\tau B}e^{\tau A/2}e^{\tau C(\tau )},`$ (B18) where $`C(\tau )`$ is Hermitian. Clearly we have $`e^{\tau (A+B)}e^{\tau C(\tau )}=R(\tau ).`$ (B19) We already have an upperbound on $`R(\tau )`$ and now want to use this knowledge to put an upperbound on the difference in eigenvalues of $`C(\tau )`$ and $`A+B`$. In general, for two Hermitian matrices $`U`$ and $`V`$ with eigenvalues $`\{u_n\}`$ and $`\{v_n\}`$ respectively, both sets sorted in non-decreasing order, we have $`\|u_nv_n\|UV,n.`$ (B20) Denoting the eigenvalues of $`A+B`$ and $`C(\tau )`$ by $`x_n(0)`$ and $`x_n(\tau )`$ respectively, combining Eq. (B16) and (B20) yields $`\|e^{\tau x_n(0)}e^{\tau x_n(\tau )}\|s\|\tau \|^3e^{\|\tau \|(A+B)}.`$ (B21) To find an upperbound on $`\|x_n(0)x_n(\tau )\|`$ we first assume that $`x_n(0)x_n(\tau )`$ and take $`\tau 0`$. It follows from (B9) that $`e^{\tau (x_n(\tau )x_n(0))}1s\tau ^3e^{\tau (A+B)\tau x_n(0)},`$ (B22) For $`x0`$, $`e^x1x`$ and we have $`x_n(0)A+BA+B`$. Hence we find $`x_n(\tau )x_n(0)s\tau ^2e^{2\tau (A+B)}.`$ (B23) An upperbound on the difference in the eigenvalues between $`C(\tau )`$ and $`A+B`$ can equally well be derived by considering the inverse of the exact and approximate time-evolution operator (B18). This is useful for the case $`x_n(0)>x_n(\tau )`$: Instead of using (B7) we start from $`\mathrm{exp}(\tau (A+B))\mathrm{exp}(\tau C(\tau ))=R(\tau )`$ ($`\tau 0`$). Note that the set of eigenvalues of a matrix and its inverse are the same and that the matrices we are considering here, i.e. matrix exponentials, are nonsingular. Making use of Eq. (B16) for $`R(\tau )`$ gives $`\|e^{\tau x_n(0)}e^{\tau x_n(\tau )}\|s\|\tau \|^3e^{\|\tau \|(A+B)},`$ (B24) and proceeding as before we find $`\tau (x_n(0)x_n(\tau ))`$ $`e^{\tau (x_n(0)x_n(\tau ))}1s\tau ^3e^{2\tau (A+B)}.`$ (B25) Putting the two cases together we finally have $`\|x_n(\tau )x_n(0)\|s\tau ^2e^{2\tau (A+B)}.`$ (B26) Clearly (B14) proves that the differences in the eigenvalues of $`A+B`$ and $`C(\tau )`$ vanish as $`\tau ^2`$. We now consider the case of the real-time algorithm ($`z=i\tau `$). For Hermitian matrices $`A`$ and $`B`$ the matrix exponentials are unitary matrices and hence their norm equals one. This simplifies the derivation of the upperbound on $`R(i\tau )`$. One finds $`R(i\tau )_E`$ $`s\|\tau \|^3,`$ (B27) where $`A_E^2TrA^{}A`$ denotes the Euclidean norm of the matrix $`A`$ . In general the eigenvalues of a unitary matrix are complex valued and therefore the strategy adopted above to use the bound on $`R(\tau )`$ to set a bound on the difference of the eigenvalues no longer works. Instead we invoke the Wielandt-Hoffman theorem : If $`U`$ and $`V`$ are normal matrices with eigenvalues $`u_i`$ and $`v_i`$ respectively, then there exists a suitable rearrangement (a permutation $`\varrho `$ of the numbers $`1,\mathrm{},n`$) of the eigenvalues so that $`{\displaystyle \underset{j=1}{\overset{N}{}}}\|u_jv_{\varrho (j)}\|^2UV_E^2.`$ (B28) Let $`U`$ and $`V`$ denote the exact and approximate real-time evolution operators respectively. The eigenvalues of $`A+B`$ and $`C(\tau )`$ are $`x_n(0)`$ and $`x_n(\tau )`$ respectively. All the $`x_n`$’s and $`x_n(\tau )`$’s are real numbers. According to the Wielandt-Hoffman theorem $`{\displaystyle \underset{j=1}{\overset{N}{}}}\|e^{i\tau x_j(0)}e^{i\tau y_j(\tau )}\|^2R(i\tau )_E^2s^2\tau ^6.`$ (B29) where $`y_j(\tau )=x_{\varrho (j)}(\tau )`$, $`\varrho `$ being the permutation such that inequality (B29) is satisfied. We see that Eq. (B29) only depends on $`(\tau x_j(0)mod2\pi )`$ and $`(\tau y_j(\tau )mod2\pi )`$, but so does the DOS (see Eq. (16)). Since the inequality (B17) and the DOS only depend on these “angles” modulo $`2\pi `$, there is no loss of generality if we make the choice $`0\|\tau (x_j(0)y_j(\tau ))\|\pi .`$ (B30) Rewriting the sum in (B17) we have $`{\displaystyle \underset{j=1}{\overset{N}{}}}\|e^{i\tau x_j(0)}e^{i\tau y_j(\tau )}\|^2`$ $`={\displaystyle \underset{j=1}{\overset{N}{}}}(22\mathrm{cos}(\tau (x_j(0)y_j(\tau ))))`$ (B31) $`=4{\displaystyle \underset{j=1}{\overset{N}{}}}\mathrm{sin}^2(\tau /2(x_j(0)y_j(\tau ))).`$ (B32) As we have $`\mathrm{sin}^2x{\displaystyle \frac{4x^2}{\pi ^2}},\text{ for }0\|x\|\pi /2,`$ (B33) the restriction Eq. (B30) allows us to write $`{\displaystyle \underset{j=1}{\overset{N}{}}}(x_j(0)y_j(\tau ))^2{\displaystyle \frac{\pi ^2s^2}{4}}\tau ^4,`$ (B34) implying $`\|x_j(0)y_j(\tau )\|{\displaystyle \frac{\pi s}{2}}\tau ^2.`$ (B35) In summary, we have shown that in the real-time case there exists a permutation of the approximate eigenvalues such that the difference with the exact ones vanishes as $`\tau ^2`$. Finally we note that both upperbounds (B35) and (B26) hold for arbitrary Hermitian matrices $`A`$ and $`B`$ and are therefore rather weak. Except for the fact that they provide a sound theoretical justification for the real and imaginary-time method, they are of little practical value.
warning/0004/astro-ph0004219.html	ar5iv	text	# Structure and Star Formation in NGC 9251footnote 11footnote 1to appear in the August 2000 Astronomical Journal ## 1 Introduction This paper is part of an ongoing series of papers studying the general H I and optical properties of late-type barred spiral galaxies (Sbc-Sd, hereafter LTBS). We will examine the stellar and star forming properties of NGC 925 based on observations from the WIYN 3.5m telescope<sup>2</sup><sup>2</sup>2The WIYN observatory is a joint facility of the University of Wisconsin-Madison, Indiana University, Yale University, and the National Optical Astronomy Observatories.. We previously observed NGC 925 and NGC 1744, both LTBS, in H I as part of a study of the gaseous properties of this class of galaxies (Pisano et al., 1998, hereafter Paper I), and to determine the pattern speed of their bars (Elmegreen et al. 1998). NGC 925 was chosen for study because it is a nearby ($``$9.3 Mpc, Silbermann et al. 1996), apparently prototypical LTBS which is well oriented in inclination for dynamical and morphological study. This value for the distance will be used throughout this paper. In paper I, we found that NGC 925 had a weak spiral pattern and bar as indicated by the small streaming motions in the spiral arms and bar. Asymmetries are quite prevalent in NGC 925. NGC 925’s southern spiral arm has reasonably coherent streaming motions, while the northern arms do not. The center of the bar is also slightly offset from the dynamical center of the galaxy by $``$950 pc. These asymmetries could possibly be related to the presence of a small H I cloud (M$`{}_{HI}{}^{}`$10<sup>7</sup>M) interacting with the main galaxy. The strength of the interaction, however, is most likely not enough to drive the asymmetries in NGC 925. The presence of an off-center bar, and a single coherent spiral arm are typically viewed as characteristic properties of barred Magellanic (SBm) galaxies (de Vaucouleurs & Freeman 1972), and have been shown to be potentially long-lasting (Levine & Sparke 1998, Noordermeer et al. 2000) While NGC 925 does not illustrate these properties as dramatically as the LMC, for example, it certainly indicates that these characteristic asymmetries do not suddenly appear in SBm’s, but are rather a continuous change in morphology throughout late-type disk galaxies. A summary of basic properties of NGC 925 is given in table 1. LTBS are an understudied and, hence, a poorly understood class of galaxies. While bars may play an important role in galaxy evolution by generating gas inflow (e.g., Matsuda & Nelson 1977, Athanassoula 1992), driving spiral structure (e.g., Sanders & Huntley 1976; James & Sellwood 1978), and flattening the radial abundance gradient (Martin & Roy 1994), these effects are shown to exist primarily in early-type barred spirals (S0-Sb). LTBS do not appear to have strong enough bars to cause such effects (see paper I). Previous studies of LTBS have found that late-type bars tend to have exponential light profiles, as opposed to flat profiles, and end well within corotation, as opposed to just inside it (Elmegreen & Elmegreen 1985, Elmegreen et al. 1998). Furthermore, late-type bars are not likely to be driving the spiral pattern (Sellwood & Sparke 1988), in contrast to early-type bars (Elmegreen & Elmegreen 1989). Star formation properties are also different, with late-type bars tending to be gas-rich and have star formation throughout the bar, while early-type bars have star formation only near the ends of the bar (Phillips 1995). Our observations in paper I of NGC 925 provided no evidence to contradict any of these properties. In this paper, we will investigate the optical properties of NGC 925, with a specific eye towards the nature of star formation in the galaxy. We will see if there are other signatures of asymmetry in NGC 925 that are typically associated with SBm’s. Finally, we will try to devine what role the bar plays in determining the properties of the galaxy as a whole. In section 2 of this paper we discuss the observations and reduction. In section 3, we discuss the optical properties of galaxy as a whole. Section 4 discusses the star formation properties of NGC 925 both across the entire galaxy, and in the bar and spiral arms specifically. We conclude in section 5. ## 2 Observations and Reductions We observed NGC 925 with the WIYN 3.5m telescope during the nights of October 6 & 7, 1996. The conditions were superb; all three nights had subarcsecond seeing and were photometric. We took the images through the Harris B, V, and R broadband filters and on and off band H$`\alpha `$ filters using the WIYN imager, a thinned 2048<sup>2</sup> STIS CCD with 21$`\mu `$m pixels. The pixels correspond to $`0.196^{\mathrm{}}`$ on the sky. The gain was 2.8 e<sup>-</sup> per ADU, and the read noise was 8 e<sup>-</sup>. The imager has a total field of view of 6.5. Since NGC 925 has a major axis of approximately 10, we had to mosaic two images to get full coverage of the galaxy. We observed the east and west sides of the galaxy separately with overlap in the bar region. The parameters of the observations are listed in table 2. We reduced the data in the usual manner. We took bias images during the afternoon before each night. We visually inspected the biases, and statistically compared them and found them to be identical for all three nights. This being the case, we combined all 44 bias frames into one image for bias subtraction to maximize the signal to noise. Each afternoon, we took $``$5 well exposed (10000-35000 counts) dome-flats per filter. We discarded any flats which were significantly deviant from the average. We averaged the remaining flats together to correct for the response of the CCD. We switched filter wheels during each night’s observations which sometimes caused the dust features to change in location and/or intensity. As some of the flats changed between nights, we only used the flats from the same night to flatten the program images. Aside from the off-line H$`\alpha `$ images, all images were flat to better than 5%. The off-line H$`\alpha `$ images were flat to 10%. The poor quality of the flat-fielding for the H$`\alpha `$-off filter may require sky-flats to improve. To flux calibrate our images we observed broadband standards in the PG0231+051 and SA92 fields from Landolt (1992), and G191B2B, a spectrophotometric standard from Massey et al. (1988). The photometric solutions for the broadband observations were found using the apphot package in IRAF. They are as follows: $$V=OB_V+1.08+0.17\times X0.04\times (BV)$$ (1) $$B=(OB_{(BV)}+OB_V)+1.28+0.22\times X0.04\times (BV)$$ (2) $$R=(OB_VOB_{(VR)})+0.95+0.10\times X0.04\times (VR)$$ (3) where OB is the instrumental magnitude and X is the airmass. We then applied these solutions to the bias subtracted, flat-fielded, sky subtracted images using IRAF to produce surface brightness maps. The errors on these equations are roughly 0.05 mag. As this is comparable to our flat-fielding errors, we combine the errors to obtain our total error on the broadband magnitudes, which is approximately 0.07 mag. Sky subtraction was difficult given the small amount of sky present in our images. There was typically only a small corner of the CCD frame which appeared to contain mostly sky emission. For each side of the galaxy we used imstat to determine a mean sky level, and subtracted it from the image before mosaicing. This value is quite uncertain given the lack of sky present in the images. Because we were not able to reliably determine the sky values in our images, we did not attempt to determine total magnitudes for NGC 925. The photometric solutions for the line and continuum images were calculated accounting for the airmass of the observations and the response of the line and continuum filters as determined for our observations of the spectrophotometric standard. Using the respective exposure times and the fractional contribution of the two \[NII\] lines to the H$`\alpha `$ image, we then can get the continuum subtracted H$`\alpha `$ image. Based on the spectra in Martin & Roy (1994), we assumed the H$`\alpha `$ line accounted for 81% of the emission in the line filter, with the two \[NII\] lines accounting for 14% and 5% of the emission. We applied these solutions to each filter and each side of the galaxy separately. The line and continuum images were aligned to better than 1 pixel with similar, but not exactly the same, seeing conditions. Once each image was fully calibrated, we aligned and averaged together the east and west images of NGC 925 to form a total image of the galaxy. To align the images, we identified stars present in both images, calculated the average shift needed to bring the two sides into alignment, and shifted the images using imshift. Images were aligned to better than 1 pixel in each dimension. A similar procedure was followed to align the full galaxy B, V, R, and H$`\alpha `$ images to each other. At this point we began our analysis. Figure 1 shows our combined B, V, R image of NGC 925. ## 3 Global Photometric Properties One of the key thrusts of investigations of late-type spirals is the degree to which their disks are lopsided. Extreme late-types such as the Magellanic spirals have dynamical centers that are significantly offset from the center of the outer isophotes (de Vaucouleurs & Freeman 1972, Odewahn 1996). We can measure a photometric center from our mosaiced image of NGC 925 and a dynamical center from the HI velocity field (Paper 1). NGC 925 is also a barred galaxy, so we can compare the photometric and dynamical center with the apparent center of the bar. Magellanic spirals in particular host bars that are greatly offset from the photometric and dynamical centers. We chose to use the isophote fitting routine in the IRAF package, ellipse. We first removed the stars and saturated columns from the image by replacing the central regions of bright stars within a radius of 2.4<sup>′′</sup> with the mean from an annulus with inner radius of 3<sup>′′</sup> and width 1<sup>′′</sup>. While some stars were still visible in the image the residuals were only slightly higher than the background level of the galaxy. We fit a series of elliptical annuli 0.<sup>′′</sup>4 in width to the galaxy. The semi-major axis of each annulus was 4<sup>′′</sup> larger than the previous one. The task ceased to produce fits which converged at a radius of 5, near the edge of our field of view. We present the results of our elliptical fits in Figures 2 & 3. Note that the scale on the surface brightness plot comes from our photometric calibrations of the background sky level (although this may not be the true sky). A number of features stand out in the isophotal fits in Figure 2. Overall, the results of the fits are very similar in all bands even though the center, position angle, and ellipticity of the galaxy all vary widely as a function of radius. Within a radius of about 60<sup>′′</sup> (i.e. the bar region) the center is fit consistently in all bands. Simply taking the means from each band we define the bar center to be: 2<sup>h</sup>24<sup>m</sup>17.0$`{}_{}{}^{s}\pm `$0.2<sup>s</sup>, 33<sup>o</sup>2115$`{}_{}{}^{\prime \prime }\pm `$1<sup>′′</sup>. The epoch for all of the coordinates in this paper is 1950, unless otherwise stated. There is a sharp change in the fitted declination of the center just beyond a radius of 60<sup>′′</sup>, coincident with the position of the prominent spiral arm on the southern side of the galaxy. The width of the “bump” in the fitted declination of the center is consistent with the width of the arm. The fitted declination then returns to a value comparable with that derived from the innermost parts of the galaxy. So if we were to ignore the effect of the spiral arm we find that the fitted photometric center of NGC 925 is only slowly varying with radius out to $`170^{\prime \prime }`$ which is close to edge of the optical disk. Beyond about 200<sup>′′</sup> the center, position angle, and ellipticity remain relatively constant. While the three bands return inconsistent fits for the center of the galaxy, we take the mean center derived from fits to the annulus 210<sup>′′</sup>-270<sup>′′</sup> to be the center of the outer isophotes. This yields an isophotal center at 2<sup>h</sup>24<sup>m</sup>14.8$`{}_{}{}^{s}\pm `$0.6<sup>s</sup>, 33<sup>o</sup>2123$`{}_{}{}^{\prime \prime }\pm `$2.5<sup>′′</sup> which is different from the bar center. Studies of lopsided galaxies usually use the center of the outermost isophotes to define a photometric center of the galaxy (e.g. Odewahn 1991). NGC 925 shows that this center will be a function of the band with which one chooses to carry out the photometry. The fitted position angle of NGC 925 in Figure 2c slowly increases from -80<sup>o</sup> to -55<sup>o</sup> from the center to a radius of 180<sup>′′</sup>; beyond that it drops abruptly. Similar trends are seen in Figure 2d where we plot the fitted ellipticity as a function of radius. It is slowly decreasing from the center out to 180<sup>′′</sup> where it also drops abruptly before leveling off. We list our adopted values for the center, position angle, and ellipticity in Table 3 In figure 3a we plot the azimuthally averaged surface brightness distribution for NGC 925. We fit a double exponential to this distribution and we show the residuals in Figures 3b and 3c, respectively. Residuals from our fit solely to the exponential outer-disk (Fig 3b) clearly illustrate the presence of a second exponential at radii within 60$`\mathrm{}`$, and a systematic surface brightness enhancement between 140$`\mathrm{}`$ and 180$`\mathrm{}`$. The latter is much clearer after we subtract our exponential fit to the bar and re-scale the residuals in the lower panel. We chose to fit an exponential to the bar, both because it represents the data well, and because Elmegreen & Elmegreen (1985) found that late-type galaxies including NGC 925 have exponential bars. Both of our fits are shown in the top panel. The derived scale-lengths and central brightnesses for the bar and disk are listed in table 3. One very interesting property illustrated by these scale-lengths is that the disk tends to get redder with increasing radius, as the B scale-length is much smaller than the R scale-length. This is contrary to what is found for most field galaxies (Vennik et al. 1996), and could be indicative of a differing star formation history for NGC 925 as compared to most other galaxies. This could also indicate the presence of a very young inner disk and bar in NGC 925. By looking at all the results of our isophotal fitting, we can pick out some distinct features in NGC 925 which are evident in both a contour plot of the model (figure 4), and in the graphs of the fit parameters (figures 2 & 3). The inner bar region, within 30$`\mathrm{}`$, is characterized by widely varying values for the center, position angle, and ellipticity that are most likely due to patchy star formation and dust in the center of the bar. We get consistent results for the rest of the bar. The southern arm strongly affects the fits in the remainder of the disk. This is evident from the drastic change in center, position angle, and ellipticity at 180$`\mathrm{}`$ where the spiral arm fades (as shown in the surface brightness fits). The slow increase in position angle with increasing radius illustrates the change from the bar dominated region of the galaxy to the region where the spiral arm dominates. Other variations in the fits and surface brightness residuals can be explained by the small scale variations in brightness caused by H II regions and patchy dust within NGC 925. The bar of NGC 925, based on our fits shown in figure 3, extends out to approximately 60$`\mathrm{}`$ (2.7 kpc). It is at this point where the exponential profile of the bar begins to meet the exponential disk. The bar is centered at 2<sup>h</sup>24<sup>m</sup>17<sup>s</sup> in right ascension, and 33<sup>o</sup>21$`\mathrm{}`$15$`\mathrm{}`$ in declination, which is not coincident with the center of the outer isophotes. This is the mean from the isophotal fits of the center within the inner 60$`\mathrm{}`$. The bar has an ellipticity of $``$0.69 (b/a = 0.3), and, with a position angle of -75<sup>o</sup>, it is aligned within a few degrees of the major axis of the galaxy. The alignment of the bar position angle with the galaxy position angle, and the origination of the spiral arms from the end of the bar may indicate that the bar has had some effect on the disk of the galaxy. We now have four centers derived for NGC 925: the peak of the brightness distribution (from the RC3), the dynamical center (from Paper I), the center of the bar, and the center of the “outer” isophotes (from above). In figure 5, we compare the locations of these centers as a measure of the asymmetry in the galaxy. The size of the boxes are representative of the fitting errors; the ellipse is the beam size from the H I data used to model the rotation curve. We find that the center of the brightness distribution is nearly coincident with the center of the bar, while the dynamical and outer isophote centers are coincident with each other. The bar center is NOT coincident, however, with either the dynamical or outer isophotal center. This is indicative of an asymmetric galaxy and is quite similar to what is seen in SBm’s such as NGC 4618 (Odewahn 1991) which has a 662 pc displacement between the bar center and the center of the outer isophotes. For NGC 925 the dynamical and outer isophotal centers are offset parallel to the bar, whereas most SBm’s show bars offset perpendicular to their major axis (Odewahn 1991). NGC 925 may be exhibiting a different type of phenomenon, but this property of SBm’s requires further study. While we were unable to calculate total magnitudes for NGC 925, we used published values for M<sub>B</sub> to examine mass-to-light ratios for NGC 925. From the RC3, we find that NGC 925 has an M<sub>B</sub>=-19.9 mag. Using the total and H I masses we derived in Paper I, we calculated M<sub>tot</sub>/L<sub>B</sub> to be 5.0, and M$`_{HI}`$/L<sub>B</sub> to be 0.36, consistent with what we expect from its classification as an Scd spiral (Roberts & Haynes 1994), but also similar to values found for later-type spiral galaxies. ## 4 Star Formation in NGC 925 ### 4.1 Global Properties In a seminal paper, Kennicutt (1989) observationally examined the dependence of star formation on the properties of the gas in the galaxy. For his study he used the expression for the critical density for a instability to grow in a thin isothermal gas disk: $$\mathrm{\Sigma }_{crit}=\alpha \frac{\kappa c}{3.36G}$$ (4) where c is the velocity dispersion of the gas, G is the gravitational constant, $`\alpha `$ is a dimensionless constant near unity which accounts for the deviation of the disk from a thin, isothermal gas disk, and $`\kappa `$ is the epicyclic frequency given by: $$\kappa =1.41\frac{V}{R}(1+\frac{V}{R}\frac{dV}{dR})^{\frac{1}{2}}$$ (5) where V is the rotation velocity at a radius R. In our case the rotation curve is given by a Brandt rotation curve derived from a fit to the H I data in paper I: $$V(R)=\frac{V_{max}}{R_{max}}\frac{R}{[\frac{1}{3}+\frac{2}{3}(\frac{R}{R_{max}})^n]^{\frac{3}{2n}}}$$ (6) where V<sub>max</sub>=118 km s<sup>-1</sup>, R<sub>max</sub>=17.8 kpc, and n=1.46 (paper I). Kennicutt examined how star formation compared to the gas density and the expected critical density within a collection of galaxies. He found that star formation, as traced by H$`\alpha `$ emission, was correlated with the gas surface density (as traced by H I emission) and that star formation was suppressed when the gas surface density fell below a critical value. Based on data from paper I, and using our H$`\alpha `$ images, we carried out a similar analysis on NGC 925. From the H I data we derive an average velocity dispersion, c, of 10 km s<sup>-1</sup> for NGC 925. The value of $`\alpha `$ should be close to one, because galaxies should be close to being thin, isothermal disks, although Kennicutt (1989) finds an $`\alpha `$ of 0.67 fits his data better. In figure 6 we show the critical surface density ($`\mathrm{\Sigma }_{crit}`$) for three values of $`\alpha `$ and surface density of H I multiplied by 1.47, $`\mathrm{\Sigma }_{gas}`$, (following the example of Kennicutt 1989 to account for the molecular component of the gas). We also plot the H$`\alpha `$ emission normalized to its peak. The H$`\alpha `$ and H I surface densities shownare azimuthal averages. We find the H$`\alpha `$ surface brightness to be relatively constant throughout the galaxy before cutting off abruptly at about 15 kpc, whereas the gas surface densitydecreases slowly with radius out to 18 kpc. This suggests that star formation in NGC 925 is not directly correlated with the gas surface density on a global scale. While the abrupt decline of H$`\alpha `$ surface brightness occurs at the edge of our field of view, we can use the value of $`\mathrm{\Sigma }_{gas}`$/$`\mathrm{\Sigma }_{crit}`$ at this point to yield $`\alpha `$. Looking figure 6, we see that a value of $`\alpha `$ of near 0.3 works best for NGC 925 as opposed to $`\alpha `$=1 and $`\alpha `$=0.67. In reality, the H$`\alpha `$ emission may continue beyond the edge of our field of view, implying a lower value of $`\alpha `$ still. Regardless of the extent of H$`\alpha `$ emission beyond the region we imaged, the star formation criterion found by Kennicutt (1989), $`\alpha `$=0.67, does not hold for NGC 925; instead, a value of $`\alpha `$ closer to 0.3 is more appropriate. This has been found to be the case for other galaxies as well. Values of $`\alpha `$ around 0.3 have been found for dIrr galaxiesby Hunter, Elmegreen, & Baker (1998), and for low surface brightness dwarf galaxies by van Zee et al. (1997). This is evidence that star formation in NGC 925 is governed more by local conditions than by global instabilities; a property more typical of later type galaxies than NGC 925. This is not surprising as star formation can occur in any type galaxy even when such a criterion is not met (see Ferguson et al. 1998 for a discussion). Ferguson et al. find H II regions outside the cited limits to the optical disk in a sample of spiral galaxies. Their outer H II regions, and ours, lie on coherent spiral arms even at large radii. A better analysis of star formation criteria in NGC 925 could be done using CO data to trace the molecular component of the gas and using a larger field-of-view CCD to better trace the full extent of the H$`\alpha `$ emission in NGC 925. ### 4.2 Star Formation in the spiral arms of NGC 925 Comparing the H I distribution of Paper I with the H$`\alpha `$ and broadband images of NGC 925 allows us to examine the nature of star formation in the galaxy. The star formation in NGC 925, as traced by the H$`\alpha `$ emission, is found in two main regions. The largest star forming regions are located at the end of the southern spiral arm and in the bar. In addition, star formation occurs at a lower level throughout the southern spiral arm (although it may be partially obscured by dust). The northern arm has sporadic star formation near the bar, but very little elsewhere, except, perhaps, near the end of the arm on the east side of the galaxy. The lack of coherent star formation along the north arms is not surprising given the weak definition of the stellar and gaseous arms. Figure 7 shows the H$`\alpha `$ in red on the H I in blue. In general, the star formation is occurring on top of H I peaks and the H$`\alpha `$ emission traces the H I spiral arms very closely. The coincidence between the two distributions is highly correlated across the entire galaxy, but the exceptions are quite interesting. Specifically, there is a sizable H II region on the northwest spiral arm near the bar which is sitting in a local H I depression. This is the only major H II region which is not associated with a H I peak. This could be due to the star formation depleting the H I in this region, ionizing it, or the H I hole could be filled with molecular gas. Looking at the broadband BVR image with the H I contours (figure 8), we see that the peak H I does trace the optical emission (as well as the H$`\alpha `$ emission), as expected. This figure clearly shows that the southern spiral arm is much stronger and more coherent than the northern arm in H I and optical emission. The north side of the galaxy shows more flocculent spiral structure and less intense star formation as compared to the south. The southern arm also has a prominent gap in optical emission in the middle of the arm, while it still has a large amount of H I. This may be due to a large dust lane crossing the arm at this position. These are not caused by inclination, as the southern part of the galaxy is the far side of the galaxy. You can also see the bluest regions of the galaxy are those where large H II regions exist, specifically at the end of the southern arm, and to the northeast of the bar (which is the end of one of the northern arms). As our data is of sufficient resolution and photometric quality to address the question of triggering of star formation in spiral arms of late-type galaxies, we took cuts through various regions of the spiral arms and bar in NGC 925 (figure 10) to look for offsets between the peaks of emission in B, R, H$`\alpha `$, and H I (see Beckman & Cepa 1990). Figure 9 shows the location of these cuts. These cuts were 4$`\mathrm{}`$ wide and of varying length (as shown in figure 9). They are roughly centered on the arms or bar. The emission is an average over the width of the cut and is normalized to the peak intensity of that band in the cuts. Triggering is defined as an increase in the star formation efficiency in the arm versus the inter-arm region; it is not simply an increase in total star formation. An offset in the peaks would suggest there is an age gradient across the cut (in a dust free world). Beckman & Cepa (1990), for instance, see such a color gradient in NGC 7479 between the B and I bands. Dust, however, can produce a similar gradient due simply to an extinction gradient within the spiral arm (Elmegreen 1995). In addition, we can not determine if triggering is actually occurring in the arm solely from an age (or color) gradient; it is entirely possible to see a gradient without any triggering (cf. Elmegreen 1995). What we can say is that the absence of an offset between the peaks in B and R implies the lack of an age gradient in the underlying stellar population. Cuts along the southern arm (figure 10a-d) show the optical emission is well-aligned with the H I across the arm for most of the length of the arm. Peaks in B and R are well-aligned. In general, the H$`\alpha `$ peaks correlate with the broadband peaks, but not always with the H I peaks. There are regions with no H$`\alpha `$ peaks, but with strong broadband emission. In addition, there are two noticeable instances of H$`\alpha `$ peaks in H I depressions. These occur at $`\pm `$ 40 in cut c and in cut d at 40. Star formation in the south arm is located in large H II regions which tend to lie on the peak of the H I distribution. Cuts across the north arms (figure 10e-g) of NGC 925 show a decidedly different phenomenon than those of the south arm. Cut e shows optical emission peaking in an H I depression at -50$`\mathrm{}`$. There is also some H$`\alpha `$ emission peaking with the H I at 20$`\mathrm{}`$. Other cuts across the north arms, such as cut f, show a correlation between the optical and gas peaks, while cut g shows an offset between the optical emission and the H I peaks. In the north arms, the star formation appears less prolific than in the southern arm and occurs in smaller H II regions which are less coherently placed than in the south of NGC 925. This suggests that the nature of the spiral pattern, and hence the star formation, in the north and south arms of NGC 925 are different, if not fundamentally, then in our perception of it. This type of asymmetry, with one spiral arm being dominant in a galaxy, is quite typical in later-type spiral galaxies, such as SBm’s (Odewahn 1991, de Vaucouleurs & Freeman 1972). NGC 925 illustrates that these properties are not limited to SBm’s, but are evident in earlier type spiral galaxies, although less prominently. ### 4.3 Star Formation in the bar of NGC 925 In general, LTBS tend to have gas-rich bars with star formation occurring throughout the bar (Phillips 1995). Friedli & Benz (1995) suggest that star formation occurs along the major axis of dynamically young bars, and is concentrated in the center or in rings of older bars. NGC 925’s bar is quite gas-rich and has star formation occurring all along its major axis. Thus, the Friedli & Benz (1995) model would imply that NGC 925 has a young bar. The extremely blue center of NGC 925, compared with its outer region (see section 3), could indicate a dynamically young bar as well, one which has just recently started prolific star formation. On the other hand, it is possible that weak bars, such as NGC 925’s, are simply inefficient at funneling gas into the galactic nucleus. Furthermore, because corotation is far out in the disk (Elmegreen et al. 1998), the spiral arms have a large reservoir of gas to drive towards the center of NGC 925. As there is no inner Lindblad resonance in NGC 925 (Elmegreen et al. 1998) where the gas would pile up, the bar would be well-supplied with gas for long-term star formation, and it may not be young at all. The star formation in the bar appears to be offset somewhat from the H I peaks. Examining the bar region in figure 7, we find that there is a great deal of star formation on the north side of the bar that is systematically displaced away from the H I peaks. This is better illustrated by cuts h-k in figure 10. Furthermore, the minor axis cuts (i-k) show that the B and R peaks are systematically offset to the north of the H I peak. While there is a H$`\alpha `$ peak aligned with the H I peak in cut j, the main optical emission is offset. This could be explained by dust extinction, depletion of the H I by star formation or ionization, the gas being in molecular form, or the bar being a wave phenomenon. For this offset to be an extinction feature, some process must pile up the dust preferentially along one side of the bar. Dust lanes have been found in the bars of other galaxies, and their coherence is usually attributed to shocks occurring in these bars (see Athanassoula 1992 for a nice discussion). The bar of NGC 925, however, does not have strong streaming motions (paper I) which would be indicative of shocks. The dust appears to be in filaments in the bar region, and not in well defined lanes. Furthermore, some of the cuts perpendicular to the major axis of the bar do show star formation correlated with the H I peaks, but all of the cuts indicate that the stellar emission is downstream of the H I peaks. Extinction should affect both H$`\alpha `$ emission and R band emission to a similar extent. As this is not evident in the cuts, it is improbable that dust is causing the apparent offset. The major axis cut (h) shows the H$`\alpha `$ peaking throughout the bar, with the gaps in emission likely due to the prolific dust content of the bar. However, the broadband optical peaks are concentrated in a H I depression. This may be due to the recent star formation and young stars evacuating or ionizing the H I in this region of the bar. The H I depression in the bar could be filled with molecular gas, and not truly be depleted of all gas. Along the major axis of the bar the H I seems to be piled up at the bar ends, particularly near the dynamical center of the galaxy at 50$`\mathrm{}`$. Spectra of the bar taken with the DensePak instrument on WIYN are in hand and should yield insight into the nature of the ionized gas in the bar. Higher spatial resolution H I data will also help to untangle the dynamics of gas in the bar. If the offset of H I south of the H$`\alpha `$ in the bar is not caused by dust, or by young stars ionizing or otherwise clearing out the H I in the bar, or by the gas being in molecular form, then the bar could be a wave phenomenon. This may be linked to the offset dynamical center of the galaxy as reported in Paper I. The dynamical center is located on the west end of the bar, so the geometry and dynamics make sense for this scenario, with the star formation occurring downstream of the gaseous peaks and the stars even further downstream, provided that star formation is occurring as the gas moves through the wave ## 5 Discussion & Conclusions We have observed NGC 925 using the WIYN telescope in B, V, R, and H$`\alpha `$ filters to better understand the stellar distribution and star formation properties of the galaxy. We have compared these observations with previously described H I observations (Paper I) to get a nearly complete picture of the properties of NGC 925. The global properties of NGC 925 are typical for a late-type spiral galaxy. It has an obvious bar, and two large spiral arms. The southern spiral arm is apparently much stronger than the northern arm, which is quite flocculent. The relative dominance of the southern arm is apparent both optically and in the H I (paper I). This difference between the northern and southern spiral arms is not the only asymmetry present in the galaxy. The center of the galaxy, as derived from the outer isophotes, is coincident with the dynamical center (Paper I), but not with the center of the bar. These properties are typically associated with barred Magellanic spirals (de Vaucouleurs & Freeman 1972), although in the case of SBm’s the offset bar and single, dominant spiral arm tend to be much more pronounced than they are in the case of NGC 925. Nevertheless, we see that such properties are not limited to Magellanic spirals, but are present in more subtle ways in earlier-type spirals. The optical properties of NGC 925, such as mass-to-light ratio and absolute magnitude, are as would be expected for a typical late-type galaxy. From our isophote fitting, we find that NGC 925’s surface brightness distribution can be characterized by a double exponential; one for the bar and one for the outer disk. The bar is quite easily distinguished from the rest of the galaxy in the fits, by having a well-defined center, position angle, and ellipticity differing from the outer disk. The bar also has a exponential brightness distribution with a scale-length much smaller than the rest of the galaxy. There is little evidence that the bar has any effect on the structure of the rest of the galaxy. The isophote fits show quite a bit of structure outside of the bar region. This is almost certainly due to the bright southern spiral arm, which shows up clearly as a surface brightness enhancement in the fits. At the same radius as this enhancement, we can see distinct changes in the fits for the center, ellipticity, and position angle. The fits do not seem to be affected by the weaker northern arm. This is further evidence for NGC 925’s similarity to the SBm one-armed spiral pattern. The fits also show that NGC 925 gets redder with increasing radius, atypical of normal spiral galaxies in the field. This could indicate a recent enhancement of star formation in the inner galaxy. From our analysis of the critical density for massive star formation in NGC 925, using the technique of Kennicutt (1989), we find that the galaxy has widespread star formation occurring outside of the radius where it should be suppressed. NGC 925 has a very flat H$`\alpha `$ surface brightness distribution, with emission extending all the way to the edge of our image at 15 kpc. Kennicutt (1989) using a sample of 15 Sc galaxies found that star formation was suppressed when $`\alpha `$=$`\mathrm{\Sigma }_{gas}`$/$`\mathrm{\Sigma }_{crit}`$ fell below 0.67. For NGC 925 star formation continues past this radius and is more consistent with an $`\alpha `$0.3, like what was found by Hunter et al. (1998) and van Zee et al. (1997) for later-type galaxies. The exact value of $`\alpha `$ is unclear, but probably lower, as the H$`\alpha `$ emission could continue beyond our field of view. As mentioned above, the spiral arms in NGC 925 are quite different. The north arm is quite weak, lacks coherence, and has only patchy star formation. The southern arm, on the other hand, can be traced all the way from the bar out to the edge of the galaxy and has plenty of star formation. Aside from a small gap, the arm is quite coherent. There is a long dust lane associated with the southern arm, which probably obscures some star formation. If there is dust along the northern arm, it is not nearly so well-behaved. Taking cuts across the northern and southern arms, we can measure the relative positions of the stars, H II regions, and H I. In the southern arm, the strongest H$`\alpha `$ and stellar emission is closely aligned with the peak of the H I. Large scale coherence breaks down in the northern half of the disk. At some places (40-50<sup>′′</sup> in cut f) star formation is well correlated with peaks in the H I distribution. At other locations ($``$40<sup>′′</sup> in cut g) star formation is displaced off of the H I peaks while elsewhere the H I and H$`\alpha `$ are anti-correlated (e.g. $``$50<sup>′′</sup> in cut e). The difference between the north and south arms suggests that there may be some fundamental difference in the nature of the two arms. Using the same cuts, we looked for evidence of triggering by the spiral arms. If there is triggering in the arms, we might expect a color gradient (implying an underlying age gradient) perpendicular to the arms. Such a color gradient can also be caused by dust extinction. The lack of a color gradient would imply the lack of an age gradient, and, hence a lack of triggering. While the cuts of the north arms show that the stellar emission is typically offset from the H$`\alpha `$ and H I emission, there is no offset present between the B & R peaks. Similarly, the southern arm shows no offset between the broadband emission peaks. The lack of an offset shows that there is no color gradient, and hence it is unlikely that an age gradient due to triggering in the spiral arm is occurring. This suggests that there is no larger scale organization of the star formation in the northern disk–a result reminiscent of irregular galaxies. Star formation must be driven by local conditions and processes in the northern half, while larger scale phenomena organize star formation in the bar and southern arm. Finally, we examined the nature of the NGC 925’s bar and the star formation occurring within it. As is expected for late-type bars (Phillips 1995), NGC 925 has star formation occurring all along the bar. Cuts perpendicular to the major axis of the bar show that the star formation is occurring along the major axis, but the broadband optical emission is systematically offset to the north of the H I peaks. While this could be a dust effect, we would expect the H$`\alpha `$ emission to be similarly affected, but it is not. This could also be due to the massive star formation in the bar ionizing or clearing out the H I in the bar. The systematic offset may be related to the dynamical center being offset from the bar center, and hence the bar could be a wave pattern moving through the gaseous medium in NGC 925. We might expect to see the peak of the stellar emission downstream of the star formation, as is the case here. It may be that the bar of NGC 925 is not actually a bar, but has more in common with the spiral arms of NGC 925. Overall, this study has shown that NGC 925 is a prototypical late-type spiral galaxy, with some properties that are characteristic of Magellanic spirals, such as an off-center bar and a dominant spiral arm. These traits are less pronounced in NGC 925 than in an SBm indicating a probable smooth transition of these properties from Scd galaxies to Sm’s. Massive star formation in NGC 925 persists beyond the radius predicted by Kennicutt (1989) for late-type spirals, but behaves more like dIrr’s (Hunter et al. 1998). This result may be exacerbated by wider field imaging of NGC 925, however, and star formation almost certainly continues at a low level beyond this point (see Ferguson et al. 1998 for an example). The north and south arms are not only different in brightness, but in coherence and in the distribution of stars across the arms. The north arm has stars offset from H I peaks, while the south arm has both star formation and stellar emission coincident with the neutral gas peaks. Finally, the bar of NGC 925 appears to be a typical late-type bar being gas-rich and having star formation throughout, but the stars are offset from the H$`\alpha `$ and H I peaks. This fact coupled with the offset dynamical center of NGC 925 suggests that the bar may be more of a wave phenomenon, similar to a spiral arm. The authors would like to thank the staff at the WIYN observatory and on Kitt Peak for helping make our observing run go smoothly, and for their excellent assistance in producing the wonderful data from the WIYN telescope. The authors also thank the anonymous referee for comments improving the quality of this paper. This work was partially supported by NSF grant AST 96-16907 to E.M.W. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. Figure Captions A combined B,V,R image from the WIYN 3.5m telescope. B is represented by blue, V by green, and R by red. The results from the isophote fitting of NGC 925. Counter-clockwise from upper left they are Right Ascension of center, Declination of center, ellipticity, and position angle. Blue represents fits to the B band image, Green for V, and Red for R. The results of the isophote fits to the surface brightness in B,V,R for NGC 925. The colors are the same as figure 2. The upper panel shows the data on an arbitrary scale (points) and the fits to the bar and disk (solid lines). The middle panel shows the residuals from just the disk fit. The lower panel shows the residuals from both the bar & disk fits. The solid line in the lower 2 panels mark the zero level. A representation of the model isophotes as fit to the R band image of NGC 925. Ellipses are plotted for one out of every three annuli. They illustrate the relative location and orientation of each annulus, but do not reflect the surface brightness distribution. This plot shows 4 centers of NGC 925: the center of NGC 925 according to the RC3 (RC3), the center of the bar (B), the center of the outer isophotes (ISO), and the dynamical center (DYN). The boxes represent the derived errors to the fitted centers. The ellipse is the FWHM of the beam from the H I observations in paper I. a) Top panel is a plot of gas surface density (solid line), as measured by H I observations from paper I, and multiplied by 1.47 to account for molecular gas. The dashed line is the critical gas surface density for collapse an $`\alpha `$=1. The dot-dash line is the radial H$`\alpha `$ distribution normalized to its peak. The vertical dotted line indicates the nearest edge of the chip from the center of the bar of NGC 925. The bottom panel indicates the ratio of the measured gas surface density to the critical density. The horizontal dotted line is the level when massive star formation should cease. The vertical horizontal line is the same as in the top panel. b) same as in a, but for $`\alpha `$=0.67. c) same as a, but for $`\alpha `$=0.3. This figure shows the H$`\alpha `$ emission from NGC 925 (in red) and the H I emission (in blue) overlaid on each other. Regions that are white are where there are peaks in both H$`\alpha `$ and H I. This figure shows the H I contours from paper I (in yellow) overlaid on the combined B,V,R WIYN image from figure 1. Hatched contours indicate a H I depression. The H I beam is shown in the lower right. This is an optical R band image of NGC 925, with the location of the cuts marked on it. The rectangular boxes are show the region over which the cut was taken. The letters correspond to the cuts shown in figure 10, and also marks the most positive point in the cut. The results of the cuts across the southern arm, northern arms, and bar of NGC 925. Each band of observations is normalized to its peak. The solid line is B band, the dashed line is R, the dot-dash line is H$`\alpha `$, and the dotted line is H I column density. Cuts a-d are for the southern arm. Cuts e-g are for the northern arms. Cut h is the bar major axis, and cuts i-k are parallel to the bar minor axis. The most positive points in the cuts are marked by the letters in figure 9.
warning/0004/math0004149.html	ar5iv	text	# On the second moment for primes in an arithmetic progression ## 1. Introduction and Statement of results In this paper we calculate a lower bound, of the same order of magnitude as conjectured, for the second moment of primes in an arithmetic progression. Specifically we examine $$I(x,h,q,a):=_x^{2x}\left(\psi (y+h;q,a)\psi (y;q,a)\frac{h}{\varphi (q)}\right)^2𝑑y$$ $`1.1`$ where $$\psi (x;q,a)=\underset{\genfrac{}{}{0pt}{}{nx}{na\left(modq\right)}}{}\mathrm{\Lambda }(n),$$ $`1.2`$ and $`\mathrm{\Lambda }`$ is the von Mangoldt function. We will take $$(a,q)=1,x2,1qhx,$$ $`1.3`$ (other ranges not being interesting). We shall assume the truth of the Generalized Riemann Hypothesis (GRH), which implies, in particular, $$E(x;q,a):=\psi (x;q,a)\frac{x}{\varphi (q)}x^{\frac{1}{2}}\mathrm{log}^2x,(qx).$$ $`1.4`$ The idea of our method originates from the work of Goldston for the case of all primes, corresponding in the present formulation to $`q=1`$. An improved and generalized version of this result appeared in as ###### Theorem A Assume GRH. Then for any $`ϵ>0`$ and $`1\frac{h}{q}\frac{x^{\frac{1}{3}}}{q^ϵ\mathrm{log}^3x}`$ we have $$\underset{\genfrac{}{}{0pt}{}{a\left(modq\right)}{(a,q)=1}}{}I(x,h,q,a)\frac{1}{2}xh\mathrm{log}\left(\left(\frac{q}{h}\right)^3x\right)O(xh(\mathrm{log}\mathrm{log}x)^3).$$ $`1.5`$ Moreover, for almost all $`q`$ with $`h^{3/4}\mathrm{log}^5xqh`$ we have $$I(x,h,q)xh\mathrm{log}(\frac{xq}{h}).$$ $`1.6`$ For an individual arithmetic progression Özlük proved unconditionally ###### Theorem B For $`1q(\mathrm{log}x)^{1\delta }`$, and $`h(\mathrm{log}x)^c`$ ($`\delta `$ and $`c`$ are any fixed positive numbers) satisfying $`qh`$, we have $$I(x,h,q,a)>(\frac{1}{2}ϵ)\frac{xh}{\varphi (q)}\mathrm{log}x$$ $`1.7`$ for any $`ϵ`$ and $`xX(ϵ,c)`$. We shall see below that the GRH implies a result of the type in Theorem B for much wider ranges of $`q`$ and $`h`$. An asymptotic estimate for $`I(x,h,q,a)`$ in certain ranges was shown by Yıldırım to be implied by GRH and a pair correlation conjecture for the zeros of Dirichlet’s $`L`$-functions. ###### Theorem C Assume GRH. Let $`\alpha _1,\alpha _2,\eta `$ be fixed and satisfying $`0<\eta <\alpha _1\alpha _21`$, and let $`\delta =x^\alpha `$ where $`\alpha _1\alpha \alpha _2`$. Assume, as $`x\mathrm{}`$, uniformly for $$q\mathrm{min}(x^{\frac{1}{2}}\delta ^{\frac{1}{2}}\mathrm{log}^Ax,\delta ^1x^\eta )\text{(}q\text{: prime or }1\text{)}$$ $`1.8`$ and $$\frac{x^{\alpha _1}}{\varphi (q)}\mathrm{log}^3xT\varphi (q)x^{\alpha _2}\mathrm{log}^3x$$ $`1.9`$ that, for $`(a,q)=1`$, $$\underset{\chi _1,\chi _2(modq)}{}\overline{\chi _1}(a)\chi _2(a)\underset{\genfrac{}{}{0pt}{}{0<\gamma _1,\gamma _2T}{\genfrac{}{}{0pt}{}{L({\scriptscriptstyle \frac{1}{2}}+i\gamma _1,\chi _1)=0}{L({\scriptscriptstyle \frac{1}{2}}+i\gamma _2,\chi _2)=0}}}{}x^{i(\gamma _1\gamma _2)}\frac{4}{4+(\gamma _1\gamma _2)^2}\varphi (q)\frac{T}{2\pi }\mathrm{log}qT.$$ $`1.10`$ Then $$_x^{2x}\left(\psi (u+u\delta ;q,a)\psi (u;q,a)\frac{u\delta }{\varphi (q)}\right)^2𝑑u\frac{3}{2}\frac{\delta x^2}{\varphi (q)}\mathrm{log}\frac{q}{\delta }$$ $`1.11`$ uniformly for $`x^{\alpha _2}\delta x^{\alpha _1}`$ and $`q`$ as in (1.8). It was also shown in that the left-hand side of (1.10) is $`\varphi (q)\frac{T}{2\pi }\mathrm{log}x`$ for $`1qx^{\frac{1}{2}}\mathrm{log}^3x`$ when $`\frac{x}{q}\mathrm{log}xTe^{x^{\frac{1}{4}}}`$. These asymptotic values are what the diagonal terms ($`\chi _1=\chi _2`$) would contribute, so the assumption (1.10) is a way of expressing that the zeros of different Dirichlet $`L`$-functions are uncorrelated. Theorem C is a generalization of one half of a result of Goldston and Montgomery for the case $`q=1`$, where an equivalence between the pair correlation conjecture for $`\zeta (s)`$ and the second moment for primes was established. Since the argument in works reversibly, a suitable converse to Theorem C is also provable. The restriction to prime $`q`$ was made in order to avoid the presence of imprimitive characters. The formula (1.11) involving differences $`u\delta `$ which vary with $`u`$ can be converted to a formula involving a fixed-difference $`h`$. Our main result is the following theorem. ###### Theorem 1 Assume GRH. Then for any $`ϵ>0`$ and $$qh(xq)^{\frac{1}{3}ϵ},$$ $`1.12`$ we have $$I(x,h,q,a)\frac{1}{2}\frac{xh}{\varphi (q)}\mathrm{log}(\frac{xq}{h^3})O(\frac{xh}{\varphi (q)}(\mathrm{log}\mathrm{log}3q)^3).$$ $`1.13`$ Notice that the conditions in (1.12) imply that both $`h`$ and $`q`$ are $`x^{\frac{1}{2}ϵ}`$. The proof of the theorem uses some new results on the function $`\lambda _R(n)`$ used as an approximation for the von Mangoldt function in our earlier work. Propositions 2, 3, and 4 embody these results, and we expect they will have further applications to other problems. ## 2. Preliminaries We shall need the following in our calculations. Let $$f(n,x,h)=_{[x,2x][nh,n)}1𝑑y=\{\begin{array}{cc}nx,\hfill & \text{for }xn<x+h\hfill \\ h,\hfill & \text{for }x+hn2x\hfill \\ 2xn+h,\hfill & \text{for }2x<n2x+h\hfill \\ 0,\hfill & \text{elsewhere}.\hfill \end{array}$$ $`2.1`$ ###### Lemma 1 For real numbers $`a_n`$ and $`b_n`$ we have $`{\displaystyle _x^{2x}}\left({\displaystyle \underset{y<ny+h}{}}a_n\right)`$ $`\left({\displaystyle \underset{y<my+h}{}}b_m\right)dy={\displaystyle \underset{x<n2x+h}{}}a_nb_nf(n,x,h)`$ $`+`$ $`{\displaystyle \underset{0<kh}{}}\left({\displaystyle \underset{x<n2x+hk}{}}(a_nb_{n+k}+a_{n+k}b_n)f(n,x,hk)\right).`$ ###### Lemma 2 Let $`C(x)=_{nx}c_n`$. Then we have $$\underset{x<n2x+h}{}c_nf(n,x,h)=_{2x}^{2x+h}C(u)𝑑u_x^{x+h}C(u)𝑑u.$$ Lemma 1 and Lemma 2 were proved in . We take this opportunity to correct a minor error in Lemma 1 of . In that lemma an extraneous term $`h(c_{x+h}c_{2x})`$ was incorrectly included and should be removed. This term then contributed an unnecessary error term in equations (2.7),(2.14), and (2.15) of . However these same error terms correctly occurred for a different reason in equation (2.9) so that starting with equation (2.16) these error terms were correctly included in the rest of . Calling $`\mathrm{\Delta }\psi =\psi (y+h;q,a)\psi (y;q,a)`$ for brevity, we have from (1.1) $$I(x,h,q,a)=_x^{2x}(\mathrm{\Delta }\psi )^2𝑑y\frac{2h}{\varphi (q)}_x^{2x}(\mathrm{\Delta }\psi )𝑑y+\frac{h^2x}{\varphi ^2(q)}.$$ By the above lemmas and (1.4) we obtain $`{\displaystyle _x^{2x}}(\mathrm{\Delta }\psi )𝑑y`$ $`={\displaystyle \underset{\genfrac{}{}{0pt}{}{x<n2x+h}{na\left(modq\right)}}{}}\mathrm{\Lambda }(n)f(n,x,h)`$ $`={\displaystyle \frac{xh}{\varphi (q)}}+{\displaystyle _{2x}^{2x+h}}E(y;q,a)𝑑y{\displaystyle _x^{x+h}}E(y;q,a)𝑑y`$ $`={\displaystyle \frac{xh}{\varphi (q)}}+O(hx^{\frac{1}{2}}\mathrm{log}^2x),`$ so that $$I(x,h,q,a)=_x^{2x}(\mathrm{\Delta }\psi )^2𝑑y\frac{xh^2}{\varphi ^2(q)}+O(\frac{x^{\frac{1}{2}}h^2\mathrm{log}^2x}{\varphi (q)}).$$ $`2.2`$ The integral $`(\mathrm{\Delta }\psi )^2`$ leads to sums of the sort $`\mathrm{\Lambda }(n)\mathrm{\Lambda }(n+k)`$ which are in the territory of the twin prime conjecture. In the uninteresting case $`1hq`$, only the sum $`\mathrm{\Lambda }^2(n)`$ is present, giving easily the evaluation $$I(x,h,q,a)=\frac{xh}{\varphi (q)}\mathrm{log}x\frac{xh}{\varphi (q)}\frac{xh^2}{\varphi ^2(q)}+O(x^{\frac{1}{2}}h\mathrm{log}^3x),(hqx).$$ $`2.3`$ Now let $`\lambda _R(n)`$ be any arithmetical function, and set $$\psi _R(y;q,a)=\underset{\genfrac{}{}{0pt}{}{ny}{na\left(modq\right)}}{}\lambda _R(n);\mathrm{\Delta }\psi _R=\psi _R(y+h;q,a)\psi _R(y;q,a).$$ $`2.4`$ Trivially $`(\mathrm{\Delta }\psi \mathrm{\Delta }\psi _R)^20`$, so that $$_x^{2x}(\mathrm{\Delta }\psi )^2𝑑y2_x^{2x}(\mathrm{\Delta }\psi )(\mathrm{\Delta }\psi _R)𝑑y_x^{2x}(\mathrm{\Delta }\psi _R)^2𝑑y.$$ $`2.5`$ We apply Lemma 1 and Lemma 2 to these integrals to obtain $`{\displaystyle _x^{2x}}(\mathrm{\Delta }\psi )`$ $`(\mathrm{\Delta }\psi _R)dy={\displaystyle \underset{\genfrac{}{}{0pt}{}{x<n2x+h}{na\left(modq\right)}}{}}\mathrm{\Lambda }(n)\lambda _R(n)f(n,x,h)`$ $`2.6a`$$`2.6b`$ $`+{\displaystyle \underset{\genfrac{}{}{0pt}{}{0<kh}{qk}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{x<n2x+hk}{na\left(modq\right)}}{}}[\mathrm{\Lambda }(n)\lambda _R(n+k)+\mathrm{\Lambda }(n+k)\lambda _R(n)]f(n,x,hk)`$ $`=`$ $`\left({\displaystyle _{2x}^{2x+h}}{\displaystyle _x^{x+h}}\right){\displaystyle \underset{\genfrac{}{}{0pt}{}{nu}{na\left(modq\right)}}{}}\mathrm{\Lambda }(n)\lambda _R(n)du`$ $`+{\displaystyle _{2x}^{2x+h}}{\displaystyle \underset{0<\|j\|\frac{u2x}{q}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{N_1<nN_2}{na\left(modq\right)}}{}}\lambda _R(n)\mathrm{\Lambda }(n+jq)du`$ $`{\displaystyle _x^{x+h}}{\displaystyle \underset{0<\|j\|\frac{ux}{q}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{N_1<nN_2}{na\left(modq\right)}}{}}\lambda _R(n)\mathrm{\Lambda }(n+jq)du,`$ where $`N_1=\mathrm{max}(0,jq)`$ and $`N_2=\mathrm{min}(u,ujq)`$. Similarly $`{\displaystyle _x^{2x}}(\mathrm{\Delta }\psi _R)^2𝑑y=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{x<n2x+h}{na\left(modq\right)}}{}}\lambda _R^2(n)f(n,x,h)`$ $`2.7`$ $`+2{\displaystyle \underset{\genfrac{}{}{0pt}{}{0<kh}{qk}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{x<n2x+hk}{na\left(modq\right)}}{}}\lambda _R(n)\lambda _R(n+k)f(n,x,hk)`$ $`=`$ $`\left({\displaystyle _{2x}^{2x+h}}{\displaystyle _x^{x+h}}\right){\displaystyle \underset{\genfrac{}{}{0pt}{}{nu}{na\left(modq\right)}}{}}\lambda _R^2(n)du`$ $`+2{\displaystyle _{2x}^{2x+h}}{\displaystyle \underset{0<j\frac{u2x}{q}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{nujq}{na\left(modq\right)}}{}}\lambda _R(n)\lambda _R(n+jq)du`$ $`2{\displaystyle _x^{x+h}}{\displaystyle \underset{0<j\frac{ux}{q}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{nujq}{na\left(modq\right)}}{}}\lambda _R(n)\lambda _R(n+jq)du.`$ ## 3. The choice of $`\lambda _R(n)`$ and some number-theoretic sums As the auxiliary function we use $$\lambda _R(n):=\underset{rR}{}\frac{\mu ^2(r)}{\varphi (r)}\underset{d(r,n)}{}d\mu (d).$$ $`3.1`$ This function is known (, ) to exhibit behavior similar to $`\mathrm{\Lambda }(n)`$ when considered on average in arithmetic progressions, and it has been employed in related problems (, , , ). An upper bound for $`\lambda _R(n)`$ is $$\|\lambda _R(n)\|\underset{dn}{}d\underset{\genfrac{}{}{0pt}{}{rR}{dr}}{}\frac{1}{\varphi (r)}\underset{rR}{\mathrm{max}}(\frac{r}{\varphi (r)})\underset{dn}{}d\underset{\genfrac{}{}{0pt}{}{rR}{dr}}{}\frac{1}{r}d(n)\mathrm{log}R\mathrm{log}\mathrm{log}R.$$ $`3.2`$ To evaluate the sums which arise when (3.1) is used in (2.6) and (2.7) we shall need some lemmas. In the following $`p`$ will denote a prime number. ###### Lemma 3 (Hildebrand ) We have for each positive integer $`k`$, uniformly in $`R1`$, $$L_k(R):=\underset{\genfrac{}{}{0pt}{}{nR}{(n,k)=1}}{}\frac{\mu ^2(n)}{\varphi (n)}=\frac{\varphi (k)}{k}(\mathrm{log}R+c+v(k))+O(\frac{w(k)}{\sqrt{R}}),$$ $`3.3`$ where $`c:=`$ $`\gamma +{\displaystyle \underset{p}{}}{\displaystyle \frac{\mathrm{log}p}{p(p1)}};v(k):={\displaystyle \underset{pk}{}}{\displaystyle \frac{\mathrm{log}p}{p}};`$ $`3.4`$ $`w(k):=`$ $`{\displaystyle \underset{dk}{}}{\displaystyle \frac{\mu ^2(d)}{\sqrt{d}}}={\displaystyle \underset{pk}{}}(1+{\displaystyle \frac{1}{\sqrt{p}}});v(1)=0,w(1)=1.`$ ###### Lemma 4 We have $$v(k)\mathrm{log}\mathrm{log}3k,$$ $`3.5`$ $$\underset{pk}{}\frac{1}{\sqrt{p}}\frac{\sqrt{\mathrm{log}k}}{\mathrm{log}\mathrm{log}3k},$$ $`3.6`$ and $$g(k):=\underset{pk}{}(1+\frac{p}{p1})2^{\nu (k)}(\mathrm{log}\mathrm{log}3k)$$ $`3.7`$ ###### Demonstration Proof We show (3.7); the other inequalities can be proved similarly. Let $`\nu (k)`$ be the number of distinct prime factors of $`k`$, which satisfies the bound $`\nu (k){\displaystyle \frac{\mathrm{log}k}{\mathrm{log}\mathrm{log}k}}`$. We have $`\mathrm{log}g(k)`$ $`={\displaystyle \underset{pk}{}}\mathrm{log}(2+{\displaystyle \frac{1}{p1}})`$ $`<\nu (k)\mathrm{log}2+{\displaystyle \underset{pk}{}}{\displaystyle \frac{1}{p}}`$ $`<\nu (k)\mathrm{log}2+{\displaystyle \underset{p2\mathrm{log}2k}{}}{\displaystyle \frac{1}{p}}`$ $`=\nu (k)\mathrm{log}2+\mathrm{log}\mathrm{log}\mathrm{log}21k+O(1),`$ where the prime number theorem and Mertens’ theorem have been employed. Exponentiating both sides we obtain (3.7). ###### Lemma 5 We have $$\underset{\mathrm{}k}{}\frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}\mathrm{log}\mathrm{}=\frac{k}{\varphi (k)}v(k).$$ $`3.8`$ ###### Demonstration Proof $`{\displaystyle \underset{\mathrm{}k}{}}`$ $`{\displaystyle \frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}}\mathrm{log}\mathrm{}=\mathrm{log}{\displaystyle \underset{\mathrm{}k}{}}\mathrm{}^{{\displaystyle \frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}}}=\mathrm{log}{\displaystyle \underset{pk}{}}p^{{\displaystyle \underset{\mathrm{}k,p\mathrm{}}{}}{\displaystyle \frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}}}`$ $`=\mathrm{log}{\displaystyle \underset{pk}{}}p^{{\displaystyle \frac{1}{\varphi (p)}}{\displaystyle \underset{\mathrm{}^{}{\scriptscriptstyle \frac{k}{p}},(\mathrm{}^{},p)=1}{}}{\displaystyle \frac{\mu ^2(\mathrm{}^{})}{\varphi (\mathrm{}^{})}}}=\mathrm{log}{\displaystyle \underset{pk}{}}p^{{\displaystyle \frac{1}{\varphi (p)}}{\displaystyle \frac{1}{1+\frac{1}{\varphi (p)}}}{\displaystyle \underset{p^{}k}{}}\left(1+{\displaystyle \frac{1}{\varphi (p^{})}}\right)}`$ $`=\mathrm{log}{\displaystyle \underset{pk}{}}p^{{\displaystyle \frac{1}{p}}{\displaystyle \frac{k}{\varphi (k)}}}={\displaystyle \frac{k}{\varphi (k)}}{\displaystyle \underset{pk}{}}{\displaystyle \frac{\mathrm{log}p}{p}}.`$ ###### Lemma 6 (Goldston ) We have $$\underset{rR}{}\frac{\mu ^2(r)}{\varphi (r)}\underset{\genfrac{}{}{0pt}{}{dr}{(d,k)=1}}{}\frac{d\mu (d)}{\varphi (d)}=\text{S}(k)+O(\frac{kd(k)}{R\varphi (k)}),$$ $`3.9`$ where $$\text{S}(k)=\{\begin{array}{cc}\hfill 2C\underset{\genfrac{}{}{0pt}{}{p\|k}{p>2}}{}\left(\frac{p1}{p2}\right),& \text{if }k\text{ is even, }k0\text{;}\hfill \\ \hfill 0,& \text{if }k\text{ is odd;}\hfill \end{array}$$ $`3.10`$ with $$C=\underset{p>2}{}\left(1\frac{1}{(p1)^2}\right).$$ $`3.11`$ ###### Demonstration Proof The proof can be found in ; we just note that $$\underset{\genfrac{}{}{0pt}{}{dr}{(d,k)=1}}{}\frac{d\mu (d)}{\varphi (d)}=\frac{\mu (r)}{\varphi (r)}\mu ((r,k))\varphi ((r,k)),$$ $`3.12`$ so the left-hand side of (3.9) may be expressed as $$\underset{r=1}{\overset{\mathrm{}}{}}\frac{\mu (r)\mu ((r,k))\varphi ((r,k))}{\varphi ^2(r)}+O(\underset{r>R}{}\frac{\mu ^2(r)\mu ^2((r,k))\varphi ((r,k))}{\varphi ^2(r)}).$$ $`3.13`$ Here the first sum is $`\text{S}(k)`$ and the error term is $``$ the $`O`$-term in (3.9). ###### Lemma 7 (Goldston and Friedlander )We have $$\underset{0<j\frac{h}{q}}{}(hjq)\text{S}(jq)=\frac{h^2}{2\varphi (q)}\frac{h}{2}\mathrm{log}\frac{h}{q}+O(h(\mathrm{log}\mathrm{log}3q)^3).$$ $`3.14`$ ###### Lemma 8 (Hooley ) Assuming GRH, we have $$\underset{\genfrac{}{}{0pt}{}{a\left(modq\right)}{(a,q)=1}}{}\underset{ux}{\mathrm{max}}\|E(u;q,a)\|^2x\mathrm{log}^4x;\text{for}qx.$$ $`3.15`$ ###### Lemma 9 We have $`{\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)\sigma (r)}{\varphi (r)}}`$ $`R`$ $`3.16`$$`3.17`$$`3.18`$ $`{\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)\sigma _{\frac{1}{2}}(r)}{\varphi (r)}}`$ $`\sqrt{R}`$ $`{\displaystyle \underset{0<rR}{}}{\displaystyle \frac{rd(r)}{\varphi (r)}}`$ $`R\mathrm{log}2R.`$ ###### Demonstration Proof To prove (3.16), note $`{\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)\sigma (r)}{\varphi (r)}}=`$ $`{\displaystyle \underset{rR}{}}\mu ^2(r){\displaystyle \underset{pr}{}}(1+{\displaystyle \frac{2}{p1}})={\displaystyle \underset{rR}{}}\mu ^2(r){\displaystyle \underset{dr}{}}{\displaystyle \frac{2^{\nu (d)}}{\varphi (d)}}`$ $`={\displaystyle \underset{dR}{}}{\displaystyle \frac{\mu ^2(d)2^{\nu (d)}}{\varphi (d)}}{\displaystyle \underset{\mathrm{}\frac{R}{d}}{}}\mu ^2(\mathrm{})R{\displaystyle \underset{dR}{}}{\displaystyle \frac{\mu ^2(d)2^{\nu (d)}}{d\varphi (d)}}`$ $`R{\displaystyle \underset{p}{}}(1+{\displaystyle \frac{2}{p(p1)}})R.`$ The proof of (3.17) is similar, and (3.18) was shown in . ## 4. The proof of the Theorem In this section we calculate the right-hand sides of (2.6) and (2.7), and so obtain our result. ###### Proposition 1 Assuming GRH, we have $$\underset{\genfrac{}{}{0pt}{}{nN}{na\left(modq\right)}}{}\mathrm{\Lambda }(n)\lambda _R(n)=\frac{N\mathrm{log}R}{\varphi (q)}+\frac{cN}{\varphi (q)}+O(\frac{N}{\varphi (q)\sqrt{R}})+O(N^{\frac{1}{2}}\mathrm{log}^3N)+O(R\mathrm{log}N).$$ $`4.1`$ ###### Demonstration Proof Starting from the definition (3.1) and recalling $`L_k(R)`$ from (3.3), we have $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{nN}{na\left(modq\right)}}{}}\mathrm{\Lambda }(n)\lambda _R(n)`$ $`={\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)}{\varphi (r)}}{\displaystyle \underset{dr}{}}d\mu (d){\displaystyle \underset{\genfrac{}{}{0pt}{}{nN}{\genfrac{}{}{0pt}{}{na\left(modq\right)}{dn}}}{}}\mathrm{\Lambda }(n)`$ $`=L_1(R)\psi (N;q,a){\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)}{\varphi (r)}}{\displaystyle \underset{pr}{}}p\mathrm{log}p{\displaystyle \underset{\genfrac{}{}{0pt}{}{k1}{\genfrac{}{}{0pt}{}{p^kN}{p^ka\left(modq\right)}}}{}}1.`$ Here the sum over $`k`$ is trivially of size $`O({\displaystyle \frac{\mathrm{log}N}{\mathrm{log}p}})`$, so that by Lemma 3 $`{\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)}{\varphi (r)}}{\displaystyle \underset{pr}{}}p\mathrm{log}p{\displaystyle \underset{\genfrac{}{}{0pt}{}{k1}{\genfrac{}{}{0pt}{}{p^kN}{p^ka\left(modq\right)}}}{}}1`$ $`\mathrm{log}N{\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)}{\varphi (r)}}{\displaystyle \underset{pr}{}}p`$ $`=\mathrm{log}N{\displaystyle \underset{pR}{}}{\displaystyle \frac{p}{\varphi (p)}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{m{\scriptscriptstyle \frac{R}{p}}}{(m,p)=1}}{}}{\displaystyle \frac{\mu ^2(m)}{\varphi (m)}}`$ $`\mathrm{log}N{\displaystyle \underset{pR}{}}\mathrm{log}{\displaystyle \frac{R}{p}}R\mathrm{log}N.`$ By the prime number theorem we obtain (4.1). In order for the main term to dominate the error terms in (4.1) we will require that $$qR\frac{N}{\mathrm{log}N},q\frac{N^{\frac{1}{2}}}{\mathrm{log}^3N}.$$ $`4.2`$ Hence the relevant contribution to (2.6 b) will be $`\left({\displaystyle _{2x}^{2x+h}}{\displaystyle _x^{x+h}}\right)`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{nu}{na\left(modq\right)}}{}}\mathrm{\Lambda }(n)\lambda _R(n)du`$ $`4.3`$ $`=`$ $`{\displaystyle \frac{xh}{\varphi (q)}}(\mathrm{log}R+c)+O({\displaystyle \frac{xh}{\varphi (q)\sqrt{R}}})+O(x^{\frac{1}{2}}h\mathrm{log}^3x)+O(Rh\mathrm{log}x).`$ ###### Proposition 2 Assuming GRH, we have for $$1qhx,qRx,$$ $`4.4`$ that $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{0<kh}{qk}}{}}`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{x<n2x+hk}{na\left(modq\right)}}{}}[\mathrm{\Lambda }(n)\lambda _R(n+k)+\mathrm{\Lambda }(n+k)\lambda _R(n)]f(n,x,hk)`$ $`4.5`$ $`=`$ $`{\displaystyle \frac{xh^2}{\varphi ^2(q)}}{\displaystyle \frac{xh}{\varphi (q)}}\mathrm{log}{\displaystyle \frac{h}{q}}+O({\displaystyle \frac{xh}{\varphi (q)}}(\mathrm{log}\mathrm{log}3q)^3)+O({\displaystyle \frac{xh^2d(q)}{\varphi ^2(q)R}}\mathrm{log}{\displaystyle \frac{2h}{q}})`$ $`+O({\displaystyle \frac{x^{\frac{1}{2}}h^{\frac{3}{2}}R\mathrm{log}^2x}{q^{\frac{1}{2}}}})+O({\displaystyle \frac{x^{\frac{1}{2}}h^2R^{\frac{1}{2}}\mathrm{log}^2x}{q}}).`$ ###### Demonstration Proof We have $$\underset{\genfrac{}{}{0pt}{}{N_1<nN_2}{na\left(modq\right)}}{}\lambda _R(n)\mathrm{\Lambda }(n+jq)=\underset{rR}{}\frac{\mu ^2(r)}{\varphi (r)}\underset{dr}{}d\mu (d)\underset{\genfrac{}{}{0pt}{}{N_1<nN_2}{\genfrac{}{}{0pt}{}{na\left(modq\right)}{dn}}}{}\mathrm{\Lambda }(n+jq).$$ We may write the innermost sum as $$\underset{\genfrac{}{}{0pt}{}{N_1+jq<mN_2+jq}{\genfrac{}{}{0pt}{}{ma\left(modq\right)}{mjq\left(modd\right)}}}{}\mathrm{\Lambda }(m).$$ Here $`mjq=\mathrm{}d`$ for some integer $`\mathrm{}`$, and so $`a\mathrm{}d(modq)`$. Since $`(a,q)=1`$, we can include only those $`d`$’s such that $`(d,q)=1`$. Then there is a unique $`b,\mathrm{\hspace{0.17em}\hspace{0.17em}0}<b<qd`$, such that $`mb(modqd)`$. We know $`(m,q)=1`$, so that $`(m,d)=1`$ if and only if $`(j,d)=1`$. Hence the innermost sum is equal to $`\psi (N_2+jq;qd,b)`$ $`\psi (N_1+jq;qd,b)`$ $`={\displaystyle \frac{N_2N_1}{\varphi (qd)}}E_{qd,b}+E(N_2+jq;qd,b)E(N_1+jq;qd,b),`$ where $`E_{qd,b}=1`$ if $`(qd,b)=1`$, and $`E_{qd,b}=0`$ if $`(qd,b)>1`$. Thus $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{N_1<nN_2}{na\left(modq\right)}}{}}\lambda _R(n)`$ $`\mathrm{\Lambda }(n+jq)={\displaystyle \frac{u\|j\|q}{\varphi (q)}}{\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)}{\varphi (r)}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{dr}{(d,jq)=1}}{}}{\displaystyle \frac{d\mu (d)}{\varphi (d)}}`$ $`4.6`$ $`+`$ $`{\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)}{\varphi (r)}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{dr}{(d,q)=1}}{}}d\mu (d)[E(N_2+jq;qd,b)E(N_1+jq;qd,b)],`$ where the first term on the right-hand side is the main term, its value settled by Lemma 6, and its contribution to (2.6 b) will be $`{\displaystyle _{2x}^{2x+h}}`$ $`{\displaystyle \underset{0<\|j\|\frac{u2x}{q}}{}}{\displaystyle \frac{u\|j\|q}{\varphi (q)}}[\text{S}(jq)+O({\displaystyle \frac{jqd(jq)}{R\varphi (jq)}})]du`$ $`4.7`$ $`{\displaystyle _x^{x+h}}{\displaystyle \underset{0<\|j\|\frac{ux}{q}}{}}{\displaystyle \frac{u\|j\|q}{\varphi (q)}}[\text{S}(jq)+O({\displaystyle \frac{jqd(jq)}{R\varphi (jq)}})]du`$ $`=`$ $`{\displaystyle \frac{2x}{\varphi (q)}}{\displaystyle \underset{0<j\frac{h}{q}}{}}(hjq)\text{S}(jq)+O({\displaystyle \frac{xhqd(q)}{R\varphi ^2(q)}}{\displaystyle \underset{0<j\frac{h}{q}}{}}{\displaystyle \frac{jd(j)}{\varphi (j)}})`$ $`=`$ $`{\displaystyle \frac{xh^2}{\varphi ^2(q)}}{\displaystyle \frac{xh}{\varphi (q)}}\mathrm{log}{\displaystyle \frac{h}{q}}+O({\displaystyle \frac{xh}{\varphi (q)}}(\mathrm{log}\mathrm{log}3q)^3)+O({\displaystyle \frac{xh^2d(q)}{\varphi ^2(q)R}}\mathrm{log}{\displaystyle \frac{2h}{q}}),`$ by Lemma 7 and (3.18). For the second term in the right-hand side of (4.6), if we use (1.4) directly, we will get the upper bound $`Rx^{\frac{1}{2}}\mathrm{log}^2x`$, by (3.16). This will lead to a contribution of $`O({\displaystyle \frac{x^{\frac{1}{2}}h^2R}{q}}\mathrm{log}^2x)`$ in (2.6 b). Instead, in view of the averaging over $`j`$ in (2.6 b), we will use Hooley’s estimate quoted as Lemma 8 above. To do this note that some of the $`d`$’s may not be coprime to $`b`$, but we can discard them (from the $`j`$\- and $`n`$-summations) with an error $$\underset{\genfrac{}{}{0pt}{}{0<\left\|j\right\|{\scriptscriptstyle \frac{h}{q}}}{(j,d)>1}}{}\psi (3x;qd,b)\underset{0<\|j\|\frac{h}{q}}{}\underset{pd}{}\underset{\genfrac{}{}{0pt}{}{n3x}{pn}}{}\mathrm{\Lambda }(n)\frac{h}{q}\mathrm{log}^2x,$$ and this leads to an error of $`O({\displaystyle \frac{h^2R}{q}}\mathrm{log}^2x)`$ in (2.6 b). Hence the contribution to (2.6 b) from the second term on the right-hand side of (4.6) is $``$ $`{\displaystyle \frac{h^2R}{q}}\mathrm{log}^2x+\left({\displaystyle _{2x}^{2x+h}}{\displaystyle _x^{x+h}}\right){\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)}{\varphi (r)}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{dr}{(d,q)=1}}{}}d{\displaystyle \underset{\genfrac{}{}{0pt}{}{0<\left\|j\right\|{\scriptscriptstyle \frac{h}{q}}}{(j,d)=1}}{}}\underset{u2x+h}{\mathrm{max}}\|E(u;qd,b)\|du`$ $`4.8`$ $``$ $`{\displaystyle \frac{h^2R}{q}}\mathrm{log}^2x+h{\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)}{\varphi (r)}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{dr}{(d,q)=1}}{}}d({\displaystyle \frac{h}{q}})^{\frac{1}{2}}({\displaystyle \underset{\genfrac{}{}{0pt}{}{0<\left\|j\right\|{\scriptscriptstyle \frac{h}{q}}}{(j,d)=1}}{}}\underset{u2x+h}{\mathrm{max}}\|E(u;qd,b)\|^2)^{\frac{1}{2}}`$ $``$ $`{\displaystyle \frac{h^2R}{q}}\mathrm{log}^2x+{\displaystyle \frac{h^{\frac{3}{2}}}{q^{\frac{1}{2}}}}{\displaystyle \underset{rR}{}}{\displaystyle \frac{\mu ^2(r)}{\varphi (r)}}{\displaystyle \underset{dr}{}}d(1+({\displaystyle \frac{h}{qd}}))^{\frac{1}{2}}({\displaystyle \underset{\genfrac{}{}{0pt}{}{j\left(modd\right)}{(j,d)=1}}{}}\underset{u3x}{\mathrm{max}}\|E(u;qd,b)\|^2)^{\frac{1}{2}}.`$ In the last sum as $`j`$ runs through the reduced residues modulo $`d`$, $`b`$ runs through those elements of the set $`\{a,a+q,\mathrm{},a+(d1)q\}`$ which are relatively prime to $`d`$ (note that $`a\mathrm{}d(modq)`$ and $`(a,q)=1`$ implies $`(a,d)=1`$), and this correspondence is one-to-one. This is because $`ma(modq)`$ and $`mjq(modd)`$ if and only if $`mn_1da+n_2jq^2(modqd)`$ where $`n_i`$ satisfy $`dn_11(modq),qn_21(modd)`$, and we have $`n_1da+n_2jq^2a+tq(modqd)`$ if and only if $`jtan_2(modd)`$. Hence we may replace the $`j`$-sum in (4.8) by $$\underset{\genfrac{}{}{0pt}{}{t\left(modd\right)}{(a+tq,d)=1}}{}\underset{1u3x}{\mathrm{max}}\|E(u;qd,a+tq)\|^2.$$ $`4.9`$ Although the last sum is over only $`{\displaystyle \frac{1}{\varphi (q)}}`$ of the reduced residue classes modulo $`qd`$, we shall use Hooley’s estimate as is. One would want to get a Hooley-type estimate for (4.9) itself, thereby saving a factor of $`\varphi (q)`$, but this seems to require some estimates for certain integrals involving pairs of $`L`$-functions. We do not follow this path now. Recall that Theorem C, which gives an asymptotic estimate for our integral already rests upon such an assumption, (1.10), about $`L`$-functions. By Lemma 8, we take $`x\mathrm{log}^4x`$ as upper bound for (4.9) on the condition that $`qRx`$, and on applying Lemma 9 we obtain that the expression in (4.8) is $$\frac{x^{\frac{1}{2}}h^{\frac{3}{2}}R}{q^{\frac{1}{2}}}\mathrm{log}^2x+\frac{x^{\frac{1}{2}}h^2R^{\frac{1}{2}}}{q}\mathrm{log}^2x.$$ $`4.10`$ This completes the proof of Proposition 2. ###### Proposition 3 For $`(a,q)=1`$ we have $$\underset{\genfrac{}{}{0pt}{}{nN}{na\left(modq\right)}}{}\lambda _R^2(n)=\frac{N}{\varphi (q)}[\mathrm{log}R+c+O(v(q))+O(R^{\frac{1}{2}+ϵ})]+O(R^2).$$ $`4.11`$ ###### Proposition 4 For $`(a,q)=1`$ and $`j0`$ we have $$\underset{\genfrac{}{}{0pt}{}{nN}{na\left(modq\right)}}{}\lambda _R(n)\lambda _R(n+jq)=\frac{N}{\varphi (q)}\text{S}(jq)+O(\frac{Ng(q)}{\varphi (q)R}\frac{jd(j)}{\varphi (j)})+O(R^2).$$ $`4.12`$ ###### Demonstration Proof The beginning of the proof of Proposition 3 may be incorporated into that of Proposition 4 upon a notational stipulation for the case $`j=0`$. When the positive integer $`t`$ satisfies $`tj`$, if $`j=0`$ we will understand that $`t`$ can be any positive integer; and we will take $`(t,0)=t`$. By definition (3.1), $$\underset{\genfrac{}{}{0pt}{}{nN}{na\left(modq\right)}}{}\lambda _R(n)\lambda _R(n+jq)=\underset{r,r^{}R}{}\frac{\mu ^2(r)\mu ^2(r^{})}{\varphi (r)\varphi (r^{})}\underset{\genfrac{}{}{0pt}{}{dr}{er^{}}}{}d\mu (d)e\mu (e)\underset{\genfrac{}{}{0pt}{}{nN}{\genfrac{}{}{0pt}{}{na\left(modq\right)}{dn,en+jq}}}{}1.$$ $`4.13`$ In the innermost sum, the conditions on $`n`$, $`d`$, and $`e`$ imply $`(q,de)=1,(d,e)j`$, and thus $`n`$ belongs to a unique residue class modulo $`[q,d,e]`$. Hence we have $$\underset{\genfrac{}{}{0pt}{}{nN}{\genfrac{}{}{0pt}{}{na\left(modq\right)}{dn,en+jq}}}{}1=\frac{N}{[q,d,e]}+O(1).$$ $`4.14`$ The contribution of the $`O(1)`$-term in (4.14) to (4.13) is $$\underset{r,r^{}R}{}\frac{\mu ^2(r)\mu ^2(r^{})}{\varphi (r)\varphi (r^{})}\underset{\genfrac{}{}{0pt}{}{dr}{er^{}}}{}de=(\underset{rR}{}\frac{\mu ^2(r)\sigma (r)}{\varphi (r)})^2R^2$$ $`4.15`$ by (3.16), and this is where the $`O(R^2)`$-term in (4.11) and (4.12) comes from. Hence $$\underset{\genfrac{}{}{0pt}{}{nN}{na\left(modq\right)}}{}\lambda _R(n)\lambda _R(n+jq)=\frac{N}{q}\underset{r,r^{}R}{}\frac{\mu ^2(r)\mu ^2(r^{})}{\varphi (r)\varphi (r^{})}\underset{\genfrac{}{}{0pt}{}{d{\scriptscriptstyle \frac{r}{(r,q)}}}{\genfrac{}{}{0pt}{}{e{\scriptscriptstyle \frac{r^{}}{(r^{},q)}}}{(d,e)j}}}{}\mu (d)\mu (e)(d,e)+O(R^2).$$ $`4.16`$ Let $`(d,e)=\delta ,d=d^{}\delta ,e=e^{}\delta `$, so that $`(d^{},e^{})=1`$. The inner sums over $`d`$ and $`e`$ become $$\underset{\genfrac{}{}{0pt}{}{\delta j}{\delta ({\scriptscriptstyle \frac{r}{(r,q)}},{\scriptscriptstyle \frac{r^{}}{(r^{},q)}})}}{}\delta \underset{d^{}\frac{r}{\delta (r,q)}}{}\mu (d^{})\underset{\genfrac{}{}{0pt}{}{e^{}{\scriptscriptstyle \frac{r^{}}{\delta (r^{},q)}}}{(e^{},d^{})=1}}{}\mu (e^{}).$$ $`4.17`$ Here the innermost sum is $$\underset{\genfrac{}{}{0pt}{}{e^{}{\scriptscriptstyle \frac{r^{}}{\delta (r^{},q)}}}{(e^{},d^{})=1}}{}\mu (e^{})=\underset{\genfrac{}{}{0pt}{}{p{\scriptscriptstyle \frac{r^{}}{\delta (r^{},q)}}}{pd^{}}}{}(1+\mu (p))=\{\begin{array}{cc}& 1\text{if }\frac{r^{}}{\delta (r^{},q)}d^{}\text{ ;}\hfill \\ & 0\text{otherwise.}\hfill \end{array}$$ $`4.18`$ Next the sum over $`d^{}`$ becomes $$\underset{\genfrac{}{}{0pt}{}{d^{}{\scriptscriptstyle \frac{r}{\delta (r,q)}}}{{\scriptscriptstyle \frac{r^{}}{\delta (r^{},q)}}d^{}}}{}\mu (d^{})=\{\begin{array}{cc}\hfill \mu (\frac{r}{\delta (r,q)})& \text{if }\frac{r^{}}{(r^{},q)}=\frac{r}{(r,q)}\text{ ;}\hfill \\ \hfill 0& \text{otherwise,}\hfill \end{array}$$ $`4.19`$ so the main term of (4.16) is $$\frac{N}{q}\underset{\genfrac{}{}{0pt}{}{r,r^{}R}{\frac{r}{(r,q)}=\frac{r^{}}{(r^{},q)}}}{}\frac{\mu ^2(r)\mu ^2(r^{})}{\varphi (r)\varphi (r^{})}\mu (\frac{r}{(r,q)})\underset{\delta (\frac{r}{(r,q)},j)}{}\delta \mu (\delta ).$$ $`4.20`$ Since $$\underset{\delta (\frac{r}{(r,q)},j)}{}\delta \mu (\delta )=\mu ((\frac{r}{(r,q)},j))\varphi ((\frac{r}{(r,q)},j)),$$ $`4.21`$ the main term is $$\frac{N}{q}\underset{\genfrac{}{}{0pt}{}{r,r^{}R}{\frac{r}{(r,q)}=\frac{r^{}}{(r^{},q)}}}{}\frac{\mu ^2(r)\mu ^2(r^{})}{\varphi (r)\varphi (r^{})}\mu (\frac{r}{(r,q)})\mu ((\frac{r}{(r,q)},j))\varphi ((\frac{r}{(r,q)},j)).$$ $`4.22`$ Writing $`(r,q)=\mathrm{},(r^{},q)=m,r=\mathrm{}s,r^{}=ms`$ where $`(s,q)=1`$, (4.22) takes the form $$\frac{N}{q}\underset{\mathrm{}q}{}\frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}\underset{mq}{}\frac{\mu ^2(m)}{\varphi (m)}\underset{\genfrac{}{}{0pt}{}{s\mathrm{min}({\scriptscriptstyle \frac{R}{\mathrm{}}},{\scriptscriptstyle \frac{R}{m}})}{(s,q)=1}}{}\frac{\mu (s)}{\varphi ^2(s)}\mu ((s,j))\varphi ((s,j)).$$ $`4.23`$ The $`j=0`$ case: We rewrite (4.23) as $$\frac{N}{q}\underset{\mathrm{}q}{}\frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}\underset{mq}{}\frac{\mu ^2(m)}{\varphi (m)}\underset{\genfrac{}{}{0pt}{}{s\mathrm{min}({\scriptscriptstyle \frac{R}{\mathrm{}}},{\scriptscriptstyle \frac{R}{m}})}{(s,q)=1}}{}\frac{\mu ^2(s)}{\varphi (s)}.$$ $`4.24`$ It is convenient to regard (4.24) as $$\frac{N}{q}\underset{\mathrm{}q}{}\frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}\underset{mq}{}\frac{\mu ^2(m)}{\varphi (m)}\underset{\genfrac{}{}{0pt}{}{s{\scriptscriptstyle \frac{R}{\mathrm{}}}}{(s,q)=1}}{}\frac{\mu ^2(s)}{\varphi (s)}\frac{N}{q}\underset{\mathrm{}q}{}\frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}\underset{\genfrac{}{}{0pt}{}{mq}{m>\mathrm{}}}{}\frac{\mu ^2(m)}{\varphi (m)}\underset{\genfrac{}{}{0pt}{}{{\scriptscriptstyle \frac{R}{m}}<s{\scriptscriptstyle \frac{R}{\mathrm{}}}}{(s,q)=1}}{}\frac{\mu ^2(s)}{\varphi (s)}.$$ $`4.25`$ For the first term of (4.25), we observe that $$\underset{mq}{}\frac{\mu ^2(m)}{\varphi (m)}=\frac{q}{\varphi (q)},$$ $`4.26`$ and by Lemma 3 and Lemma 5 we have $`{\displaystyle \frac{N}{q}}`$ $`{\displaystyle \underset{\mathrm{}q}{}}{\displaystyle \frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}}{\displaystyle \underset{mq}{}}{\displaystyle \frac{\mu ^2(m)}{\varphi (m)}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{s{\scriptscriptstyle \frac{R}{\mathrm{}}}}{(s,q)=1}}{}}{\displaystyle \frac{\mu ^2(s)}{\varphi (s)}}`$ $`4.27`$ $`={\displaystyle \frac{N}{\varphi (q)}}(\mathrm{log}R+c+v(q)){\displaystyle \frac{N}{q}}{\displaystyle \underset{\mathrm{}q}{}}{\displaystyle \frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}}\mathrm{log}\mathrm{}+O({\displaystyle \frac{Nw(q)}{\varphi (q)\sqrt{R}}}{\displaystyle \underset{\mathrm{}q}{}}{\displaystyle \frac{\mu ^2(\mathrm{})\sqrt{\mathrm{}}}{\varphi (\mathrm{})}})`$ $`={\displaystyle \frac{N}{\varphi (q)}}(\mathrm{log}R+c)+O({\displaystyle \frac{N}{\varphi (q)\sqrt{R}}}{\displaystyle \underset{pq}{}}(1+{\displaystyle \frac{1}{\sqrt{p}}})(1+{\displaystyle \frac{\sqrt{p}}{p1}})).`$ The last product has logarithm $$\underset{pq}{}\mathrm{log}(1+\frac{1}{\sqrt{p}})+\mathrm{log}(1+\frac{\sqrt{p}}{p1})2\underset{pq}{}\frac{1}{\sqrt{p}}+O(1)\frac{\sqrt{\mathrm{log}q}}{\mathrm{log}\mathrm{log}3q},$$ $`4.28`$ by (3.6). Hence the first term of (4.25) is $$\frac{N}{\varphi (q)}(\mathrm{log}R+c)+O(\frac{N}{\varphi (q)R^{\frac{1}{2}ϵ}})$$ $`4.29`$ for any arbitrarily small and fixed $`ϵ>0`$. Using $$\underset{\genfrac{}{}{0pt}{}{{\scriptscriptstyle \frac{R}{m}}<s{\scriptscriptstyle \frac{R}{\mathrm{}}}}{(s,q)=1}}{}\frac{\mu ^2(s)}{\varphi (s)}=\frac{\varphi (q)}{q}\mathrm{log}\frac{m}{\mathrm{}}+O(w(q)\sqrt{\frac{m}{R}}),$$ $`4.30`$ which is implied by Lemma 3, the second term of (4.25) is expressed as $`{\displaystyle \frac{N\varphi (q)}{q^2}}{\displaystyle \underset{\mathrm{}q}{}}{\displaystyle \frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}}`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{mq}{m>\mathrm{}}}{}}{\displaystyle \frac{\mu ^2(m)}{\varphi (m)}}\mathrm{log}m{\displaystyle \frac{N\varphi (q)}{q^2}}{\displaystyle \underset{\mathrm{}q}{}}{\displaystyle \frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}}\mathrm{log}\mathrm{}{\displaystyle \underset{\genfrac{}{}{0pt}{}{mq}{m>\mathrm{}}}{}}{\displaystyle \frac{\mu ^2(m)}{\varphi (m)}}`$ $`4.31`$ $`+O({\displaystyle \frac{Nw(q)}{q\sqrt{R}}}{\displaystyle \underset{\mathrm{}q}{}}{\displaystyle \frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{mq}{m>\mathrm{}}}{}}{\displaystyle \frac{\mu ^2(m)}{\varphi (m)}}\sqrt{m}).`$ Each term of (4.31) is majorized by deleting the restriction $`m>\mathrm{}`$. Then, by (4.26) and (3.8), the first two terms are each $`{\displaystyle \frac{Nv(q)}{\varphi (q)}}`$, and the error term is the same as that of (4.27). Hence we have shown (4.11). The $`j0`$ case: Since $$\underset{\genfrac{}{}{0pt}{}{s=1}{(s,q)=1}}{\overset{\mathrm{}}{}}\frac{\mu (s)}{\varphi ^2(s)}\mu ((s,j))\varphi ((s,j))=\underset{\genfrac{}{}{0pt}{}{pj}{pq}}{}(1+\frac{1}{p1})\underset{pjq}{}(1\frac{1}{(p1)^2})=\text{S}(jq)\frac{\varphi (q)}{q},$$ $`4.32`$ by (3.10) and (3.11), (4.23) is $$\frac{N}{\varphi (q)}\text{S}(jq)+O(\frac{N}{q}\underset{\mathrm{}q}{}\frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}\underset{mq}{}\frac{\mu ^2(m)}{\varphi (m)}\underset{s>\mathrm{min}(\frac{R}{\mathrm{}},\frac{R}{m})}{}\frac{\mu ^2(s)\mu ^2((s,j))\varphi ((s,j))}{\varphi ^2(s)}).$$ $`4.33`$ The sum over $`s`$ was encountered before in (3.13) and majorized as in (3.9), so the $`O`$-term in (4.33) is $`{\displaystyle \frac{N}{qR}}{\displaystyle \frac{jd(j)}{\varphi (j)}}{\displaystyle \underset{\mathrm{}q}{}}{\displaystyle \frac{\mu ^2(\mathrm{})}{\varphi (\mathrm{})}}{\displaystyle \underset{mq}{}}{\displaystyle \frac{\mu ^2(m)}{\varphi (m)}}\mathrm{max}(\mathrm{},m)`$ $`4.34`$ $`{\displaystyle \frac{Ng(q)}{\varphi (q)R}}{\displaystyle \frac{jd(j)}{\varphi (j)}}.`$ This completes the proof of Proposition 4. We now return to the proof of the theorem. The relevant contributions to (2.7) are $$\underset{\genfrac{}{}{0pt}{}{x<n2x+h}{na\left(modq\right)}}{}\lambda _R^2(n)f(n,x,h)=\frac{xh}{\varphi (q)}(\mathrm{log}R+c+O(v(q))+O(R^{\frac{1}{2}+ϵ}))+O(hR^2),$$ $`4.35`$ and $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{0<kh}{qk}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{x<n2x+hk}{na\left(modq\right)}}{}}\lambda _R(n)\lambda _R(n+k)f(n,x,hk)`$ $`4.36`$ $`={\displaystyle \frac{x}{\varphi (q)}}{\displaystyle \underset{0<j\frac{h}{q}}{}}\text{S}(jq)+O({\displaystyle \frac{xh}{R\varphi (q)}}g(q){\displaystyle \underset{j\frac{h}{q}}{}}{\displaystyle \frac{jd(j)}{\varphi (j)}})`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \frac{xh^2}{\varphi ^2(q)}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{xh}{\varphi (q)}}\mathrm{log}{\displaystyle \frac{h}{q}}+O({\displaystyle \frac{xh}{\varphi (q)}}(\mathrm{log}\mathrm{log}3q)^3)+O({\displaystyle \frac{xh^2g(q)}{Rq\varphi (q)}}\mathrm{log}{\displaystyle \frac{2h}{q}})+O({\displaystyle \frac{h^2R^2}{q}}).`$ Now we put together equations (2.2), (2.5), (2.6), (2.7), (4.3), (4.5), (4.35), (4.36), subject to (1.3), (4.2) and $`qR^2x`$, to obtain $`I(x,h,q,a)`$ $`{\displaystyle \frac{xh}{\varphi (q)}}\mathrm{log}{\displaystyle \frac{Rq}{h}}+O({\displaystyle \frac{xh}{\varphi (q)}}(\mathrm{log}\mathrm{log}3q)^3)+O({\displaystyle \frac{xh^2}{R\varphi (q)}}({\displaystyle \frac{d(q)}{\varphi (q)}}+{\displaystyle \frac{g(q)}{q}}))`$ $`4.37`$ $`+O({\displaystyle \frac{x^{\frac{1}{2}}h^{\frac{3}{2}}R\mathrm{log}^2x}{q^{\frac{1}{2}}}})+O({\displaystyle \frac{x^{\frac{1}{2}}h^2R^{\frac{1}{2}}\mathrm{log}^2x}{\varphi (q)}})+O({\displaystyle \frac{h^2R^2}{q}}).`$ Recall that $`d(q)`$ and $`g(q)`$ are both $`q^ϵ`$. Here we pick $$R=(\frac{x}{hq\mathrm{log}^4x})^{\frac{1}{2}}.$$ $`4.38`$ This choice of $`R`$ makes all the error terms $`o({\displaystyle \frac{xh}{q}})`$ provided that $`h(xq)^{\frac{1}{3}ϵ}`$ .
warning/0004/hep-lat0004025.html	ar5iv	text	# Visualizations of the QCD Vacuum ## 1 Overview The numerical approach to resolving the features of QCD is occasionally criticized via statements<sup>b</sup><sup>b</sup>bSuch statements overlook the fact that one can map out the properties of QCD by exploring the quark-mass and temperature dependence of observables and by resolving the individual quark sector contributions to various observables. suggesting the knowledge of QCD lies “in the silicon.” What these speakers are recognizing is that there are massive amounts of data processed by today’s supercomputers in arriving at the final few bytes of information reported as QCD observables. The focus of this investigation is to further probe the features of QCD by using visualization techniques which efficiently convey the content of these massive amounts of data. Here we will examine the properties of typical pure-gauge QCD-vacuum field configurations created on a $`24^3\times 36`$ space-time lattice using the standard Wilson action at $`\beta =6.0`$, which provides a lattice spacing of 0.1 fm. In addition, we consider an $`𝒪(a^2)`$-improved action on a $`16^3\times 32`$ lattice at $`\beta =4.38`$ providing a lattice spacing of 0.17 fm. To remove the short-range noisy perturbative fluctuations, the field configurations are smoothed by a local algorithm designed to minimize the gauge action at each update. This algorithm known as cooling $`^\mathrm{?}`$ is a well established method for locally suppressing quantum fluctuations. The locality of the method allows topologically nontrivial field configurations to survive numerous iterations of the cooling algorithm. Visualization techniques may be used to answer questions such as the following: To what extent do the QCD vacuum field configurations resemble a randomly oriented multi-(anti)instanton field configuration. Do the instantons display a remnant of spherical symmetry expected for isolated instantons, or do the nonperturbative interactions completely distort this picture? While it has been established that anti-instanton instanton pairs have an attractive interaction and will annihilate during cooling, it is not clear how this process takes place. Do the pairs simply come together or do they wrap around each other in the process of annihilation? One can also search for visual evidence of a polarization phenomena $`^\mathrm{?}`$ where large sized instantons tend to have on average the same sign and are over screened by smaller instantons which tend to have the opposite topological charge of the larger instantons. Visualization techniques have also been used to examine the microscopic effects of various smoothing algorithms.$`^\mathrm{?}`$ ## 2 Visualizations Fig. 1 illustrates a three-dimensional slice of the QCD action density from the $`24^3\times 36`$ lattice after 30 cooling sweeps.<sup>c</sup><sup>c</sup>cEvery link on the lattice is updated once during a single sweep. Here the blue isosurface connects all points having the same action density. Tri-linear interpolation is used to smooth the surface. Volume rendering is done within the surface to illustrate changes in the action density. Outside the surface at low action densities, no volume rendering is done in order to allow one to see within the field configuration. Sharp peaks in the action density have been clipped to aid in the illustration. The regions in which the action density has survived have non-trivial topological charge density as illustrated in Fig. 2. Here the yellow (green) isosurface connects points having equal positive (negative) topological charge density. Volume rendering of the topological charge density inside the isosurfaces illustrates the changes in the positive (negative) topological charge densities by red (blue) extremes in colour. Here we see that the (anti)instantons that are somewhat isolated do show a remnant of spherical symmetry in the action and topological charge densities. Indeed similar plots in which the isosurface explores very high action or topological charge density show an elliptically shaped isosurface. However, closely spaced (anti)instantons can show extreme deviations from spherical symmetry. The action and topological charge densities illustrated in the upper center of Figs. 1 and 2 correspond to instanton anti-instanton pairs in the process of annihilation. By 100 iterations of the cooling algorithm, the action and topological charge densities have largely vanished from this region. In the process of annihilation, one can see regions in which fingers of significant topological charge density are wrapping around negative topological charge density. Between these isosurfaces, the topological charge density is rapidly going to zero and therefore is not rendered. The action density simply shows a region of significant interaction. Animations (in animated-gif format suitable for viewing within a browser) illustrating this process of annihilation are available on the web at http://www.physics.adelaide.edu.au/theory/staff/leinweber/VisualQCD/QCDvacuum/ The process of cooling smoothes and extends the size of (anti)instantons such that the action and charge densities decrease. If one produces an animation with a fixed isosurface value one will see the rendered regions shrink in size and eventually disappear as the local density drops below the isosurface threshold. One would incorrectly reach the conclusion that instantons shrink during cooling and eventually fall through the lattice. These animations illustrate the local action and charge densities relative to the total action and topological charge to more correctly illustrate the distribution of action and charge on the lattice. ## 3 Over-Improved Cooling It is well known that errors in the standard Wilson action eventually destroy (anti)instanton configurations. Errors in the standard Wilson action underestimate the action of a field configuration such that application of a cooling sweep on a single instanton configuration will result in an action less than the one-instanton bound. The instanton is spoiled by the cooling process and eventually the remaining action will be eliminated via the cooling procedure. Hence we also examine the five-loop over-improved action of De Forcrand et al. $`^\mathrm{?}`$ designed to render instantons stable over several hundreds of sweeps. This approach includes extended planar paths combined to reproduce the classical action with no $`𝒪(a^2)`$ nor $`𝒪(a^4)`$ discretization errors and with coefficients fine-tuned to stabilize instantons over many applications of the algorithm. To investigate these algorithms based on improved actions, we consider pure-gauge configurations from an $`𝒪(a^2)`$-improved $`16^3\times 32`$ lattice with a lattice spacing of 0.17 fm. Fig. 3 illustrates a three-dimensional slice of the action density after 30 sweeps of cooling using the standard one-loop Wilson action. The population of (anti)instantons is very sparse and the isolated instantons are quite spherical. Fig. 4 illustrates the same gauge-field configuration, this time cooled with a parallel implementation $`^{\mathrm{?},\mathrm{?}}`$ of the five-loop over-improved action. While there is a correspondence between the instantons surviving 30 sweeps of the one-loop action and those of the five-loop action, it is clear that many more (anti)instantons have survived the improved cooling algorithm. Hence one can see the difficulty in cooling with the standard one-loop Wilson action. The density and size of (anti)instantons is dependent on the number of cooling sweeps.<sup>d</sup><sup>d</sup>dWhile these facts are fairly well known, we find the visualizations illustrating these difficulties very compelling. The true vacuum is much denser as suggested by the five-loop action density. In short, the visualizations reveal rich structure in typical field configurations of the QCD vacuum. The action and topological charge densities presented here display long-range non-perturbative correlations between strongly interacting instantons and anti-instantons. ## Acknowledgments Thanks to John Ahern, Sundance Bilson-Thompson, Frederic Bonnet, Patrick Fitzhenry, Greg Kilcup, Mark Stanford, Tony Williams and James Zanotti for their contributions to making these visualizations possible. Additional thanks to Francis Vaughan of the SACPC for generous supercomputer support and the DHPC Group for support in the development of parallel algorithms. This research is supported by the Australian Research Council. ## References
warning/0004/cond-mat0004475.html	ar5iv	text	# Optical Studies of a Layered Manganite La1.2Sr1.8Mn2O7 : Polaron Correlation Effect \[ ## Abstract Optical conductivity spectra of a cleaved ab-plane of a La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> single crystal exhibit a small polaron absorption band in the mid-infrared region at overall temperatures. With decreasing temperature ($`T`$) to Curie temperature ($`T_C`$), the center frequency of the small polaron band moves to a higher frequency, resulting in a gap-like feature, and that it collapses to a lower frequency below $`T_C`$. Interestingly, with decreasing $`T`$, the stretching phonon mode hardens above $`T_C`$ and softens below $`T_C`$. These concurring changes of lattice and electronic structure indicate that short range polaron correlation exist above $`T_C`$ but disappear with a magnetic ordering. \] Recent studies on manganites have shown that there exist strong coupling among spin, charge, orbital, and lattice degrees of freedom. The relative coupling strength of those degrees of freedom can be sensitively affected by variation of physical parameters, such as amounts of carrier doping, and/or structural modification. For example, the structure of cubic perovskite (La,Sr)MnO<sub>3</sub> can be modified into a layered one by inserting a rock-salt-type block layer (La,Sr)<sub>2</sub>O<sub>2</sub> into every n-MnO<sub>2</sub> sheets, i.e., by forming the Ruddlesden-Popper compound, (La,Sr)<sub>n+1</sub>Mn<sub>n</sub>O<sub>3n+1</sub>. With the variation of structures from single$``$ (n=1: K<sub>2</sub>NiF<sub>4</sub> structure), double$``$ (n=2) and $`\mathrm{}`$ (cubic perovskite) MnO<sub>2</sub> sheet, physical properties of these materials are sensitively varying. In addition, the effective low dimensionality of the reduced n system can enhance charge and spin fluctuations to induce more localized tendency than the cubic one. La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub>, which has the double MnO<sub>2</sub> sheets, is a prototypical material that exibits intriguing interplays of various degrees of freedom. Although it becomes a ferromagnetic (FM) metal below Curie temperature ($`T_C`$) at $`126`$ K, earlier studies showed significant local Jahn-Teller (J-T) lattice distortions at overall temperatures and short range antiferromagnetic spin order. A more recent study also provided a clear evidence of lattice polaron formation above $`T_C`$ by showing diffuse X-ray scattering around the Bragg peaks. At the same time, the scattering experiments indicated an existence of short range polaron ordering by showing incommesurate satellite peaks. Optical spectra for this bilayered system have been reported already. However, there are no systematic optical investigations how polaron effects with short range correlation become manifest in the optical spectra. In this report, we present detailed optical conductivity spectra which reveal interplays of spin, charge, and lattice degrees of freedom in La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub>. With decreasing $`T`$, the mid-infrared small polaron band moves to a higher frequency up to $`T_C`$ and then becomes collapsed to a lower frequency below $`T_C`$. And, the stretching phonon mode hardens above $`T_C`$ and softens below $`T_C`$. These concurring changes of lattice and electronic structure support the existence of the enhanced polaron (charge) correlation above $`T_C`$ and its sudden collapse below $`T_C`$. A single crystal of La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> was grown by the floating-zone method using a mirror furnace. The sample was characterized by resistivity and magnetization measurements. For optical reflectivity measurements, a cleaved ab-plane was freshly prepared. Details for the reflectivity measurements were described in our previous report. Using the Kramers-Kronig relation, we obtained optical conductivity spectra $`\sigma (\omega )`$ from reflectivity spectra $`R(\omega )`$. Figure 1 shows $`T`$-dependent $`R(\omega )`$ of the cleaved ab-plane of La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub>. At 290 K, there are three sharp peaks originating from optic phonon modes in the far-infrared region. With $`T`$ approaching $`T_C`$, $`R(\omega )`$ below 0.4 eV decrease, which is consistent with the dc resistivity behavior shown in the inset of Fig. 1. As $`T`$ decreases further below $`T_C`$, $`R(\omega )`$ start to increase significantly, apporaching to a metallic response. $`T`$-dependent $`\sigma (\omega )`$ both above and below $`T_C`$ are displayed in Figs. 2 (a) and (b), respectively. Above $`T_C`$, $`\sigma (\omega )`$ show a broad absorption band around 1.2 eV. The shape of $`\sigma (\omega )`$ at 290 K looks similar to that at 250 K. However, as $`T`$ approaches 130 K, $`\sigma (\omega )`$ below 0.5 eV become suppressed to show a finite gap-like feature. At the same time, the height of the broad band near 1.2 eV increases to form a sharper band. On the other hand, as $`T`$ is lowered below $`T_C`$, the spectral weight moves suddenly to a lower energy region as shown in Fig. 2 (b). The shape becomes rather asymmetric and its magnitude below 1.0 eV increases significantly, indicating that a large spectral weight become transferred from a higher energy region with the onset of the magnetic ordering. To get further insights, we compared the $`T`$-dependent $`\sigma (\omega )`$ of La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> with those of Nd<sub>0.7</sub>Sr<sub>0.3</sub>MnO<sub>3</sub> (NSMO), which is a cubic perovskite with a similar metal-insulator transition. Compared with the other cubic perovskite manganites such as La<sub>0.7</sub>Ca<sub>0.3</sub>MnO<sub>3</sub> and La<sub>0.7</sub>Sr<sub>0.3</sub>MnO<sub>3</sub>, NSMO has a relatively low $`T_C`$ around $`200`$ K, due to a reduced electron bandwidth. It is now well established that small polaron plays an important role in the paramagnetic insulating regime of the perovskite manganites. Furthermore, a recent X-ray and neutron scattering studies on La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> showed the existence of polarons in the paramagnetic state. Therefore, it is quite reasonable that optical spectra of La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> above $`T_C`$ can be analyzed by the polaron picture. There exist some differences in the spectra of La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> and NSMO. First, $`\sigma (\omega )`$ of NSMO and other cubic perovskites were reported to be nearly $`T`$-independent above $`T_C`$. However, $`\sigma (\omega )`$ of La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> below 2 eV show a systematic $`T`$-dependence above $`T_C`$. As $`T`$ approaches $`T_C`$, spectral weight below 1.0 eV decreases, but that above 1.0 eV increases. Second, $`\sigma (\omega )`$ of La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> do not show a Drude-like peak even at 10 K (see the inset of Fig. 2 (b)). In the FM metallic states of the cubic perovskite manganites, the small polaron spectral weight was transferred to a lower energy to form an asymmetric mid-infrared band and a finite Drude-like peak at very low T, which were interpreted as incoherent and coherent absorption bands of a large polaron, respectively. The lack of the Drude peak in the $`\sigma (\omega )`$ of La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> can be related to its effective low dimensionality, induced by a decrease of the number of the MnO<sub>2</sub> sheet (i.e. $`n=2`$). Third, the lineshape of the La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> absorption band at 10 K $`(0.1T_C`$, $`T_C=126`$ K$`)`$ is quite similar to that of NSMO at 180 K $`(0.9T_C`$, $`T_C=198`$ K$`)`$, as shown in Fig. 2 (b). Note that $`\sigma (\omega )`$ of NSMO at 180 K is close to that of 200 K (above $`T_C`$) in its shape, without showing the Drude-like peak. These observations indicate that La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> remain as the small polaron state even far below $`T_C`$. Therefore, it is reasonable that the $`\sigma (\omega )`$ of La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> can be analyzed in the small polaron picture at all temperatures. To get a more quantitative information on the small polaron absorption, we analyzed the experimental $`\sigma (\omega )`$ in terms of $`\sigma (\omega )=\sigma _{ms}(\omega )+\sigma _L(\omega )`$. Here, $`\sigma _{ms}(\omega )`$ represent the conductivity contribution of the two mid-gap states below 2.0 eV and $`\sigma _L(\omega )`$ correspond to the charge-transfer transition between $`O`$ 2p and Mn $`e_g`$ levels, centered around 4.0 eV. By fitting with a Lorentz oscillator, we determined $`\sigma _L(\omega )`$ and subtracted it from the experimental $`\sigma (\omega )`$ to obtain $`\sigma _{ms}(\omega )`$. For fitting $`\sigma _{ms}(\omega )`$, we used two Gaussian functions as shown in Fig. 3. One is located below 1.0 eV (Peak I) and the other is around 1.5 eV (Peak II). Peak I corresponds to the small polaron absorption related to a nearest neighbor hopping from Mn<sup>3+</sup> to Mn<sup>4+</sup>, and Peak II corresponds to on-site $`dd`$ transition. Interestingly, $`\sigma _{ms}(\omega )`$ above and just below $`T_C`$ can be well described when only the center frequency of Peak I, $`\omega _\text{I}`$, is assumed to be $`T`$ -dependent, while the other parameters such as the strength and the width of Peak I are fixed. And, Peak II are nearly $`T`$ -independent within 3 %. Fig. 3 shows the fitting results above and just below $`T_C`$. With lowering $`T`$ further, the best fitting was obtained when the strength of Peak I as well as $`\omega _\text{I}`$ was assumed to change with a slight variation of the strength of Peak II. Fig. 4 (a) shows $`T`$-dependence of $`\omega _\text{I}`$ obtained by the fitting process. As $`T`$ becomes lower in the paramagnetic region, $`\omega _\text{I}`$ clearly increases from 0.8 to 1.0 eV. With magnetic ordering at $`T_C`$, $`\omega _\text{I}`$ starts to decrease abruptly to reach a finite value of 0.58 eV at 10 K. In case of an adiabatic small polaron, $`\omega _\text{I}`$ corresponds to two times of the small polaron binding energy. Therefore, above results indicate that the coupling between charge and lattice should exist far above $`T_C`$ and that its strength be enhanced near $`T_C`$. And, the coupling strength suddenly decreases to a lower value with the influence of spin ordering. The increase of $`\omega _\text{I}`$ at the high $`T`$ region is responsible for the apparent suppression of $`\sigma (\omega )`$ below 0.4 eV, shown in Fig. 2 (a). The suppression of the $`\sigma (\omega )`$ produces a finite gap-like tail below 0.4 eV. The tail moves systematically to a higher energy between 250 K and $`T_C`$ and the gap-like behavior becomes evident around 130 K just above $`T_C`$. This behavior is reminiscent of a finite charge gap formation in the materials with a clear long range charge ordering (CO) at low $`T`$. Because there is no evidence for the long range CO in La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub>, the peculiar behavior may suggest an existence of short range charge correlation above $`T_C`$. In reality, a recent scattering experiment confirmed the existence of this short range charge and polaron correlation that grows up near $`T_C`$ and diminishs below $`T_C`$. To quantify the short range charge order above $`T_C`$, we define $`\mathrm{\Delta }^{}`$ as a crossing point energy with the $`\sigma (\omega )=0`$ line when a steeply increasing part of $`\sigma (\omega )`$ is linearly extrapolated. Fig. 4 (a) summarizes the $`T`$-dependent values of $`\mathrm{\Delta }^{}`$. The $`T`$-dependence of $`\mathrm{\Delta }^{}`$ is quite similar to that of $`\omega _\text{I}`$. With decreasing $`T`$ to $`T_C`$, $`\mathrm{\Delta }^{}`$ gradually increases from 0.15 to 0.28 eV. The $`\mathrm{\Delta }^{}`$ suddenly starts to decrease at $`T_C`$ and becomes zero below $``$100 K. These experimental results of $`\omega _\text{I}`$ and $`\mathrm{\Delta }^{}`$ strongly support that polaron and charge correlations grow up to $`T_C`$ and collapse due to the FM ordering. The stretching optical phonon mode, related to the lattice degree of freedom, also reflects the existence of short range charge and polaron correlations above $`T_C`$. Figure 5 presents $`T`$-dependence of the phonon mode $`\omega _{\text{TO}}`$ located around 615 cm<sup>-1</sup> at 290 K. With lowering $`T`$ to $`T_C`$, $`\omega _{\text{TO}}`$ shows a significant hardening. When $`T`$ is further lowered, the phonon mode softens. Figure 4 (b) shows the values of $`\omega _{\text{TO}}`$ determined by fitting with the Lorentz oscillator. Between 290 and 130 K, $`\omega _{\text{TO}}`$ increases by about 10 cm<sup>-1</sup>. With lowering $`T`$, it is clearly shown that $`\omega _{\text{TO}}`$ decreases abruptly to 616 cm<sup>-1</sup> at 10 K. The frequency shift of the stretching mode reflects that there exist significant modulations of local Mn-O bond lengths. In long range CO systems, the stretching phonon mode hardens near CO temperature, $`T_{\text{CO}}`$: the observed frequency shifts were about 15 cm<sup>-1</sup> and 25 cm<sup>-1</sup> in case of La<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub> and Bi<sub>0.18</sub>Ca<sub>0.82</sub>MnO<sub>3</sub>, respectively. The hardening behavior of the stretching mode in La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> above $`T_C`$ is somewhat similar to that observed in La<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub> and Bi<sub>0.18</sub>Ca<sub>0.82</sub>MnO<sub>3</sub> near $`T_{\text{CO}}`$. This observation is also consistent with the results of Fig. 4 (a), showing the existence of charge correlation effects in La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub>. In addition, the abrupt softening of $`\omega _{\text{TO}}`$ below $`T_C`$ should be understood in terms of the melting of the short range spatial correlation influenced by the spin ordering. All our experimental findings suggest that there should be intimate coupling among charge, spin, and lattice degrees of freedom (through polaron) and that they interplays with each other in La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub>. Especially, short range spatial correlation effects of the charge, spin, and lattice degrees of freedom can be used to explain the $`T`$-dependence of $`\sigma (\omega )`$, such as the gap-like behaviors and increase of small polaron peak frequency in the mid-infrared region. \[And, the dynamic fluctuations of those various degrees of freedom can be also important in a similar T window.\] In summary, optical conductivities in La<sub>1.2</sub>Sr<sub>1.8</sub>Mn<sub>2</sub>O<sub>7</sub> indicate that the short range correlation of polarons exist above $`T_C`$, and that the sample remains as small polaron state even at 10 K. Subtle balance and competition among the spin, charge, and lattice degrees of freedom should be considered in understanding optical properties of the layered manganite. We thank to H. K. Lee, Y. S. Lee, and Dr. Y. Chung for useful discussion and helpful experiments. This work was supported by Ministry of Science and Technology through the Nanostructure Technology Project and by the BK-21 Project of the Ministry of Education.
warning/0004/math0004018.html	ar5iv	text	# Symmetry Reduction of Discrete Lagrangian Mechanics on Lie groups ## 1. Introduction Background. This paper continues our development of discrete Lagrangian mechanics on a Lie group introduced in Marsden, Pekarsky, and Shkoller \[MPeS 99\]. In our earlier paper, using the context of the Veselov method for discrete mechanics, discrete analogues of Euler-Poincaré and Lie-Poisson reduction theory (see, for example, Marsden and Ratiu \[MR 99\]) were developed for systems on finite dimensional Lie groups $`G`$ with Lagrangians $`L:TG`$ that are $`G`$-invariant. The resulting discrete equations provide “reduced” numerical algorithms which manifestly preserve the symplectic structure. The manifold $`G\times G`$ is used as the discrete approximation of $`TG`$, and a discrete Lagrangian $`𝕃:G\times G`$ is constructed from a given Lagrangian $`L`$ in such a way that the $`G`$-invariance property is preserved. Reduction by $`G`$ results in a new “variational” principle for the reduced Lagrangian $`\mathrm{}:G(G\times G)/G`$, which then determines the discrete Euler-Poincaré (DEP) equations. Reconstruction of these equations is consistent with the usual Veselov discrete Euler-Lagrange equations developed in \[WM 97, MPS 98\], which are naturally symplectic-momentum algorithms. Furthermore, the solution of the DEP algorithm leads directly to a discrete Lie-Poisson (DLP) algorithm. For example, when $`G=\text{SO}(n)`$, the DEP and DLP algorithms for a particular choice of the discrete Lagrangian $`𝕃`$ are equivalent to the Moser-Veselov \[MoV 91\] scheme for the generalized rigid body. Main Results of this Paper. We show that when a discrete Lagrangian $`𝕃:G\times G`$ is $`G`$-invariant, a Poisson structure on (a subset) of one copy of the Lie group $`G`$ can be defined which governs the corresponding discrete reduced dynamics. The symplectic leaves of this structure become dynamically invariant manifolds which are manifestly preserved under the structure preserving discrete Euler-Poincaré algorithm (see Section 2.1). Moreover, starting with a discrete Euler-Poincaré system on $`G`$ one can readily recover, by means of the Legendre transformation, the corresponding Lie-Poisson Hamilton-Jacobi system on $`𝔤^{}`$ analyzed by Ge and Marsden \[GM 88\]; the relationship between the discrete Euler-Lagrange and discrete Euler-Poincaré equations and the Lie-Poisson Hamilton-Jacobi equations was examined from a different point of view in our companion paper \[MPeS 99\]. We also apply Weinstein’s results on Lagrangian mechanics on groupoids and algebroids \[W 96\] to the setting of regular Lie groups. The groupoid-algebroid setting reveals new and interesting connections between discrete and continuous dynamics. ## 2. Discrete reduction In this section we review the discrete Euler-Poincaré reduction of a Lagrangian system on $`G\times G`$ considered in detail in \[MPeS 99\]. We approximate $`TG`$ by $`G\times G`$ and form a discrete Lagrangian $`𝕃:G\times G`$ from the original Lagrangian $`L:TG`$ by $$𝕃(g_k,g_{k+1})=L(\kappa (g_k,g_{k+1}),\chi (g_k,g_{k+1})),$$ where $`\kappa `$ and $`\chi `$ are functions of $`(g_k,g_{k+1})`$ which approximate the current configuration $`g(t)G`$ and the corresponding velocity $`\dot{g}(t)T_gG`$. We choose discretization schemes for which the discrete Lagrangian $`𝕃`$ inherits the symmetries of the original Lagrangian $`L`$: $`𝕃`$ is $`G`$-invariant on $`G\times G`$ whenever $`L`$ is $`G`$-invariant on $`TG`$. In particular, the induced right (left) lifted action of $`G`$ onto $`TG`$ corresponds to the diagonal right (left) action of $`G`$ on $`G\times G`$. Having specified the discrete Lagrangian, we form the action sum $$𝕊=\underset{k=0}{\overset{N1}{}}𝕃(g_k,g_{k+1}),$$ which approximates the action integral $`S=L𝑑t`$, and obtain the discrete Euler-Lagrange (DEL) equations $$D_2𝕃(g_{k1},g_k)+D_1𝕃(g_k,g_{k+1})=0,$$ (2.1) as well as the discrete symplectic form $`\omega _𝕃`$, given in coordinates on $`G\times G`$ by $$\omega _𝕃=\frac{^2𝕃}{g_k^ig_{k+1}^j}dg_k^idg_{k+1}^j.$$ In (2.1), $`D_1`$ and $`D_2`$ denote derivatives with respect to the first and second argument, respectively. The algorithm (2.1) as well as $`\omega _𝕃`$ are obtained by extremizing the action sum $`𝕊:G^{N+1}`$ with arbitrary variations. Using this variational point of view, it is known that the flow $`𝔽_t`$ of the DEL equations preserves this discrete symplectic structure. This result was obtained using a discrete Legendre transform and a direct computation in \[V 88, V 91, WM 97\] and a proof using the variational structure directly was given in \[MPS 98\]. ###### Remark 2.1. We remark that the discrete symplectic structure $`\omega _𝕃`$ is not globally defined, but rather need only be nondegenerate in a neighborhood of the diagonal $`\mathrm{\Delta }`$ in $`G\times G`$, i.e., whenever $`g_k`$ and $`g_{k+1}`$ are nearby. Section 3 of \[MPS 98\] shows that $`\omega _𝕃`$ arises from the boundary terms of the discrete action sum restricted to the space of solutions of the discrete Euler-Lagrange equations; an implicit function theorem argument relying on the regularity of the discrete Lagrangian $`𝕃`$ is required in order to obtain solutions to the discrete Euler-Lagrange equations, and this regularity need only hold in a neighborhood of the diagonal $`\mathrm{\Delta }G\times G`$. ### 2.1. The discrete Euler-Poincaré algorithm The discrete reduction of a right-invariant system proceeds as follows (see \[MPeS 99\] for details). The case of left invariant systems is similar. Of course, some systems such as the rigid body are left invariant. The induced group action on $`G\times G`$ by an element $`\overline{g}G`$ is simply right multiplication in each component: $$\overline{g}:(g_k,g_{k+1})(g_k\overline{g},g_{k+1}\overline{g}),$$ for all $`g_k,g_{k+1}G.`$ The quotient map is given by $$\pi _d:G\times G(G\times G)/GG,(g_k,g_{k+1})g_kg_{k+1}^1.$$ (2.2) One may alternatively use $`g_{k+1}g_k^1`$ instead of $`g_kg_{k+1}^1`$ as the quotient map; the projection map (2.2) defines the reduced discrete Lagrangian $`\mathrm{}:G`$ for any $`G`$-invariant $`𝕃`$ by $`\mathrm{}\pi _d=𝕃`$, so that $$\mathrm{}(g_kg_{k+1}^1)=𝕃(g_k,g_{k+1}),$$ and the reduced action sum is given by $$s=\underset{k=0}{\overset{N1}{}}\mathrm{}(f_{kk+1}),$$ where $`f_{kk+1}g_kg_{k+1}^1`$ denotes points in the quotient space. A reduction of the DEL equations results in the discrete Euler-Poincaré (DEP) equations $$R_{f_{kk+1}}^{}\mathrm{}^{}(f_{kk+1})L_{f_{k1k}}^{}\mathrm{}^{}(f_{k1k})=0$$ (2.3) for $`k=1,\mathrm{},N1`$, where $`R_f^{}`$ and $`L_f^{}`$, for $`fG`$ are the right and left pull-backs by $`f`$, respectively, defined as follows: for $`\alpha _gT_g^{}G`$, $`R_f^{}\alpha _g𝔤^{}`$ is given by $`R_f^{}\alpha _g,\xi =\alpha _g,TR_f\xi `$ for any $`\xi 𝔤`$, where $`TR_f`$ is the tangent map of the right translation map $`R_f:GG`$; $`hhf`$, with a similar definition for $`L_f^{}`$. Also, $`\mathrm{}^{}:GT^{}G`$ is the differential of $`\mathrm{}`$ defined as follows. Let $`g^ϵ`$ be a smooth curve in $`G`$ such that $`g^0=g`$ and $`(d/dϵ)\|_{ϵ=0}g^ϵ=v`$. Then $$\mathrm{}^{}(g)v=(d/dϵ)\|_{ϵ=0}\mathrm{}(g^ϵ).$$ For the other choice of the quotient in (2.2) given by $`h_{k+1k}g_{k+1}g_k^1`$, the DEP equations are $$L_{h_{k+1k}}^{}\mathrm{}^{}(h_{k+1k})R_{h_{kk1}}^{}\mathrm{}^{}(h_{kk1})=0$$ (2.4) ###### Remark 2.2. In the case that $`𝕃`$ is left invariant, the discrete Euler-Poincaré equations take the form $$L_{f_{k+1k}}^{}\mathrm{}^{}(f_{k+1k})R_{f_{kk1}}^{}\mathrm{}^{}(f_{kk1})=0$$ (2.5) where $`f_{k+1k}g_{k+1}^1g_k`$ is in the left quotient $`(G\times G)/G`$. Notice that equations (2.4) and (2.5) are formally the same. We may associate to any $`C^1`$ function $`F`$ defined on a neighborhood $`𝒱`$ of $`\mathrm{\Delta }G\times G`$ its Hamiltonian vector field $`X_F`$ on $`𝒱\mathrm{\Delta }`$ satisfying $`X_F\text{ }\text{ }\omega _𝕃=dF`$, where $`dF`$, the differential of $`F`$, is a one-form. The symplectic structure $`\omega _𝕃`$ naturally defines a Poisson structure on a neighborhood $`𝒱`$ of $`\mathrm{\Delta }`$ (which we shall denote $`\{,\}_{G\times G}`$) by the usual relation $$\{F,H\}_{G\times G}=\omega _𝕃(X_F,X_H).$$ Theorem $`2.2`$ of \[MPeS 99\] states that if the action of $`G`$ on $`G\times G`$ is proper, the algorithm on $`G`$ defined by the discrete Euler-Poincaré equations (2.3) preserves the induced Poisson structure $`\{,\}_G`$ on $`𝒰G`$ given by $$\{f,h\}_G\pi _d=\{f\pi _d,h\pi _d\}_{G\times G}$$ (2.6) for any $`C^1`$ functions $`f,h`$ on $`𝒰`$, where $`𝒰=\pi _d(𝒱)`$. Using the definition $`f_{kk+1}=g_kg_{k+1}^1`$, the DEL algorithm can be reconstructed from the DEP algorithm by $$(g_{k1},g_k)(g_k,g_{k+1})=(f_{k1k}^1g_{k1},f_{kk+1}^1g_k),$$ (2.7) where $`f_{kk+1}`$ is the solution of (2.3). Indeed, $`f_{kk+1}^1g_k`$ is precisely $`g_{k+1}`$. Similarly one shows that in the case of a left $`G`$ action, the reconstruction of the DEP equations (2.5) is given by $$(g_{k1},g_k)(g_k,g_{k+1})=(g_{k1}f_{kk1}^1,g_kf_{k+1k}^1).$$ ### 2.2. The discrete Lie-Poisson algorithm In addition to reconstructing the dynamics on $`G\times G`$, one may use the coadjoint action to form a discrete Lie-Poisson algorithm approximating the dynamics on $`𝔤^{}`$ \[MPeS 99\] $$\mu _{k+1}=\mathrm{Ad}_{f_{kk+1}}^{}\mu _k,$$ (2.8) where $`\mu _k:=\mathrm{Ad}_{g_k^1}^{}\mu _0`$ is an element of the dual of the Lie algebra, $`\mu _0`$ is the constant of motion (the momentum map value), and the sequence $`\{f_{kk+1}\}`$ is provided by the DEP algorithm on $`G`$. The corresponding discrete Lie-Poisson equation for the left invariant system is given by $$\mathrm{\Pi }_{k+1}=\mathrm{Ad}_{f_{k+1k}^1}^{}\mathrm{\Pi }_k,$$ (2.9) where $`\mathrm{\Pi }_k:=\mathrm{Ad}_{g_k}^{}\pi _0`$ and $`\pi _0`$ is the constant momentum map value. Henceforth, we shall use the notation $`\mu 𝔤^{}`$ for the right invariant system and $`\mathrm{\Pi }𝔤^{}`$ for the left. ## 3. Poisson structure and invariant manifolds on Lie groups Discretization of an Euler-Poincaré system on $`TG`$ results in a system on $`G\times G`$ defined by a Lagrangian $`𝕃`$. If it is regular, the Legendre transformation (in the sense of Veselov) $`F𝕃`$ define a symplectic form (and, hence, a Poisson structure) on $`𝒱G\times G`$ via the pull-back of the canonical form from $`T^{}G`$. Then, general Poisson reduction applied to these discrete settings defines a Poisson structure on the reduced space $`𝒰=\pi _d(𝒱)G`$. This approach was adopted in Theorem 2.2 of \[MPeS 99\]. Alternatively, without appealing to the reduction procedure, a Poisson structure on a Lie group can be defined using ideas of Weinstein \[W 96\] on Lagrangian mechanics on groupoids and their algebroids. The key idea can be summarized in the following statements. A smooth function on a groupoid defines a natural (Legendre type) transformation between the groupoid and the dual of its algebroid. This transformation can be used to pull back a canonical Poisson structure from the dual of the algebroid, provided the regularity conditions are satisfied. The ideas outlined in this section can be easily expressed using the groupoid-algebroid formalism. Such a formalism is suited to the discrete gauge field theory generalization as well as to discrete semidirect product theory; nevertheless, the theory of groupoids and algebroids is not essential for the derivations, but rather contributes nicely to the elegance of the exposition. ### 3.1. Dynamics on groupoids and algebroids In this subsection, we show that our discrete reduction methodology is consistent with Weinstein’s groupoid-algebroid construction; the contents of this subsection are not essential for the remainder of the paper. We briefly summarize results from Weinstein \[W 96\] and refer the reader to the original paper for details of proofs and definitions. Let $`\mathrm{\Gamma }`$ be a groupoid over a set $`M`$, with $`\alpha ,\beta :\mathrm{\Gamma }M`$ being its source and target maps, with a multiplication map $`m:\mathrm{\Gamma }_2\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }_2\{(g,h)\mathrm{\Gamma }\times \mathrm{\Gamma }\|\beta (g)=\alpha (h)\}`$. Denote its corresponding algebroid by $`𝒜`$. The Lie groupoids relevant to our exposition are the Cartesian product $`G\times G`$ of a Lie group $`G`$, with multiplication $`(g,h)(h,k)=(g,k)`$, and the group $`G`$ itself. The corresponding algebroids are the tangent bundle $`TG`$ and the Lie algebra $`𝔤`$, respectively. The dual bundle to a Lie algebroid carries a natural Poisson structure. This is the Poisson bracket associated to the canonical symplectic form on $`T^{}G`$ and the Lie-Poisson structure on $`𝔤^{}`$, respectively. Lagrangian mechanics on a groupoid $`\mathrm{\Gamma }`$ is defined as follows. Let $``$ be a smooth, real-valued function on $`\mathrm{\Gamma }`$, $`_2`$ the restriction to $`\mathrm{\Gamma }_2`$ of the function $`(g,h)(g)+(h)`$. ###### Definition 3.1. Let $`\mathrm{\Sigma }_{}\mathrm{\Gamma }_2`$ be the set of critical points of $`_2`$ along the fibers of the multiplication map $`m`$; i.e. the points in $`\mathrm{\Sigma }_{}`$ are stationary points of the function $`(g)+(h)`$ when $`g`$ and $`h`$ are restricted to admissible pairs with the constraint that the product $`gh`$ is fixed \[W 96\]. A solution of the Lagrange equations for the Lagrangian function $``$ is a sequence $`\mathrm{},g_2,g_1,g_0,g_1,g_2,\mathrm{}`$ of elements of $`\mathrm{\Gamma }`$, defined on some “interval” in $``$, such that $`(g_j,g_{j+1})\mathrm{\Sigma }_{}`$ for each $`j`$. The Hamiltonian formalism for discrete Lagrangian systems is based on the fact that each Lagrangian submanifold of a symplectic groupoid determines a Poisson automorphism on the base Poisson manifold. Recall that the cotangent bundle $`T^{}\mathrm{\Gamma }`$ is, in addition to being a symplectic manifold, a groupoid itself, the base being $`𝒜^{}`$; notice that both manifolds are naturally Poisson. The source and target mappings $`\stackrel{~}{\alpha },\stackrel{~}{\beta }:T^{}\mathrm{\Gamma }𝒜^{}`$ are induced by $`\alpha `$ and $`\beta `$. ###### Definition 3.2. Given any smooth function $``$ on $`\mathrm{\Gamma }`$, a Poisson map $`\mathrm{\Lambda }_{}`$ from $`𝒜^{}`$ to itself, which may be said to be generated by $``$ is defined by the Lagrangian submanifold $`d(\mathrm{\Gamma })`$ (under a suitable hypothesis of nondegeneracy) \[W 96\]. The appropriate “Legendre transformation” $`F`$ in the groupoid context is given by $`\stackrel{~}{\alpha }d:\mathrm{\Gamma }𝒜^{}`$ or $`\stackrel{~}{\beta }d:\mathrm{\Gamma }𝒜^{}`$, depending on whether we consider right or left invariance (through the definition of maps $`\stackrel{~}{\alpha }`$ and $`\stackrel{~}{\beta }`$). The transformation $`F`$ relates the mapping on $`\mathrm{\Gamma }`$ defined by $`\mathrm{\Sigma }_{}`$ with the mapping $`\mathrm{\Lambda }_{}`$ on $`𝒜^{}`$. $`F`$ also pulls back the Poisson structure from $`𝒜^{}`$ to $`\mathrm{\Gamma }`$, which, in general, is defined only locally on some neighborhood $`𝒰\mathrm{\Gamma }`$. In the context of a Lie group, this means that any regular function $`\mathrm{}:G`$ defines a Poisson structure on $`𝒰`$. We shall address this issue in the next subsections. The reader is referred to \[W 96\] for an application of the above ideas to the groupoid $`M\times M`$ when the manifold $`M`$ does not necessarily have group structure. ### 3.2. DEP equations as generators of Lie-Poisson Hamilton-Jacobi equations A Lie group $`G`$ is the simplest example of a groupoid with the base being just a point. Its algebroid is the corresponding Lie algebra $`𝔤`$, with the dual being $`𝔤^{}`$. Consider left invariance and let a general function $``$ on the group be specified by the discrete reduced Lagrangian $`\mathrm{}:G`$. Then, the Legendre transform defined above is given by $$F\mathrm{}=L_g^{}d\mathrm{}:G𝔤^{},$$ where $`d\mathrm{}:GT^{}G`$. Using these transformations we define $$\mathrm{\Pi }_{k1}F\mathrm{}(f_{kk1})=L_{f_{kk1}}^{}d\mathrm{}(f_{kk1}).$$ Recall the DEP equation (2.5) for left-invariant systems : $$L_{f_{k+1k}}^{}d\mathrm{}(f_{k+1k})R_{f_{kk1}}^{}d\mathrm{}(f_{kk1})=0,$$ where we have identified the notations $`\mathrm{}^{}`$ and $`d\mathrm{}`$. The latter equation can be rewritten as a system $$\{\begin{array}{c}\mathrm{\Pi }_k=L_f^{}d\mathrm{}(f),\hfill \\ \mathrm{\Pi }_{k+1}=R_f^{}d\mathrm{}(f),\hfill \end{array}$$ (3.1) where the first equation is to be solved for $`f`$ (which stands for $`f_{k+1k}`$) which then is substituted into the second equation to compute $`\mathrm{\Pi }_{k+1}`$. This system is precisely the Lie-Poisson Hamilton-Jacobi system described in \[GM 88\] with the reduced discrete Lagrangian $`\mathrm{}`$ playing the role of the generating function. This means that there is no need to find an approximate solution of the reduced Hamilton-Jacobi equation \[GM 88\]. Notice also that the DLP equation (2.9) is a direct consequence of the system (3.1): $$\mathrm{\Pi }_{k+1}=\mathrm{Ad}_{f_{k+1k}^1}^{}\mathrm{\Pi }_k.$$ The following diagrams relate the dynamics on $`G`$ and on $`𝔤^{}`$: $$\begin{array}{ccc}G& \stackrel{\mathrm{\Sigma }_{\mathrm{}}}{}& G\\ F\mathrm{}& & F\mathrm{}\\ 𝔤^{}& \stackrel{\mathrm{\Lambda }_{\mathrm{}}}{}& 𝔤^{}\end{array}\begin{array}{ccc}f_{kk1}& \stackrel{\mathrm{\Sigma }_{\mathrm{}}}{}& f_{k+1k}\\ F\mathrm{}& & F\mathrm{}\\ \mathrm{\Pi }_{k1}& \stackrel{\mathrm{\Lambda }_{\mathrm{}}}{}& \mathrm{\Pi }_k,\end{array}$$ (3.2) where $`\mathrm{\Sigma }_{\mathrm{}}`$ and $`\mathrm{\Lambda }_{\mathrm{}}`$ are given in Definitions 3.1 and 3.2. ### 3.3. Some Advantages of Structure-preserving Integrators As we mentioned above, the “Legendre transform” $`F\mathrm{}`$ allows us to put a Poisson structure on the Lie group $`G`$, which, of course, depends on the discrete Lagrangian $`𝕃`$ on $`G\times G`$, and hence on the original Lagrangian $`L`$ on $`TG`$ (if we consider this from the discrete reduction point of view). It follows that the reduction of the discrete Euler-Lagrange dynamics on $`G\times G`$ is necessarily restricted to the symplectic leaves of this Poisson structure, so that these leaves are invariant manifolds, and correspond (under $`F\mathrm{}`$) to the symplectic leaves (coadjoint orbits) of the continuous reduced system on $`𝔤^{}`$. These ideas are the content of the following theorems. ###### Theorem 3.1. Let $`L`$ be a right invariant Lagrangian on $`TG`$ and let $`𝕃`$ be the Lagrangian of the corresponding discrete system on $`𝒱G\times G`$. Assume that $`𝕃`$ is regular, in the sense that the Legendre transformation $`F𝕃:𝒱F𝕃(𝒱)T^{}G`$ is a local diffeomorphism, and let the quotient maps be given by $$\pi _d:G\times G(G\times G)/GG\text{and}\pi :T^{}G(T^{}G)/G𝔤^{}.$$ Let $`\mathrm{}`$ be the reduced Lagrangian on $`G`$ defined by $$𝕃=\mathrm{}\pi _d,$$ and let $$F\mathrm{}:𝒰G𝔤^{}$$ be the corresponding Legendre transform. Then the following diagram commutes: $$\begin{array}{ccc}𝒱G\times G& \stackrel{F𝕃}{}& T^{}G\\ \pi _d& & \pi \\ 𝒰G& \stackrel{F\mathrm{}}{}& 𝔤^{}.\end{array}$$ (3.3) ###### Proof. First, we choose coordinate systems on each space. Let $`(g_k,g_{k+1})G\times G`$ and $`(g,p)T^{}G`$, so that the discrete quotient map (2.2) is given by $`\pi _d:(g_k,g_{k+1})f_{kk+1}=g_kg_{k+1}^1`$, and the continuous quotient map by $`\pi :(g,p)\mu =R_g^{}p`$. Recall that the fiber derivative $`F𝕃`$ in these coordinates has the following form (see, e.g., \[WM 97\]) $$F𝕃:G\times GT^{}G;(g_k,g_{k+1})(g_k,D_1𝕃(g_k,g_{k+1})).$$ Then the above diagram is given by: $$\begin{array}{ccc}(g_k,g_{k+1})& \stackrel{F𝕃}{}& (g_k,p_k=\frac{𝕃}{g_k})\\ \pi _d& & \pi \\ f=R_{g_{k+1}^1}g_k& & \mu =R_{g_k}^{}p_k,\end{array}$$ (3.4) where $`f`$ stands for $`f_{kk+1}=g_kg_{k+1}^1`$. To close this diagram and to verify the arrow determined by $`F\mathrm{}`$ compute the derivative of $`𝕃`$ using the chain rule $$\mu =R_{g_k}^{}p_k=R_{g_k}^{}\frac{(\mathrm{}\pi )}{g_k}=R_{g_k}^{}\left(R_{g_{k+1}^1}^{}\frac{\mathrm{}}{f}\right)=R_f^{}\frac{\mathrm{}}{f}=R_f^{}\mathrm{}^{}(f),$$ (3.5) where we have used that, according to the definition of $`f`$, the partial derivative $`{\displaystyle \frac{f}{g_k}}`$ is given by the linear operator $`TR_{g_{k=1}^1}`$. (3.5) is precisely the Legendre transformation $`F\mathrm{}`$ for a right invariant system (see the previous subsection). ∎ ###### Corollary 3.1. Reconstruction of the discrete Lie-Poisson (DLP) dynamics on $`𝔤^{}`$ by $`\pi ^1`$ corresponds to the image of the discrete Euler-Lagrange (DEL) dynamics on $`G\times G`$ under the Legendre transformations $`F𝕃`$ and results in an algorithm on $`T^{}G`$ approximating the continuous flow of the corresponding Hamiltonian system. ###### Proof. The proof follows from the results of the previous subsection, in particular, diagram (3.2) relates the DLP dynamics on $`𝔤^{}`$ with the DEP dynamics on $`𝒰G`$ which, in turn, is related to the DEL dynamics on $`𝒱G\times G`$ via the reconstruction (2.7). ∎ An important remark to this corollary which follows from the results in \[KMO 99\] (see also \[KMOW 99)\]) is that, in general, to get a corresponding algorithm on the Hamiltonian side which is consistent with the corresponding continuous Hamiltonian system on $`T^{}G`$, one must use the time step $`h`$-dependent Legendre transform given by the map $$(g_k,g_{k+1})(g_k,hD_1𝕃(g_k,g_{k+1})).$$ The results of this paper are not effected, however, as we assume $`h`$ to be constant and so we would simply add a constant multiplier to the corresponding symplectic and Poisson structures. For variable time-stepping algorithms, this remark is crucial and must be taken into account. ###### Theorem 3.2. The Poisson structure on the Lie group $`G`$ obtained by reduction of the Lagrange symplectic form $`\omega _𝕃`$ on $`𝒱G\times G`$ via $`\pi _d`$ coincides with the Poisson structure on $`𝒰G`$ obtained by the pull-back of the Lie-Poisson structure $`\omega _\mu `$ on $`𝔤^{}`$ by the Legendre transformation $`F\mathrm{}`$ (see diagram (3.3) above). ###### Proof. The proof is based on the commutativity of diagram (3.3) and the $`G`$ invariance of the unreduced symplectic forms. Notice that $`G`$ and $`𝔤^{}`$ in (3.3) are Poisson manifolds, each being foliated by symplectic leaves, which we denote $`\mathrm{\Sigma }_f`$ and $`𝒪_\mu `$ for $`fG`$ and $`\mu 𝔤^{}`$, respectively. Denote by $`\omega _f`$ and $`\omega _\mu `$ the corresponding symplectic forms on these leaves. Below we shall prove the compatibility of these structures under the diagram (3.3). Repeating this proof leaf by leaf establishes then the equivalence of the Poisson structures and proves the theorem. Recall that the Lagrange $`2`$-form $`\omega _𝕃`$ on $`𝒱G\times G`$ derived from the variational principle coincides with the pull-back of the canonical $`2`$-form $`\omega _{\text{can}}`$ on $`T^{}G`$ (see, e.g., \[MPS 98, WM 97\]). Recall also that for a right-invariant system, reduction of $`T^{}G`$ to $`𝔤^{}`$ is given by right translation to the identity $`eG`$, i.e. any $`pT_g^{}G`$ is mapped to $`\mu =R_g^{}p𝔤^{}T_e^{}G`$. Thus, for any $`g\pi ^1(\mu )`$, where $`\mu 𝔤^{}`$, $$\pi ^1\|_{T^{}G}=R_{g^1}^{}:𝔤^{}T_g^{}G,$$ so that $`(\pi ^1)^{}=(R_{g^1}^{})^{}`$ pulls back $`\omega _{\text{can}}`$ to $`\omega _\mu `$. Henceforth, $`\pi ^1`$ shall denote the inverse map of $`\pi `$ restricted to $`T_gG^{}`$. Let us write down using the above notations how the symplectic forms are being mapped under the transformations in diagram (3.3); we see that $$\begin{array}{ccc}𝒱G\times G& \stackrel{F𝕃}{}& T^{}G\\ \pi _d& & \pi \\ 𝒰G& \stackrel{F\mathrm{}}{}& 𝔤^{}\end{array}\begin{array}{ccc}\omega _𝕃& \stackrel{F𝕃^{}}{}& \omega _{\text{can}}\\ & & \left(\pi ^1\right)^{}\\ \omega _f& \stackrel{F\mathrm{}^{}}{}& \omega _\mu .\end{array}$$ (3.6) Then, using the coordinate notations of diagram (3.4), for any $`fG`$ and $`u,vT_f\mathrm{\Sigma }_f`$, $$\omega _f(f)(u,v)\omega _\mu (\mu )(TF\mathrm{}(u),TF\mathrm{}(v)),$$ (3.7) where $`\mu =F\mathrm{}(f)𝔤^{}`$. Continuing this equation using diagram (3.6), we have that $$\begin{array}{c}\omega _f(f)(u,v)=\omega _{\text{can}}((g_k,p_k))(T\pi ^1TF\mathrm{}(u),T\pi ^1TF\mathrm{}(v))\hfill \\ \hfill =\omega _𝕃((g_k,g_{k+1}))(TF𝕃^1T\pi ^1TF\mathrm{}(u),TF𝕃^1T\pi ^1TF\mathrm{}(v)),\end{array}$$ (3.8) where $`(g_k,p_k)\pi ^1(\mu )`$ and $`T\pi ^1`$ denotes $`TR_{g^1}^{}`$. Using (3.3), it follows that $$F\mathrm{}\pi _d=\pi F𝕃$$ and, hence, for the tangent maps, we have that $$TF\mathrm{}T\pi _d=T\pi TF𝕃.$$ So, if $`u,v`$ in (3.7) are images of some $`G`$ invariant vector fields $`U,V`$ on $`𝒱G\times G`$, i.e. $`u=T\pi _d(U),v=T\pi _d(V)`$, then from (3.8) it follows that $$\omega _f(f)(u,v)=\omega _𝕃((g_k,g_{k+1}))(U,V),$$ where $`(g_k,g_{k+1})=\pi _d^1(f)`$ and $`U,VT_{(g_k,g_{k+1})}G\times G`$. The last equation precisely means that $`\omega _f`$ is the discretely reduced symplectic form, i.e. the image of $`\omega _𝕃`$ under the quotient map $`\pi _d`$. ∎ Analogous theorems hold for the case of left invariant systems. More General Configuration Spaces. Similar ideas carry over to the integration of systems defined on a general configuration space $`M`$ with some symmetry group $`G`$. In this case, the reduced discrete space $`(M\times M)/G`$ inherits a Poisson structure from the one defined on $`M\times M`$ (analogously to (2.6)). Its symplectic leaves again become dynamically invariant manifolds for structure-preserving integrators and can be viewed as images of the symplectic leaves of the reduced Poisson manifold $`T^{}M/G`$ under appropriately defined “Legendre transformations”. This is a topic of ongoing research that builds on recent progress in Lagrangian reduction theory; see \[MRS 99\]. ### 3.4. Poisson structures of the rigid body As an example of applications of the above ideas, we consider the dynamics of the rigid body and its associated reduction and discretization (see, e.g. \[MR 99, MoV 91, LS 96, MPeS 99\] for more details). The Basic Set Up. The configuration space of the system is $`\mathrm{SO}(3)`$. The corresponding Lagrangian is determined by a symmetric positive definite operator $`J:𝔰o(3)𝔰o(3)`$, defined by $`J(\xi )=\mathrm{\Lambda }\xi +\xi \mathrm{\Lambda }`$, where $`\xi 𝔰o(3)`$ and $`\mathrm{\Lambda }`$ is a diagonal matrix satisfying $`\mathrm{\Lambda }_i+\mathrm{\Lambda }_j>0`$ for all $`ij`$. The left invariant metric on $`\mathrm{SO}(3)`$ is obtained by left translating the bilinear form at $`e\mathrm{SO}(3)`$ given by $$(\xi ,\xi )=\frac{1}{4}\mathrm{Tr}\left(\xi ^TJ(\xi )\right).$$ The operator $`J`$, viewed as a mapping $`𝒥:𝔰o(3)𝔰o(3)^{}`$, has the usual interpretation of the inertia tensor, and the $`\mathrm{\Lambda }_i`$ correspond to the sums of certain principal moments of inertia. The rigid body Lagrangian is the kinetic energy of the system $$L(g,\dot{g})=\frac{1}{4}g^1\dot{g},𝒥(g^1\dot{g})=\frac{1}{4}\xi ,𝒥\xi )=l(\xi ),$$ (3.9) where $`\xi =g^1\dot{g}𝔰o(3)`$ and $`,`$ is the pairing between the Lie group and its dual. Poisson Structures and Casimir Functions. The Lie algebra dual $`𝔰o(3)^{}`$ has a well-known Lie-Poisson structure with a Casimir $`C_{𝔰𝔬(3)^{}}(\mu )=\mathrm{Tr}(\mu ^2)`$, where $`\mu 𝔰o(3)^{}`$. Upon identification with $`^3`$, its generic symplectic leaves become concentric spheres with Kirilov-Kostant symplectic form being proportional to the area form. If $`y`$ denotes coordinates on $`^3𝔰o(3)^{}`$, then the above Casimir function is given by $`C_{𝔰o(3)^{}}(y)=y^2`$. Following Section 5 of \[MPeS 99\] on discrete Euler-Poincaré reduction, we obtain the reduced form of the Moser-Veselov Lagrangian on the group SO$`(3)`$ given by $$\mathrm{}(f)=\mathrm{Tr}(f\mathrm{\Lambda }),$$ where $`f\mathrm{SO}(3)`$ and $`\mathrm{SO}(3)`$ is embedded into the linear space $`𝔤𝔩(3)`$. Then, the Legendre transform $`F\mathrm{}`$ takes the form $$F\mathrm{}(f)=L_f^{}d\mathrm{}(f)=\mathrm{skew}(f\mathrm{\Lambda })=f\mathrm{\Lambda }\mathrm{\Lambda }f^T:\mathrm{SO}(3)𝔰o(3)^{},$$ where the constraint that $`f`$ be in $`\mathrm{SO}(3)`$ has been enforced. The pull-back of $`C_{𝔰o(3)^{}}`$ under $`F\mathrm{}^{}`$ defines a Casimir function on the group, which up to a constant term and a sign, is given by $$C_{\mathrm{SO}(3)}(f)=\mathrm{Tr}(f\mathrm{\Lambda }f\mathrm{\Lambda })f\mathrm{SO}(3).$$ (3.10) Its symplectic leaves constitute the invariant manifolds of the reduced discrete dynamics corresponding to the Lagrangian (3.9). Analogously, one can define a Poisson structure on the Lie algebra $`𝔰o(3)`$ using the duality between Lie-Poisson and Euler-Poincaré reduced systems on $`𝔰o(3)^{}`$ and $`𝔰o(3)`$, respectively. The Lagrangian (3.9) defines the Legendre transformations $`Fl`$ from $`𝔰o(3)`$ to $`𝔰o(3)^{}`$ given by $`\mu ={\displaystyle \frac{l}{\xi }}=𝒥(\xi )`$. Then, the pull-back by $`Fl^{}`$ defines a Casimir function on $`𝔰o(3)`$: $$C_{𝔰o(3)}(\xi )=Fl^{}C_{𝔰o(3)^{}}(\xi )=𝒥(\xi ),𝒥(\xi ),$$ where the metric on the dual is induced by the one on the algebra, i.e. by the symmetric positive definite operator $`J`$. If $`x`$ denotes coordinates on $`^3𝔰o(3)`$, then the above Casimir function is given by $`C_{𝔰o(3)}(x)=𝒥(x)^2`$. Thus, the corresponding symplectic leaves are ellipsoids of $`𝒥^2`$. They do not coincide with adjoint orbits, which are spheres in $`^3`$. The dynamic orbits are obtained by intersecting these ellipsoids, determined by $`𝒥^2`$, with the energy ellipsoids, determined by $`𝒥`$. ## Acknowledgments The authors would like to thank Alan Weinstein for pointing out the connections with the general theory of dynamics on groupoids and algebroids and Sameer Jalnapurkar for comments and discussions. SP and SS would also like to thank the Center for Nonlinear Science in Los Alamos for providing a valuable setting for part of this work. SS was partially supported by NSF-KDI grant ATM-98-73133, and JEM and SP were partially supported by NSF-KDI grant ATM-98-73133 and the Air Force Office of Scientific Research.
warning/0004/cond-mat0004032.html	ar5iv	text	# Ab-initio prediction of the electronic and optical excitations in polythiophene: isolated chains versus bulk polymer ## I Introduction Semiconducting conjugated organic polymers have received increasing interest in recent years, especially since the discovery of electroluminescence of these materials. The charge carriers and excitations in these materials have been studied extensively both experimentally and theoretically, but many important fundamental issues still remain unsolved. For instance, the magnitude of the exciton binding energy in these materials is still disputed. This is a very important quantity, since e.g. in photovoltaic devices (solar cells) one would like to have a small binding energy, which facilitates the fast separation of charges, while in electroluminescent devices such as LEDs a larger exciton binding energy, to increase the probability of fast (radiative) annihilation of electron-hole pairs, is desirable. In conventional semiconductors such as Si and GaAs the optical excitations are well described in terms of very weakly bound electron-hole pairs (so-called Wannier excitons) with a binding energy of the order of 0.01 eV. In crystals made of small organic molecules such as anthracene, the exciton is essentially confined to a single molecule (Frenkel exciton), leading to a binding energy of the order of 1 eV. The question is where exactly conjugated polymers fit in between conventional semiconductors on the one hand and molecular crystals on the other: negligibly small (0.1 eV or less ), intermediate ($``$ 0.5 eV ), and large ($``$ 1.0 eV ) binding energies have been proposed. Ab-initio calculations, on a variety of conjugated polymers, within the Local Density Approximation of Density Functional Theory (DFT-LDA) yield equilibrium structures in very good agreement with experiment. The Kohn-Sham gaps in these calculations are typically 40% smaller than the optical band gap (absorption gap). In cases where calculations for both isolated chains and the crystalline situation were performed, small differences (0.1 to 0.2 eV for gaps of $`3.0`$ eV) in Kohn-Sham band gaps were found. However, it is well known that the Kohn-Sham eigenvalues formally cannot be interpreted as excited state energies. Moreover, excitonic effects are not taken into account in these calculations. An ab-initio many-body calculation within the $`GW`$ Approximation ($`GW`$A) was performed for poly-acetylene (PA) by Ethridge et al. They claim that their quasi-particle (QP) gap, excluding excitonic effects, is in agreement with the experimental absorption gap. This result seems to be in contrast with a more recent calculation by Rohlfing and Louie of both one- (QP) and two-particle (exciton) excitation energies for PA and poly-phenylene-vinylene (PPV) chains. Their absorption gaps are in good agreement with experiments, but the inclusion of excitonic effects proves to be crucial for this. However, their exciton binding energy of 0.9 eV for PPV is much larger than recently obtained experimental values: $`0.35\pm 0.15`$ eV for an alkoxy-substituted PPV, and $`0.48\pm 0.14`$ eV for unsubstituted PPV. In a recent Letter, hereafter referred to as I, we focused on the differences in excitations between an isolated polythiophene (PT) chain, see Fig. 1, and crystalline polythiophene. For the isolated chain, we found an absorption gap in good agreement with experiment, but the energy differences between the various exciton levels were too large. After including the screening by the surrounding chains, both the optical gap and the exciton transition energies were in good agreement with the experimental values. The difference in screening between an isolated polymer chain and a condensed polymer medium can be explained as follows. For an isolated quasi-one dimensional system, such as a single polymer chain in vacuum, there is no long-range screening. A way to understand this is to realize that if we have two charges on a polymer chain with a separation larger than the width of the chain ($`7`$ a.u.), most field lines connecting the charges will be outside the chain. If the chain is embedded in a medium (possibly, but not necessarily, consisting of similar chains), the medium will provide a long-range screening of the Coulomb interaction. The screened Coulomb interaction determines both the QP energies (including the band gap) and the exciton energies. Long-range screening reduces both the QP band gap and the exciton binding energies. Apparently, there is a near cancellation between the change of the QP gap and the exciton binding energy, meaning that the optical absorption gap, which is the difference between the two, is influenced much less by introducing screening. In I, a method for the calculation of the dielectric tensor of crystalline polythiophene from the ab-initio single chain polarizability function was introduced, without giving any details. Further, some novel technical procedures in the $`GW`$A calculation were used, in particular in the handling of the Coulomb divergence both in real and reciprocal space. Details of these approaches, as well as of the calculation of the quasi-particle energies and exciton binding energies, will be explained here. The paper is organized as follows: in the next section (II) we explain the computational methods employed to calculate the quasi-particle bandstructure, to regularize Coulomb interaction, to calculate the exciton binding energies and the dielectric tensor. In Section III, we will present results for the electronic and optical excitations of both the isolated chain and bulk PT. In section IV we will discuss these results and compare them to other calculations, and draw our conclusions. ## II Computational methods Many successful ab-initio calculations of the QP band structure of conventional anorganic semiconductors have been performed within the $`GW`$ for the electronic self-energy $`\mathrm{\Sigma }`$ of the one-particle Green function. Very recently, progress has been made in the evaluation of the two-particle Green function, from which the optical properties can be obtained. This is done by solving the Bethe-Salpeter equation (BSE), which can be mapped onto a two-body Schrödinger equation for an electron and a hole forming an exciton. We will use these approaches to calculate the QP band structure and exciton binding energies of PT. The calculational scheme is as follows: first we perform a DFT-LDA-based Car-Parrinello calculation, from which we obtain atomic positions, wave functions and ground-state energies. We use these as input for the $`GW`$A calculation, which yields the QP excitation energies. With the DFT-LDA wave functions and the QP energies we calculate the two-particle excitations by solving the BSE. This scheme is first applied to the single chain and next to the crystal. We assume the same atomic geometry for the ground and excited states, i.e. the coupling between electronic and lattice degrees of freedom is neglected. Experimental data indicate that energy shifts due to lattice relaxations are of the order of 0.1 eV in PT. DFT-LDA calculations predict a hole-polaron relaxation energy of 0.04 eV for 16T ($`n`$T is an oligomer consisting of $`n`$ thiophene-rings). A similar calculation predicts a triplet exciton relaxation energy of 0.2 eV for 12T. Singlet relaxation energies are typically smaller. These values, calculated for oligomers, are upper bounds for the values in the polymer, since in the oligomers the excitation is confined, leading to a larger local deviation from the ground state density and hence to a larger relaxation energy than in the polymer. ### A The quasi-particle equation We start our calculations with a pseudo-potential plane-wave DFT-LDA calculation of a geometry-relaxed PT chain in a tetragonal supercell. The plane wave cut-off energy is 40 Ry. The length in the chain direction $`a_x`$ was optimized and found to be $`a_x=14.80`$ a.u. (experimental values range from $`14.65`$ a.u. to $`15.18`$ a.u.). In the perpendicular directions we found that a separation of $`a_y=a_z=15.0`$ a.u. is enough to consider the chains in the DFT-LDA calculation as non-interacting. The two rings in the unit cell are found to be co-planar and we choose them in the $`y=z`$ plane. We use Hartree atomic units (with the Bohr radius $`a_0`$ as unit of length and the Hartree as unit of energy) throughout this article, unless specified otherwise. The one-particle excitation energies are evaluated by solving the QP equation: $$\left[\frac{^2}{2}+V_H(\text{r})\right]\varphi _{nk}(\text{r})+\left[V_{PP}(\text{r},\text{r}^{})+\mathrm{\Sigma }(\text{r},\text{r}^{},E_{nk})\right]\varphi _{nk}(\text{r}^{})d^3r^{}=E_{nk}\varphi _{nk}(\text{r}),$$ (1) where $`V_H`$ is the Hartree potential, $`V_{PP}`$ the non-local pseudo-potential of the atomic core, and $`\mathrm{\Sigma }`$ the electronic self-energy. Since in practice the DFT-LDA wave functions and the QP wave functions are almost identical, we use the former in all calculations. In DFT-LDA, $`\mathrm{\Sigma }`$ is approximated by: $$\mathrm{\Sigma }(\text{r},\text{r}^{},\omega )=V_{xc}(\text{r})\delta (\text{r}\text{r}^{}),$$ (2) where $`V_{xc}`$ is the exchange-correlation potential of the homogeneous electron gas. In the $`GW`$A, $`\mathrm{\Sigma }`$ is approximated by the first term of the many-body expansion in terms of the one-particle Green function $`G`$ and the screened Coulomb interaction $`W`$ of the system. In order to calculate $`\mathrm{\Sigma }`$, we follow the real-space imaginary-time formulation of the $`GW`$A of Rojas et al., in its mixed-space formulation. In this formulation, we transform non-local functions $`F(\text{r},\text{r}^{})`$ to functions $`F_k(\text{r},\text{r}^{})`$ $`F_k(\text{r},\text{r}^{})`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{n_k}{}}}F(\text{r}+na_x\widehat{x},\text{r}^{})e^{ik(x+na_xx^{})},`$ (3) $`F(\text{r},\text{r}^{})`$ $`=`$ $`{\displaystyle \frac{1}{n_k}}{\displaystyle \underset{k}{}}F_k(\text{r},\text{r}^{})e^{ik(xx^{})},`$ (4) where $`n_k`$ is the number of equidistant $`k`$-points in the 1D Brillouin zone (BZ). We use periodic boundary conditions: $`F(\text{r}+n_ka_x\widehat{x},\text{r}^{})=F(\text{r},\text{r}^{}+n_ka_x\widehat{x})=F(\text{r},\text{r}^{})`$. The functions $`F_k(\text{r},\text{r}^{})`$ are fully periodic: $`F_k(\text{r}+a_x\widehat{x},\text{r}^{})=F_k(\text{r},\text{r}^{}+a_x\widehat{x})=F_k(\text{r},\text{r}^{})`$, so that r and $`\text{r}^{}`$ can be chosen in the unit cell. We calculate the one-particle Green function for imaginary times $$G_k(\text{r},\text{r}^{},i\tau )=\{\begin{array}{ccc}\hfill i& \underset{v}{}u_{vk}(\text{r})u_{vk}^{}(\text{r}^{})e^{(ϵ_{vk}ϵ_F)\tau }\hfill & \text{ for }\tau <0,\hfill \\ & & \\ \hfill i& \underset{c}{}u_{ck}(\text{r})u_{ck}^{}(\text{r}^{})e^{(ϵ_{ck}ϵ_F)\tau }\hfill & \text{ for }\tau >0,\hfill \end{array}$$ (5) where $`u_{nk}(\text{r})=\varphi _{nk}(\text{r})e^{ikx}`$ and $`ϵ_{nk}`$ (with $`n=c,v`$) are DFT-LDA wavefunctions and the corresponding energies (with $`v`$ and $`c`$ referring to valence and conduction states, respectively). $`ϵ_F`$ is the Fermi energy (set in the middle of the DFT-LDA gap). We further calculate the irreducible single-chain polarizability function in the Random Phase Approximation (RPA) $$P_k(\text{r},\text{r}^{},i\tau )=2i\underset{q}{}G_q(\text{r},\text{r}^{},i\tau )G_{qk}(\text{r}^{},\text{r},i\tau ),$$ (6) the screened Coulomb interaction $$W_k(\text{r},\text{r}^{},i\omega )=\left[\stackrel{~}{v}_k^1(\text{r},\text{r}^{})P_k(\text{r},\text{r}^{},i\omega )\right]^1,$$ (7) where $`\stackrel{~}{v}_k(\text{r},\text{r}^{})`$ is a cut-off Coulomb interaction in mixed space, discussed in the next Section, and we calculate the electronic self-energy $$\mathrm{\Sigma }_k(\text{r},\text{r}^{},i\tau )=i\underset{q}{}G_q(\text{r},\text{r}^{},i\tau )W_{kq}(\text{r},\text{r}^{},i\tau ).$$ (8) We calculate all the above two-point functions on a double $`24\times 24\times 24`$ real-space grid for r and $`\text{r}^{}`$ in the unit cell. This corresponds to a plane wave cut-off of 25 Ry. The total number of valence and conduction bands taken into account was 300. In Eqs. (6), (7) and (8) we switch between time- and frequency domain using Fourier transforms. Our imaginary-time grid has an exponential spacing (0.25 a.u. near $`\tau =0`$, up to a spacing 6 a.u. near $`\tau _{max}=32.0\text{ }\mathrm{a}.\mathrm{u}.`$) and we interpolate to a linear grid when using the Fast Fourier Transform (FFT) to imaginary frequency. A similar exponential grid is used for the imaginary frequency. We split the self-energy in an exchange part $`\mathrm{\Sigma }^x`$ and a correlation part $`\mathrm{\Sigma }^c`$: $`\mathrm{\Sigma }_k^x(\text{r},\text{r}^{})`$ $`=`$ $`{\displaystyle \underset{q}{}}iG_q(\text{r},\text{r}^{},i\delta )\stackrel{~}{v}_{qk}(\text{r}^{},\text{r}),`$ (9) $`\mathrm{\Sigma }_k^c(\text{r},\text{r}^{},i\tau )`$ $`=`$ $`{\displaystyle \underset{q}{}}iG_q(\text{r},\text{r}^{},i\tau )W_{qk}^{\mathrm{scr}}(\text{r},\text{r},i\tau ),`$ (10) where $`\delta `$ is an infinitesimally small positive time, and $`W^{\mathrm{scr}}`$ is the screening interaction: $$W_k^{\mathrm{scr}}(\text{r},\text{r}^{},i\omega )W_k(\text{r},\text{r}^{},i\omega )\stackrel{~}{v}_k(\text{r},\text{r}^{}).$$ (11) In the calculation of $`\mathrm{\Sigma }^x`$ we use a 1D Brillouin-zone sampling of 10 equally spaced $`k`$-points, and in the calculation of $`\mathrm{\Sigma }^c`$ 4 $`k`$-points (since the screening interaction is short-ranged, the convergence of $`\mathrm{\Sigma }^c`$ with respect to the number of $`k`$-points is faster than that of $`\mathrm{\Sigma }^x`$). With the above parameters, the calculated QP gap has converged to within about 0.05 eV. From a two-pole fit on the imaginary-frequency axis an analytical continuation to the real-frequency axis is obtained: $`\mathrm{\Sigma }^c(i\omega )\mathrm{\Sigma }^c(\omega )`$. Subsequently, the QP equation Eq. (1) can be solved by replacing the QP wave functions by the DFT-LDA wave functions and obtaining $`E_{nk}`$ iteratively: $$E_{nk}=ϵ_{nk}+\varphi _{nk}\left\|\mathrm{\Sigma }_k^c(E_{nk})+\mathrm{\Sigma }_k^xV^{xc}\right\|\varphi _{nk}.$$ (12) ### B Treatment of the Coulomb interaction We have developed a novel procedure to deal with the $`\text{k}=0`$ and $`\text{r}=\text{r}^{}`$ singularities of the Coulomb interaction $`v_𝐤(\text{r},\text{r}^{})`$. We will first describe the procedure for a 3D system and later explain the specific adaptations of this procedure we used for our quasi-1D system. In reciprocal space the Coulomb interaction is given by: $$v_𝐤(\text{K},\text{K}^{})=\frac{4\pi }{\|𝐤+𝐊\|^2}\delta _{𝐊,𝐊^{}}.$$ (13) We replace $`v_0(0,0)`$, which would be infinite in Eq. (13) by a finite value, which is obtained in the following way. We evaluate the integral over all space of the Coulomb interaction multiplied by a Gaussian: $`I_\alpha `$ $`=`$ $`{\displaystyle d^3\text{q}\frac{4\pi }{q^2}e^{\alpha q^2}}=8\pi ^2\sqrt{{\displaystyle \frac{\pi }{\alpha }}},`$ (14) and evaluate the corresponding sum, excluding the singularity for $`𝐊=\text{k}=0`$: $$S_\alpha =\mathrm{\Delta }V\underset{𝐤,𝐊}{}^{}v_𝐤(\text{K},\text{K})e^{\alpha \|𝐤+𝐊\|^2},$$ (15) where $`\text{k}`$ 1BZ, the first Brillouin zone of the 3D lattice. $`\mathrm{\Delta }V`$ is the volume per $`(\text{k},\text{K})`$-point and the prime indicates that $`\text{k}=\text{K}=0`$ is excluded in this sum. We now put: $$v_{𝐤=0}(𝐊=0,𝐊^{}=0)\underset{\alpha 0}{lim}\left[I_\alpha S_\alpha \right].$$ (16) Finally, we obtain $`v(\text{r}\text{r}^{})`$ by a discrete FFT of $`v_𝐤(\text{K},\text{K}^{})`$ to real space. We find a finite value for $`v(\text{r}\text{r}^{}=0)`$, solving at the same time the problem with the Coulomb singularity for $`\text{r}\text{r}^{}=0`$. In the original formulation of the space-time method, the authors used a grid for $`\text{r}^{}`$ offset with respect the r-grid in order to avoid this singularity. In order to study a truly isolated chain, which is a quasi-1D system, we have to avoid ‘crosstalk’ between periodic images of the chain in the perpendicular directions. We do this by dividing space into regions of points that are closer to the atoms of a specific chain than to those of any other. Subsequently, we cut off the Coulomb interaction $`v(\text{r}\text{r}^{})`$, obtained in the way described above, by setting it zero if r and $`\text{r}^{}`$ belong to different regions. Thus we obtain an interaction $`\stackrel{~}{v}(\text{r},\text{r}^{})`$. In the construction of the Coulomb interaction $`v(\text{r}\text{r}^{})`$, we take a regular grid of k-points with a spacing in the $`y`$\- and $`z`$-direction approximately equal to that in the $`x`$-direction. From the cut-off interaction $`\stackrel{~}{v}(\text{r},\text{r}^{})`$ we obtain $`\stackrel{~}{v}_k(\text{r},\text{r}^{})`$ in mixed space from Eq. (3) with $`k`$ now in the 1D Brillouin zone. ### C The Bethe-Salpeter equation The two-body electron-hole Schrödinger equation related to the BSE is solved by expanding the exciton wave functions $`\mathrm{\Phi }(\text{r}_e,\text{r}_h)`$ in products of conduction $`\varphi _{ck}(\text{r}_e)`$ and valence wave functions $`\varphi _{vk}(\text{r}_h)`$: $$\mathrm{\Phi }(\text{r}_e,\text{r}_h)=\underset{k,c,v}{}A_{kcv}\varphi _{ck}(\text{r}_e)\varphi _{vk}^{}(\text{r}_h).$$ (17) Here we have restricted our discussion to excitons which have zero total momentum, since only these are optically active. As we are interested in the lowest lying excitons, an expansion in the highest occupied valence ($`\pi `$) and lowest unoccupied conduction ($`\pi ^{}`$) bands is sufficient to converge the exciton energies to within 0.1 eV; energy differences are converged even better. Below, we will give all energies in eV with a precision of two decimal places. The exciton binding energies $`E_b`$ follow from the Schrödinger-like equation: $$[E_{ck}E_{vk}E_g+E_b]A_{kcv}+\underset{k^{}c^{}v^{}}{}[2V_{kcv,k^{}c^{}v^{}}^x\delta _{s,0}W_{kcv,k^{}c^{}v^{}}]A_{k^{}c^{}v^{}}=0,$$ (18) where $`E_g`$ is the QP band gap, $`E_b`$ the exciton binding energy and $`W_{kcv,k^{}c^{}v^{}}`$ are the matrix elements of the static ($`\omega =0`$) screened interaction $$W_{kcv,k^{}c^{}v^{}}=d^3\text{r}d^3\text{r}^{}\varphi _{vk}(\text{r}^{})\varphi _{ck}^{}(\text{r})W(\text{r},\text{r}^{},\omega =0)\varphi _{v^{}k^{}}^{}(\text{r}^{})\varphi _{c^{}k^{}}(\text{r}),$$ (19) and $`V_{kcv,k^{}c^{}v^{}}^x`$ the exchange matrix elements (present for singlet excitons, $`s=0`$, only) of the bare Coulomb interaction: $$V_{kcv,k^{}c^{}v^{}}^x=d^3\text{r}d^3\text{r}^{}\varphi _{v^{}k^{}}^{}(\text{r})\varphi _{ck}^{}(\text{r}^{})\stackrel{~}{v}(\text{r},\text{r}^{})\varphi _{vk}(\text{r}^{})\varphi _{c^{}k^{}}(\text{r}).$$ (20) The integrals over space in Eqs. (19) and (20) are in the calculations replaced by summations over our real-space grid. We use wave functions and energies on a grid of 8 $`k`$-points and extrapolate to a grid of 100 $`k`$-points. Formally, dynamical screening effects may only be ignored in the BSE if $`E_gE_b`$. However, since it has been shown that dynamical effects in the electron-hole screening and in the one-particle Green function largely cancel each other, this approximation is nevertheless valid, even if the relation $`E_gE_b`$ does not strictly hold. We calculate an approximate exciton size $`a_{\mathrm{ex}}`$ by fitting the exciton coefficients $`A_{kcv}`$ to the hydrogen-like form: $$A_{kcv}=\frac{A_{k=0,cv}}{(1+a_{\mathrm{ex}}^2k^2)^2}.$$ (21) Note that in fact the exciton is highly anisotropic. Nevertheless, Eq. (21) gives pretty good fits and can be used to get a qualitative impression of the (relative) size of the excitons. ### D Inclusion of interchain screening As mentioned in the introduction, in a quasi-1D system, such as an isolated chain of a polymer in vacuum, there is no long-range screening. For a meaningful comparison of our calculations to the experimental data, which are obtained from either films or bulk polymer material, both the intra- and the interchain screening are important, and only the latter is long-ranged. It would be desirable to perform a $`GW`$A and exciton calculation for a 3D crystal structure of PT, but the amount of computational work involved is as yet prohibitively large. Since PT samples prepared in many different ways show very similar optical behavior, we expect the details of the interchain screening not to be extremely important. This consideration leads us to propose the following approximation for the interchain screening interaction, defined analogously to Eq. (11): $$W_{\mathrm{inter}}^{\mathrm{scr}}(\text{r},i\omega )=(1e^{r/r_{\mathrm{inter}}})\left\{\left[\epsilon _{}^2(i\omega )x^2+\epsilon _{\|\|}(i\omega )\epsilon _{}(i\omega )(y^2+z^2)\right]^{1/2}v(\text{r})\right\},$$ (22) where $`\epsilon _{}(i\omega )`$ and $`\epsilon _{\|\|}(i\omega )`$ are the ab-initio frequency-dependent dielectric constants perpendicular and parallel to the chain, respectively. The counter-intuitive combination of dielectric constants and coordinates in Eq. (22) results from solving the Laplace equation for a point charge in a homogeneous, anisotropic medium with dielectric constants $`\epsilon _{\|\|}`$ and $`\epsilon _{}`$. The prefactor takes care of a smooth cut-off for distances smaller than $`r_{\mathrm{inter}}`$, for which the interchain screening should be replaced by the intrachain screening. Eq. (22) has the correct behaviour for distances larger than the interchain distance $`r_{\mathrm{inter}}`$, for which we take 10 a.u., which is typical for the experimental crystal structures of Refs. and . The total screened interaction for the bulk system then becomes: $$W_{\mathrm{total}}(\text{r},\text{r}^{},i\omega )=W_{\mathrm{intra}}^{\mathrm{scr}}(\text{r},\text{r}^{},i\omega )+W_{\mathrm{inter}}^{\mathrm{scr}}(\text{r}\text{r}^{},i\omega )+v(\text{r}\text{r}^{}),$$ (23) where $`W_{\mathrm{intra}}^{\mathrm{scr}}`$ is the intrachain screening already calculated with Eq. (7). The screened interaction $`W_{\mathrm{total}}`$ is correct at short range, where the interchain screening is vanishingly small compared to the $`1/r`$-divergence of the intrachain screening, and at long range, where the intrachain screening vanishes due to its quasi-1D nature. Of course, for intermediate ranges, it is not strictly allowed to simply add the parts representing long- and short-ranged screening, but we expect Eq. (22) to give a reasonable interpolation there. Note that the interchain screening part given by Eq. (22) is long-ranged by construction, and 8 $`k`$-points are now needed to converge the corresponding self-energy $`\mathrm{\Sigma }_{\mathrm{inter}}^c`$ from Eq. (10). On the other hand, the required number of real-space grid points in order to calculate $`\mathrm{\Sigma }_{\mathrm{inter}}^c`$ is less than before, because $`W_{\mathrm{inter}}^{\mathrm{scr}}`$ is a very smooth function of r; a $`12\times 12\times 12`$ real space grid turns out to be sufficient. The total self-energy can be expressed as: $$\mathrm{\Sigma }_{\mathrm{total}}=\mathrm{\Sigma }_{\mathrm{intra}}^c+\mathrm{\Sigma }_{\mathrm{inter}}^c+\mathrm{\Sigma }^x.$$ (24) Because the self-energies in this equation are additive, we can reuse the self-energies $`\mathrm{\Sigma }_{\mathrm{intra}}^c`$ and $`\mathrm{\Sigma }^x`$, which we have already calculated for the isolated chain. The overlap between wave functions, and therefore the electronic coupling between neighboring chains, is very small. This means that we can use the isolated-chain wave functions to calculate the Green function and self-energy. This obviously implies that in our exciton calculations we restrict ourselves to excitons in which we take the electron and hole are on the same chain (so-called intrachain excitons). In summary, the only, but essential, difference between our calculations for the isolated PT chain and bulk PT is in the use of an interchain screened interaction. ### E Dielectric tensor of crystalline PT In order to construct the screened interaction of Eq. (22), we have to calculate the dielectric tensor of bulk PT. We do this for the crystalline structure of Ref. , which is reproduced in Fig. 1. We use a model in which the chains are replaced by polarizable line objects with a polarizability tensor obtained from the single-chain polarizability function. The principal axes of the chain are the following: $$\widehat{x}_1=\widehat{x},\text{ }\widehat{x}_2=\frac{1}{\sqrt{2}}(\widehat{z}+\widehat{y}),\text{ }\widehat{x}_3=\frac{1}{\sqrt{2}}(\widehat{z}\widehat{y}).$$ (25) The full polarizability function $`X(\text{r},\text{r}^{},i\omega )`$ of a single chain is given by: $`X(\text{r},\text{r}^{},i\omega )`$ $`=`$ $`P(\text{r},\text{r}^{},i\omega )+{\displaystyle 𝑑\text{r}^{\prime \prime }𝑑\text{r}^{\prime \prime \prime }P(\text{r},\text{r}^{\prime \prime },i\omega )W_{\mathrm{intra}}(\text{r}^{\prime \prime },\text{r}^{\prime \prime \prime },i\omega )P(\text{r}^{\prime \prime \prime },\text{r}^{},i\omega )}`$ (26) $``$ $`X^{(0)}(\text{r},\text{r}^{},i\omega )+X^{(1)}(\text{r},\text{r}^{},i\omega ).`$ (27) The long-wavelength ($`q0`$) polarizability tensor $`\underset{¯}{\chi }`$ per unit of chain length of a single chain in the $`(x_1,x_2,x_3)`$ coordinate system is diagonal and has diagonal elements given by: $$\chi _1(i\omega )=\underset{q0}{lim}\left[\frac{1}{q^2}𝑑\text{r}𝑑\text{r}^{}e^{iq(x_1x_1^{})}X(\text{r},\text{r}^{},i\omega )\right],$$ (28) and for $`j`$=2,3: $$\chi _j(i\omega )=𝑑\text{r}𝑑\text{r}^{}x_jX(\text{r},\text{r}^{},i\omega )x_j^{}.$$ (29) The calculation of $`\chi _j(i\omega )`$ has been performed with 4 $`k`$-points, with the exception of $`\chi _1^{(0)}`$, for which it proved to be necessary to use 8 $`k`$-points. If we now approximate the chains by polarizable line objects with the above polarizability tensor, we can calculate the macroscopic dielectric tensor of the crystal of these chains. This is done by a procedure of which the details are given in the Appendix. The axes of the crystal unit cell are denoted by $`\widehat{a}`$, $`\widehat{b}`$ and $`\widehat{c}=\widehat{x}=\widehat{x}_1`$. Dropping the frequency dependence in the notation, we find the following expression for $`\epsilon _c`$: $$\epsilon _c=1+\frac{4\pi \chi _1}{A},$$ (30) where $`A`$ is the surface area per chain in the plane perpendicular to the chain. For $`\epsilon _a`$ and $`\epsilon _b`$ we find: $$\epsilon _\gamma =\frac{1}{1\frac{4\pi }{A}\stackrel{~}{\chi }_\gamma },$$ (31) where $`\gamma =a,b`$ and $`\stackrel{~}{\chi }_\gamma `$ is the effective polarizability of the chain along the $`\gamma `$-axis. In the Appendix details of the calculation of $`\epsilon _a`$, $`\epsilon _b`$, and $`\epsilon _c`$ are given. To retain the tetragonal symmetry in our calculation (in order keep the computations feasible), we average $`\epsilon _a(i\omega )`$ and $`\epsilon _b(i\omega )`$, which are not very different, to obtain $`\epsilon _{}(i\omega )`$. For $`\epsilon _{\|\|}(i\omega )`$ we take $`\epsilon _c(i\omega )`$. Note that for using the screened interaction of Eq. (22) to in the implementation of the $`GW`$A formalism presented in Section II A, we have calculated the dielectric constants $`\epsilon _{}(i\omega )`$ and $`\epsilon _{\|\|}(i\omega )`$ along the imaginary frequency axis. ## III Results ### A Isolated chain The calculated $`GW`$A QP band structure (together with the DFT-LDA band structure) is shown in Fig. 2, left panel. We find a minimal band gap $`E_g`$ at $`\mathrm{\Gamma }`$ of 3.59 eV, which is quite large compared to the DFT-LDA value of 1.22 eV. The effective masses, $`m^{}=1/\mathrm{}^2(^2E/k^2)^1`$, of the $`\pi `$ and $`\pi ^{}`$ bands at $`\mathrm{\Gamma }`$, which are 0.15 and 0.17 $`m_e`$ (with $`m_e`$ the free electron mass) in DFT-LDA, are reduced by about 15% in the $`GW`$A to 0.13 and 0.15 $`m_e`$. This corresponds to an increase of the band width from 1.96 and 1.51 eV in DFT-LDA to 2.48 and 1.81 eV in the $`GW`$A, for the $`\pi `$ and $`\pi ^{}`$ bands, respectively. In an earlier $`GW`$A study, a similar increase of the bandwidth was found for a wide variety of materials. The lowest-lying singlet exciton (<sup>1</sup>B<sub>u</sub>) has a binding energy $`E_b`$ of 1.85 eV. The size $`a_{\mathrm{ex}}`$ of this exciton, calculated using Eq. (21), is 12 a.u., i.e. less than two thiophene rings. To give an impression of the exciton wave function $`\mathrm{\Phi }(\text{r}_e,\text{r}_h)`$, we have plotted in Fig. 3 (top panel) the probability to find the hole at a distance $`x_h`$ along the chain from the electron, $$\mathrm{Prob}(x_h)𝑑y_h𝑑z_h\|\mathrm{\Phi }(\text{r}_e,\text{r}_h)\|^2,$$ (32) where the electron coordinate $`\text{r}_e`$ is taken 1 a.u. from the inversion center, in the direction perpendicular to the polymer plane (for the electron coordinate in the inversion center, this probability would be zero due to symmetry). We have plotted $`\mathrm{Prob}(x_h)`$ for both the <sup>1</sup>B<sub>u</sub> and <sup>1</sup>A<sub>g</sub> excitons. As the optical gap is given by $`E_o=E_gE_b`$, we have $`E_o=1.74`$ eV, in good agreement with the experimental value of 1.8 eV (see Table I). While there is good agreement for the optical gap, the difference between the $`{}_{}{}^{1}\mathrm{B}_{\mathrm{u}}^{}`$ and $`{}_{}{}^{1}\mathrm{A}_{\mathrm{g}}^{}`$ binding energies of the isolated PT chain is definitely not in agreement with experiment, see Table I. Moreover, the $`{}_{}{}^{1}\mathrm{B}_{\mathrm{u}}^{}`$ exciton binding energy of 1.85 eV is very large compared to values currently discussed in the literature, which range from $`0.1`$ to $`1.0`$ eV. ### B Dielectric properties We calculate the polarizabilities per unit length $`\chi _j(i\omega )`$ with Eqs. (28) and (29). The obtained $`\omega =0`$ values are listed in Table II. Note that the polarizability along the chain, i.e. in the direction of the extended carbon $`\pi `$-system, is much larger than those in the perpendicular directions. This difference is reflected in the dielectric constants $`\epsilon _\gamma (i\omega )`$ calculated using Eqs. (30) and (31); the dielectric constant along the chain is much larger than those in the perpendicular directions. In real systems with disorder the conjugation length will be finite, which will reduce $`\epsilon _{\|\|}`$. Note, however, that the perpendicular dielectric constant $`\epsilon _{}`$ plays the dominant role in the interchain screening of Eq. (22) along the chain. ### C Crystalline polythiophene The resulting band structure, calculated using the bulk screening from Eqs. (22) and (23) is given in Fig. 2, right panel. The QP gap $`E_g`$ has decreased to 2.49 eV; the <sup>1</sup>B<sub>u</sub> exciton binding energy is 0.76 eV (see Table I). Hence, the predicted optical gap is 1.73 eV, virtually unchanged from the isolated chain results of 1.74 eV and in good agreement with experiment. Note that the absorption gap of Ref. is 2.0 eV, also found in earlier work on PT, but the luminescence gap is 1.8 eV. There are two reasons why we should compare our result to the latter gap. The first reason is that absorption occurs everywhere in a sample, both in the ordered and disordered parts, but luminescence occurs predominantly in the most ordered parts with the longest conjugation lengths. This is because, prior to recombination, excitons diffuse to those parts of the sample where they have the lowest energy. The second reason is that after photoexcitation, the rings, which may be twisted around their common C-C bond, tend to co-planarize in the excited state, due to the fact that the excited state is slightly more quinoid than the aromatic ground state. As we are performing our calculations for a perfect, co-planar chain of PT, we should therefore compare our optical gap to the luminescence gap. Note that in principle it is possible that excitons trapped in defects or disordered parts of the sample to have a lower energy than in a fully conjugated, defect-free polymer. However, the luminescence spectrum of Ref. can be fully understood in terms of the <sup>1</sup>B<sub>u</sub> exciton decay and its vibronic side bands, which means that such defects are either rare or that excitons trapped by such defects decay non-radiatively. What is very important, is that the relative exciton energies (also listed in Table I) are now also in good agreement with experiment. The sizes of the excitons have increased by $`50`$%; the <sup>1</sup>B<sub>u</sub> size $`a_{\mathrm{ex}}`$ is now 18 a.u., or slightly more than two rings. In Fig. 3 (bottom panel) it is clearly seen that the excitons are larger than the corresponding excitons on the isolated chain (top panel). In order to test the sensitivity of our results to the precise value of the cutoff distance $`r_{\mathrm{inter}}`$ in Eq. (22) we performed similar calculations for $`r_{\mathrm{inter}}`$ = 8 a.u. and $`r_{\mathrm{inter}}`$ = 12 a.u. These data are also listed in Table I. The QP band gaps are 2.32 and 2.61 eV, respectively. The $`{}_{}{}^{1}\mathrm{B}_{\mathrm{u}}^{}`$ binding energies are 0.64 and 0.86 eV and hence the optical gaps are 1.68 and 1.73 eV, respectively. This means that the optical gap is quite insensitive to the choice of $`r_{\mathrm{inter}}`$. This is consistent with the fact that in the limit $`r_{\mathrm{inter}}\mathrm{}`$, which corresponds to no interchain screening, we should find the isolated chain absorption gap of 1.74 eV. The energy differences between the excitons are even less sensitive to $`r_{\mathrm{inter}}`$. The good agreement with experiment and the fact that especially the optical gap and the energy separation between the excitons are hardly influenced by varying $`r_{\mathrm{inter}}`$ are also a posteriori justifications for our model screening interaction Eq. (23). ## IV Conclusions and discussion Summing up, we have calculated the quasi-particle band structure and lowest-lying exciton binding energies of an isolated polythiophene chain and crystalline polythiophene. For the isolated chain (where there is only intrachain screening) we find a large band gap and large exciton binding energies, due to the absence of long-range screening. After including interchain screening, which is responsible for the long-range screening in bulk polythiophene, we find that both the band gap and exciton binding energies are drastically reduced. However, the optical gap is hardly affected. We suggest that these conclusions hold for conjugated polymers in general. This sheds light on the fact that the calculations by Rohlfing and Louie on isolated chains of PA and PPV yield good results for the optical gaps, whereas their lowest-lying singlet exciton binding energy of 0.9 eV for PPV is in excess of recent experimental values of $`0.35\pm 0.15`$ eV, obtained by a direct STM measurement for an alkoxy-substituted PPV, and $`0.48\pm 0.14`$ eV for unsubstituted PPV. The inclusion of interchain screening effects will drastically reduce their calculated binding energy and may well lead to agreement with this experiment. Clearly, it would also be very interesting to repeat the experiment in Ref. for polythiophene and polyacetylene. Interestingly enough, a value of 0.4 eV is obtained for the exciton binding energy in PPV by means of an effective-mass appromixation in which the electron-hole interaction is derived from a bulk dielectric tensor. The difference of about a factor of two in exciton binding energy between crystalline PT and PPV can, at least qualitatively, be explained by the differences in reduced masses $`\mu `$ ($`1/\mu =1/m_\pi +1/m_\pi ^{}`$) of PT and PPV, for which we find $`\mu ^{\mathrm{PT}}=0.08m_e`$, while $`\mu ^{\mathrm{PPV}}=0.04m_e`$ both in DFT-LDA, and by the fact that in an effective-mass approximation the binding energy is proportional to $`\mu `$. Of course, these arguments, which are qualitative only, do not take away the need for ab-initio calculations on the crystalline phase of PPV. Further, the apparent discrepancy of the results for PA by Ethridge et al. and those of Rohlfing and Louie, can be understood. The latter find, for an isolated chain, a QP gap of 2.1 eV and an exciton binding energy of 0.4 eV, yielding an absorption gap of 1.7 eV. The former find a QP gap of 1.86 eV and do not include excitonic effects. This calculation, however, is performed for one PA chain in the same volume as a PA chain in a crystal would have. Therefore, this calculation is in fact one for a bulk situation, which means that this QP gap is by our arguments expected to be smaller than that of Rohlfing and Louie. Furthermore, our arguments predict an exciton binding energy in bulk PA considerably smaller than the 0.4 eV of Rohlfing and Louie. We conclude that a correct many-body description of the electronic and optical properties of bulk polymer systems should include the effect of interchain screening. An important consequence of this conclusion is that neither Hartree-Fock nor DFT-LDA calculations should be relied upon in this context, since Hartree-Fock does not contain screening effects at all and since the exchange-correlation potential in DFT-LDA only depends on the local density and cannot describe the non-local effects due to the long-range screening. Moreover, since exciton effects play such a large role in conjugated polymers, it is essential to take these effects into account. ## Acknowledgements Financial support from NCF (Nationale Computer Faciliteiten) project SC-496 is acknowledged. G.B. acknowledges the financial support from Philips Research through the FOM-LZM (Fundamenteel Onderzoek der Materie - Laboratorium Zonder Muren) program. ## A Calculation of the crystal dielectric tensor within a line-dipole model We apply an electric field $`\text{E}_{\mathrm{appl}}(\text{r})=\text{E}_0e^{i𝐤𝐫}`$ (and we will take the limit $`k0`$), where $`\text{E}_0`$ and k are parallel to the $`a`$,$`b`$ or $`c`$-axis of the crystal (see Fig. 1) to calculate $`\epsilon _a`$, $`\epsilon _b`$ and $`\epsilon _c`$, respectively. The applied field $`\text{E}_{\mathrm{appl}}`$ leads to an induced field $`\text{E}_{\mathrm{ind}}(\text{r})`$; the total microscopic field $`\text{E}_{\mathrm{micr}}(\text{r})`$ is then given by: $$\text{E}_{\mathrm{micr}}(\text{r})=\text{E}_{\mathrm{appl}}(\text{r})+\text{E}_{\mathrm{ind}}(\text{r}).$$ (A1) We define $`\stackrel{}{\rho }=u\widehat{a}+v\widehat{b}`$ with $`\rho ^2=u^2+v^2`$. Note that there are two different chains: the $`𝒜`$ type, at the corners of the unit cell, and the $``$ type at the center of the unit cell. For the $`𝒜`$ and $``$ chain we have: $`\text{p}_𝒜(x)`$ $`=`$ $`\underset{¯}{\chi }_𝒜\text{E}_{\mathrm{micr}}^{}(x,\stackrel{}{\rho }=0)`$ (A2) $`\text{p}_{}(x)`$ $`=`$ $`\underset{¯}{\chi }_{}\text{E}_{\mathrm{micr}}^{}(x,\stackrel{}{\rho }={\displaystyle \frac{1}{2}}\widehat{a}+{\displaystyle \frac{1}{2}}\widehat{b})`$ (A3) with $`\text{p}_𝒜(x)`$ ($`\text{p}_{}(x)`$) the long-wavelength dipole moment per unit length of the $`𝒜`$ ($``$) chain and $`\underset{¯}{\chi }_𝒜`$ ($`\underset{¯}{\chi }_{}`$) the polarizability tensor of the $`𝒜`$ ($``$) chain calculated with Eqs. (28) and (29) and using the relations: $$\underset{¯}{\chi }_𝒜=\underset{¯}{U}_𝒜^1\underset{¯}{\chi }\underset{¯}{U}_𝒜,\text{ }\underset{¯}{\chi }_{}=\underset{¯}{U}_{}^1\underset{¯}{\chi }\underset{¯}{U}_{},$$ (A4) with $`\underset{¯}{U}_𝒜`$ and $`\underset{¯}{U}_{}`$ are the rotation matrices relating the ($`x_2`$,$`x_3`$) coordinate system to the ($`a`$,$`b`$) coordinate system ($`\widehat{c}=\widehat{x}_1`$): $`\underset{¯}{U}_𝒜`$ $`=`$ $`\left(\begin{array}{cc}\hfill \mathrm{cos}(\frac{\pi }{4}\alpha )& \hfill \mathrm{sin}(\frac{\pi }{4}\alpha )\\ \hfill \mathrm{sin}(\frac{\pi }{4}\alpha )& \hfill \mathrm{cos}(\frac{\pi }{4}\alpha )\end{array}\right),`$ (A7) $`\underset{¯}{U}_{}`$ $`=`$ $`\left(\begin{array}{cc}\hfill \mathrm{cos}(\frac{3\pi }{4}\alpha )& \hfill \mathrm{sin}(\frac{3\pi }{4}\alpha )\\ \hfill \mathrm{sin}(\frac{3\pi }{4}\alpha )& \hfill \mathrm{cos}(\frac{3\pi }{4}\alpha )\end{array}\right).`$ (A10) The prime in Eqs. (A2) and (A3) indicates that the field caused by the chain itself is excluded. We will refer to our model, in which a PT chain is represented by an homogeneous line with a certain dipole moment per unit length p, as a ‘line-dipole’. In CGS units the dielectric tensor $`\underset{¯}{\epsilon }`$ is defined as: $$\text{E}(\text{r})+4\pi \text{P}(\text{r})=\underset{¯}{\epsilon }\text{E}(\text{r}).$$ (A11) where $`\text{E}(\text{r})`$ is the macroscopic field, and $`\text{P}(\text{r})`$ is the macroscopic polarization. For each direction of the applied field, we will calculate $`\text{E}_{\mathrm{ind}}(\text{r})`$, evaluate the macroscopic fields $`\text{E}(\text{r})`$ and $`\text{P}(\text{r})`$ by averaging, and solve Eq. (A11) to obtain the dielectric tensor $`\underset{¯}{\epsilon }`$. ### 1 Calculation of $`\epsilon _c`$ For $`\text{E}_{\mathrm{appl}}`$ and k parallel to $`\widehat{x}`$ (and hence to $`\widehat{c}`$ and also $`\widehat{x}_1`$), we have for both the $`𝒜`$ and $``$ chain from Eqs. (A2) and (A3): $$p_x(x)=\chi _1E_x^{}(x)$$ (A12) The field induced by a line-dipole on the $`x`$-axis is given by: $$\text{E}_{\mathrm{ind}}(\text{r})=\stackrel{}{}\mathrm{\Phi }(\text{r})=\stackrel{}{}\frac{p_xe^{ikx^{}}(xx^{})}{\|\text{r}\text{r}^{}\|^3}𝑑x^{},$$ (A13) where $`\mathrm{\Phi }`$ is the electrostatic potential and we have used the fact that $`p_x(x^{})=p_xe^{ikx^{}}`$. Evaluation of Eq. (A13) yields: $$E_{\mathrm{ind},x}(\text{r})=2k^2p_xK_0(\rho k)e^{ikx},$$ (A14) where $`K_0`$ is a zeroth order Bessel function of the third kind. From here on, we omit the factor $`e^{ikx}`$. We can calculate the total microscopic field at the $`x`$-axis, due to both applied and induced fields, for a crystal of line-dipoles, by summing over all line-dipoles but the one at the origin: $$E_{\mathrm{micr},x}^{}(\stackrel{}{\rho }=0)=E_{\mathrm{appl},x}(\stackrel{}{\rho }=0)+\underset{\stackrel{}{\rho _i}0}{}E_{\mathrm{ind},x}(\stackrel{}{\rho }_i)$$ (A15) where the positions of the other chains are given by $`\stackrel{}{\rho _i}`$. In the limit $`k0`$, we can replace the sum by an integral: $`\underset{k0}{lim}{\displaystyle \underset{\stackrel{}{\rho _i}0}{}}E_{\mathrm{ind},x}(\stackrel{}{\rho _i})`$ $`=`$ $`{\displaystyle \frac{2\pi p_x}{A}}{\displaystyle _0^{\mathrm{}}}\rho ^{}𝑑\rho ^{}K_0(\rho ^{})`$ (A16) $`=`$ $`{\displaystyle \frac{4\pi }{A}}p_x,`$ (A17) where $`\rho ^{}`$ = $`\rho k`$ and $`A=ab/2`$ is the area of the two dimensional unit cell per chain. Substitution of Eq. (A17) and (A15) in Eq. (A12) yields: $$p_x=\frac{\chi _1A}{A+4\pi \chi _1}E_{\mathrm{appl},\mathrm{x}}.$$ (A18) Since $`P_x=p_x/A`$, we have: $$P_x=\frac{\chi _1}{A+4\pi \chi _1}E_{\mathrm{appl},x}.$$ (A19) The macroscopic field $`E_x`$ is the average over the two-dimensional unit cell of the microscopic field as given by Eq. (A15) for general $`\stackrel{}{\rho }`$, but now including the chain at $`\stackrel{}{\rho _i}=0`$ in the sum: $`E_x`$ $`=`$ $`E_{\mathrm{appl},x}+\underset{k0}{lim}{\displaystyle \frac{1}{2A}}{\displaystyle _{\mathrm{unit}\text{ }\mathrm{cell}}}d^2\stackrel{}{\rho }{\displaystyle \underset{\stackrel{}{\rho _i}}{}}E_{\mathrm{ind},x}(\stackrel{}{\rho }\stackrel{}{\rho _i})`$ (A20) $`=`$ $`E_{\mathrm{appl},x}+\underset{k0}{lim}{\displaystyle \frac{2\pi }{A}}{\displaystyle \rho 𝑑\rho \underset{\stackrel{}{\rho }_i}{}E_{\mathrm{ind},x}(\stackrel{}{\rho }\stackrel{}{\rho _i})}`$ (A21) $`=`$ $`E_{\mathrm{appl},x}{\displaystyle \frac{4\pi p_x}{A}}`$ (A22) Combining this with Eqs. (A11) and (A19) we obtain Eq. (30): $$\epsilon _c=1+\frac{4\pi \chi _1}{A}.$$ (A23) ### 2 Calculation of $`\epsilon _a`$ and $`\epsilon _b`$ We now take $`\text{E}_{\mathrm{appl}}(\text{r})`$ and k parallel to $`\widehat{a}`$. The derivation for $`\text{E}_{\mathrm{appl}}(\text{r})`$ and k parallel to $`\widehat{b}`$ is equivalent. The dipole moments of the chains must satisfy Eqs. (A2) and (A3). The field induced by the chain at the origin is given by: $`\text{E}_{\mathrm{ind}}(\text{r})=\mathrm{\Phi }(\text{r})`$ $`=`$ $`{\displaystyle \frac{\text{p}_𝒜\text{r}}{\|\text{r}\text{r}^{}\|^3}𝑑x^{}}`$ (A24) $`=`$ $`\underset{¯}{M}(\rho )\text{p}_𝒜,`$ (A25) where $$\underset{¯}{M}(\rho )\left(\begin{array}{cc}\frac{4u^2}{\rho ^4}\frac{2}{\rho ^2}& \frac{4uv}{\rho ^4}\\ \frac{4uv}{\rho ^4}& \frac{4v^2}{\rho ^4}\frac{2}{\rho ^2}\end{array}\right)$$ (A26) in the two dimensional ($`a`$,$`b`$) coordinate system (the dipole moment in the $`c`$-direction is zero and hence we work with $`2\times 2`$ instead of $`3\times 3`$ matrices). The microscopic electric field $`\text{E}_{\mathrm{micr}}^{}`$ at the origin, excluding the field induced by chain at the origin itself, is given by: $$\text{E}_{\mathrm{micr}}^{}(\stackrel{}{\rho }=0)=\text{E}_{\mathrm{appl}}(\stackrel{}{\rho }=0)+\underset{¯}{M}_𝒜\text{p}_𝒜+\underset{¯}{M}_{}\text{p}_{},$$ (A27) where $`\underset{¯}{M}_𝒜`$ $``$ $`\underset{k0}{lim}{\displaystyle \underset{\stackrel{}{\rho _i}𝒜,\stackrel{}{\rho _i}0}{}}\underset{¯}{M}(\stackrel{}{\rho _i})\mathrm{cos}ku_i,`$ (A28) $`\underset{¯}{M}_{}`$ $``$ $`\underset{k0}{lim}{\displaystyle \underset{\stackrel{}{\rho _i}}{}}\underset{¯}{M}(\stackrel{}{\rho _j})\mathrm{cos}ku_j,`$ (A29) These sums are evaluated in the next subsection. Substitution of Eq. (A27) in Eq. (A2) and solving yields: $$\stackrel{~}{\chi }_ap_{𝒜,a}/E_{\mathrm{appl},a}=\left(\underset{¯}{\chi }_𝒜\left[\underset{¯}{1}\underset{¯}{\chi }_𝒜\underset{¯}{M}_𝒜\underset{¯}{\chi }_{}\underset{¯}{M}_{}\right]^1\right)_{aa},$$ (A30) with $`\underset{¯}{\chi }_𝒜`$ and $`\underset{¯}{\chi }_{}`$ as defined in Eq. (A4). Analogous to the derivation given by Jackson for a point dipole, we can derive the electric field of a line-dipole at $`\stackrel{}{\rho }=0`$: $$\text{E}(\stackrel{}{\rho })=\left(\underset{¯}{M}(\stackrel{}{\rho })2\pi \delta (\stackrel{}{\rho })\underset{¯}{1}\right)\text{p}$$ (A31) where the convention in Eq. (A31) is that the field within the line-dipole at $`\stackrel{}{\rho }=0`$ is given by the term $`2\pi \delta (\stackrel{}{\rho })𝐩`$ and the Cauchy principal value of the integral should be taken in integrals across the $`1/\rho ^2`$ singularity at $`\stackrel{}{\rho }=0`$. The macroscopic field is given by the average over the microscopic field of Eq. (A27) for general $`\stackrel{}{\rho }`$ including the chain at the origin. Note that since, by symmetry, $`p_{𝒜,b}=p_{,b}`$, the $`b`$ components do not contribute to the macroscopic field. Also by symmetry, we have $`p_{𝒜,a}=p_{,a}=p_a`$. We then have for the macroscopic field $`E_a(\stackrel{}{\rho })`$ : $`E_a(\stackrel{}{\rho })`$ $`=`$ $`E_{\mathrm{appl},a}+\underset{k0}{lim}{\displaystyle \frac{1}{2A}}𝒫{\displaystyle _{\mathrm{unit}\text{ }\mathrm{cell}}}d^2\stackrel{}{\rho }{\displaystyle \underset{\rho _k}{}}M_{aa}(\stackrel{}{\rho }\stackrel{}{\rho _k})\mathrm{cos}(ku_k)p_a{\displaystyle \frac{2\pi }{A}}p_a`$ (A32) $`=`$ $`E_{\mathrm{appl},a}+\underset{k0}{lim}{\displaystyle \frac{1}{A}}𝒫{\displaystyle d^2\stackrel{}{\rho }M_{aa}(\stackrel{}{\rho })\mathrm{cos}(ku)p_a}{\displaystyle \frac{2\pi }{A}}p_a`$ (A33) $`=`$ $`E_{\mathrm{appl},a}{\displaystyle \frac{4\pi p_a}{A}}.`$ (A34) Substituting this result in Eq. (A11) and using the fact that $`P_a=p_a/A`$, we find Eq. (31): $$\epsilon _a=\frac{1}{1\frac{4\pi }{A}\stackrel{~}{\chi }_a}.$$ (A35) A similar result is obtained for $`\epsilon _b`$. ### 3 Evaluation of $`\underset{¯}{M}_𝒜`$ and $`\underset{¯}{M}_{}`$ From the symmetry of Eqs. (A26), (A28) and (A29), we see that $`\underset{¯}{M}_{𝒜,ab}=\underset{¯}{M}_{𝒜,ba}=0=\underset{¯}{M}_{,ab}=\underset{¯}{M}_{,ba}=0`$ and $`\underset{¯}{M}_{𝒜,aa}=\underset{¯}{M}_{𝒜,bb}`$ and $`\underset{¯}{M}_{,aa}=\underset{¯}{M}_{,bb}`$. This leaves us with only one element of each matrix to be determined. Considering $`\underset{¯}{M}_𝒜`$ first, we split the summation of Eq. (A28) into two parts. For $`\rho _i<R`$ (with $`R`$ large) we perform the summation explicitly (taking $`k=0`$), while for $`\rho _iR`$ we replace the summation by an integral: $$\underset{¯}{M}_{𝒜,aa}=\underset{\stackrel{}{\rho _i}𝒜\text{ }\stackrel{}{\rho _i}0\text{ }\rho _i<R}{}\underset{¯}{M}_{aa}(\stackrel{}{\rho _i})+\underset{k0}{lim}\frac{1}{2A}_R^{\mathrm{}}\rho 𝑑\rho _0^{2\pi }𝑑\varphi \underset{¯}{M}_{aa}(\rho )\mathrm{cos}(k\rho \mathrm{cos}\varphi ),$$ (A36) which is exact in the limit $`R\mathrm{}`$. The sum is evaluated numerically; its values is $`0.009677`$ $`a_0^2`$ in the limit $`R\mathrm{}`$. The integral becomes $`\pi /(2A)`$ after first taking the limit $`k0`$ and then the limit $`R\mathrm{}`$. We can calculate $`\underset{¯}{M}_{,aa}`$ in a similar way. The sum yields $`0.012035`$ $`a_0^2`$ and the integral becomes again $`\pi /(2A)`$. Therefore, $`\underset{¯}{M}_𝒜`$ and $`\underset{¯}{M}_{}`$ are: $`\underset{¯}{M}_𝒜`$ $`=`$ $`\left(\begin{array}{cc}0.030068\text{ }a_0^2\hfill & 0\hfill \\ 0\hfill & 0.030068\text{ }a_0^2\hfill \end{array}\right),`$ (A39) $`\underset{¯}{M}_{}`$ $`=`$ $`\left(\begin{array}{cc}0.008357\text{ }a_0^2\hfill & 0\hfill \\ 0\hfill & 0.008357\text{ }a_0^2\hfill \end{array}\right).`$ (A42)
warning/0004/cond-mat0004300.html	ar5iv	text	# Contributions of spontaneous phase slippage to linear and non–linear conduction near the Peierls transition in thin samples of o -TaS3 ## I Introduction Though the Peierls transition in quasi 1-dimensional conductors was discovered more than 30 years ago, its mechanism and the role of fluctuations still remain an unsettled problem . The fluctuations in quasi 1-dimensional compounds below $`T_P`$ are seen from various studies including transport properties, such as the large width of the Peierls transition in comparison with that expected from the BCS-type onset of the gap, the smeared edge of the Peierls gap revealed through the optical investigations , the appearance of spontaneous current noise, which is associated with thermally initiated phase slip (PS) developing several kelvins below $`T_P`$ . A decrease of the cross-section area $`s`$ of the samples results in growth of the fluctuations. For example, in the samples of o-TaS<sub>3</sub> with $`s10^2`$ $`\mu `$m<sup>2</sup> and below, the Peierls transition is smeared out and substantially shifted down to lower temperatures . Conductance hysteresis in such thin samples is absent within decades of kelvins below $`T_P`$; in this temperature range spontaneous PS is observed, and the conductivity strongly deviates from the Arrhenius law . Another fluctuation effect known as threshold rounding consists in smearing out of the onset of the non-linear current at the threshold field $`E_T`$ . This effect is found in NbSe<sub>3</sub>. The rounding increases both with increasing $`T`$ and decreasing thickness $`t`$ of the crystals; in the thinnest crystals the growth of conductivity starts from zero field. In phase slippage has been discussed as a possible basis of the rounding, but the authors did not find enough arguments in favor of this explanation. Another interpretation was found to be more reasonable : the rounding was attributed to the thermally–assisted creep of charge–density waves (CDWs) in the framework of the weak-pinning model . This approach implies that in very thin samples the pinning energy of the phase-correlation volume becomes comparable with $`kT`$, and activated creep of the continuous CDW within the correlation lengths $`L_{2\pi }`$ is possible. Estimates for NbSe<sub>3</sub> based on the mean-field BCS dependence for $`\mathrm{\Delta }(T)`$, showed that thermal depinning of the CDW is probable. This interpretation, however, is rather dubious for TaS<sub>3</sub>, where the mean-field approach fails near $`T_P`$: in highly anisotropic compounds such as TaS<sub>3</sub> and K<sub>0.3</sub>MoO<sub>3</sub>, the onset of the gap near $`T_P`$ does not follow the BCS dependence , and the pseudogap does not vanish tens of kelvins above $`T_P`$ . In the present paper the threshold rounding in thin samples of TaS<sub>3</sub> is reported. Independently we observe a fluctuation contribution to the linear conductivity. It is shown that creep of the continuous CDW cannot account for the threshold rounding in TaS<sub>3</sub>. Alternatively, we show that spontaneous PS observed near $`T_P`$ results in local creep of the CDW and contributes to the linear and non–linear conductivity, in agreement with our experiment. The result is generalized for large samples. We discuss the mechanism of the Peierls transition in the light of the PS – induced creep. ## II Experimental Technique and Results Thin samples of TaS<sub>3</sub> were placed on sapphire substrates. We used vacuum–deposited indium contacts . The cross–section area of the samples was estimated from the values of the room- temperature resistance ($`3\times 10^4`$ $`\mathrm{\Omega }`$cm) and the visible contact separation . Similar results are observed on 5 samples from high-quality batches. Most of the data reported here are obtained on the representative sample with the dimensions $`L=4.5`$ $`\mu `$m, $`s=0.3\times 10^3`$ $`\mu `$m<sup>2</sup>. The dependencies of conduction $`\sigma `$ on temperature and voltage $`V`$ are presented in Figs. 1 and 2 respectively. One can see (Fig. 1) that the Peierls transition is smeared out in comparison with the transitions in usual-sized samples (shown with a dotted line), in agreement with Ref. . Deviation from the Arrhenius law is observed tens of kelvins below $`T_P`$ (indicated by an arrow), the latter being considerably shifted downwards in comparison with $`T_P=220`$ K observed in thick samples. The activation dependence $`\sigma \mathrm{exp}(\mathrm{\Delta }/T)`$ with $`\mathrm{\Delta }=800`$ K extrapolated from the low temperatures is shown by the broken line; we denote the corresponding conductivity as $`\sigma _\mathrm{\Delta }`$. We shall consider the temperature and sample-size dependence of the Peierls gap, $`2\mathrm{\Delta }`$, to be insignificant, which is supported by the results of Refs. . Then, the difference $`\delta \sigma \sigma \sigma _\mathrm{\Delta }`$ can be considered as the fluctuation contribution to the conductivity. The Arrhenius plot $`\delta \sigma `$ vs. $`1/T`$ is shown in Fig. 3. The dependence is close to a straight line up to $`TT_P`$; the activation energy $`W`$ being about 2400 K, is well above $`\mathrm{\Delta }`$. For samples with higher cross–section areas we obtained somewhat larger activation energies, up to $`W(57)\times 10^3`$ K for the normal–sized samples, as it was reported earlier . Turning to the $`\sigma `$ vs. $`V`$ dependences (Fig. 2) we note that the onset of the non–linear conduction is smeared, the threshold rounding being clearly seen above $`T=120`$ K. The scale of the voltages applied corresponds to large fields, above 1 kV/cm. At lower temperatures (about 120 K) the onset of the non–linear conduction is relatively sharp, and we can estimate the threshold for collective conduction as $`V_T0.2`$ V ($`E_T400`$ V/cm), in accordance with the size effect . We shall assume that at higher temperatures the value of $`E_T`$ is approximately the same. Thus, the non–linear conduction for $`V<0.2`$ V will be referred to as subthreshold non–linear conduction, while at $`V`$ well above 0.2 V collective conduction is expected. Below we shall see that this division is not unphysical. One can see (Fig. 2) that the rounding progresses with approaching $`T_P`$. At high temperatures ($`T170`$ K) it is impossible to define a voltage range of linear conduction: the non–linearity starts from zero voltage. Fig. 3 shows the non–linear conductance, $`\sigma _{nl}\sigma (V)\sigma (0)`$, at fixed values $`V<V_T`$ as a function of $`T`$, together with $`\delta \sigma (T)`$. Evidently, $`\delta \sigma (T)`$ and $`\sigma _{nl}(T)`$ behave in similar ways up to $`T175`$ K, while at higher temperatures $`\sigma _{nl}`$ deviates downwards. One can conclude that the excess conductivity, $`\delta \sigma `$, and the threshold rounding have a common underlying mechanism at least at the lower temperatures. The possibility of coupling of the non-linear conduction below $`E_T`$ with an enhancement of low-field conductivity was also noticed in Ref. for thin samples of NbSe<sub>3</sub>. The scaling observed resembles that between the linear and non–linear conduction above $`V_T`$ known for different CDW conductors, including TaS<sub>3</sub> . Our samples also demonstrate such a scaling: the inset to Fig. 3 shows the dependence $`\sigma (T)`$ together with $`\sigma _{nl}(T)`$ at fixed $`V>V_T`$. In agreement with earlier observations, both values depend on temperature in a similar way, with the activation energy $`\mathrm{\Delta }`$. It is clear that the scaling between $`\sigma _{nl}(T)`$ ($`V<V_T`$) and $`\delta \sigma (T)`$ is quite different, as the slopes of the curves correspond to much higher activation energies, $`W\mathrm{\Delta }`$. ## III Discussion The large values of $`E_T`$ result from the small transversal dimensions of the samples, in accordance with the size effect observed in TaS<sub>3</sub> . As both transversal dimensions of our TaS<sub>3</sub> samples are of the same order of magnitude , we expect the pinning to be 1-dimensional, rather than 2-dimensional, as in the thin samples of NbSe<sub>3</sub> . Following the explanations of rounding in NbSe<sub>3</sub> , we could assume that the pinning energies of phase–correlation volumes in TaS<sub>3</sub> nanosamples are small enough to enable thermal depinning of phase–coherent volumes. The lowest-energy local depinning of the CDW results in a phase gain by the value $`2\pi `$ over the phase-correlation length . Such a deformation will cause a CDW stress, resulting in a local variation of resistance by percents for our samples . Meanwhile, metastable states cannot exist in such thin samples in the vicinity of $`T_P`$: the hysteresis loop develops only below 140 K in the representative sample (in usual-sized pure samples the hysteresis develops 5—7 K below $`T_P`$ ). So, any deformation immediately relaxes via a PS act, i.e. a plastic deformation of the CDW. The PS act is followed, sequentially, by local creep of the CDW . This results in a phase perturbation of the same order of magnitude as the initial elementary act of creep , and so PS is to be taken into account. Below we discuss in detail the conditions for the spontaneous PS and its effects on the conductivity. Remarkable is that the activation energy for the fluctuations is nearly independent of the field applied while it is below $`E_T`$: the slopes of the excess linear conductivity and of the non–linear conductivity at $`V=160`$ mV (which is quite close to $`V_T200`$ mV) are close, while the activation energy for the creep should run to zero at $`EE_T`$ . So the process initiating the fluctuation conduction is other than creep of the CDW. At the same time, at low temperatures ($`130`$ K) the non–linear conductivity becomes distinguishable only close to $`E_T`$, i.e. reveals itself as the threshold rounding. So $`E_T`$ is a characteristic field for the fluctuation conductivity, and the latter is in a way coupled to the CDW creep. This apparent contradiction is removed by the the following consideration. Evidently, the mechanism initiating the conductivity is the PS: the high activation energy is typical for PS in TaS<sub>3</sub> , and its independence of $`E`$ at $`T>120`$ K was reported in Ref. . At the same time according to Ref. each PS act is followed by temporary creep (rearrangement) of the CDW in the vicinity of the point where the PS occurred. In the presence of an external electric field the creep prevails in the direction defined by the field and provides a mechanism of the CDW conduction below $`E_T`$. At $`EE_T`$ the CDW phase-correlation length diverges , so $`E_T`$ is expected to be the critical point for the conduction. A hypothesis that phase slippage (in particular, edge dislocations) could facilitate CDW creep was also remarked in Ref. . In the case of 1-dimentional pinning (which could be applied to our samples) and PS involving the whole cross-section area, we can estimate the current induced by the PS. For simplicity let us consider the initial state to be uniform, i.e. the shift of the chemical potential $`\zeta =`$const. Entering of a new period in the absence of external field is followed by CDW creep under the internal electric fields $`E_{int}=d\zeta /dx`$. The creep proceeds while the effect of $`E_{int}`$ exceeds the effect of impurities, which we for simplicity describe by the average value, $`E_T`$. The resulting phase perturbation (Fig. 4) covers the length $$Ł_{2\pi }2\sqrt{\pi (d\zeta /dq)/E_T},$$ (1) where $`d\zeta /dq`$ characterizes the CDW elastic modulus, $`q`$ being the in-chain component of the CDW wave vector. Note, that $`L_{2\pi }`$ appears to be of the order of the phase-correlation length . Under an external electric field $`E<E_T`$ the creep proceeds asymmetrically with respect to the point of the PS nucleation giving the divergence of $`L_{2\pi }`$ at $`EE_T`$ . The new period is distributed so that the $`d\zeta /dx=E_T+E`$ from one side of the maximum remnant deformation and $`(E_TE)`$ from the other side (Fig. 4). The resulting progress of the CDW (and of the coupled charge $`2e`$ per chain) in the direction defined by $`E`$ could be estimated as $`\delta L=\frac{1}{3}(L_2L_1)`$, where $`L_1`$ and $`L_2`$ are the lengths of the phase perturbations in the two directions (Fig. 4) . With the condition that the areas under the triangles (Fig. 4) should be equal and correspond with the phase gain $`2\pi `$ we obtain from simple calculations: $$\delta L(E)=\frac{1}{3}L_{2\pi }\frac{E}{E_T}\frac{1}{\sqrt{1(E/E_T)^2}}.$$ (2) If the PS nucleation rate per unit length is $`\nu (T,E)`$, then the resulting mean current is $$I_{PS}=2e\nu \delta L$$ (3) per chain. As each PS act (fluctuator) affects the length $`L_{2\pi }`$, $`L_{2\pi }\nu f`$ may be considered as a typical frequency of switching of independent fluctuators. The temperature dependence of the PS rate could be empirically presented as $`\mathrm{exp}(W/T)`$ , where $`W(5`$ \- $`7)\times 10^3`$ K . So, Eqs. 2 and 3 give the dependence of the excess current both on $`T`$ and $`E`$. As $`EE_T`$ an unphysical divergence of $`I_{PS}`$ occurs, because in the model we have neglected the time of creep, $`\tau _{cr}`$, following each PS act. With approaching $`E_T`$ $`\tau _{cr}`$ grows together with $`L_2`$ (Fig. 4), and the PS frequency becomes dominated by $`1/\tau _{cr}`$. At low temperatures when $`f`$ is relatively small, $`I_{PS}`$ becomes noticeable only for $`E`$ close to $`E_T`$: Eqs. 2 and 3 thus feature the threshold rounding. At higher $`T`$ (and $`f`$) the area of validity of Eq. (3) shrinks down to smaller $`\|E\|`$. In the limit of small $`\|E\|`$ neglecting the dependence of $`\nu `$ on $`E`$ we obtain: $$I_{PS}=\frac{2}{3}ef\left[\frac{E}{E_T}+\frac{1}{2}\left(\frac{E}{E_T}\right)^3\right]I_l+I_{nl}.$$ (4) Thus, spontaneous PS gives contributions both to linear ($`I_l`$) and non–linear ($`I_{nl}`$) currents. Note that extrapolation of $`E`$ to $`E_T`$ gives a relation between $`I_{PS}`$ and $`f`$ that is very similar to that between the CDW current and the fundamental frequency. This is natural, because for $`EE_T`$ each pair of electrons entering the CDW via a PS act creeps along the whole sample. From Eq. (4) we obtain $$\sigma _{nl}=\frac{1}{2}\sigma _l\frac{E^2}{E_T^2},$$ (5) where $`\sigma _l`$ is the PS–induced part of the linear conductivity. At fixed $`E/E_T`$ ($`E<E_T`$) the fluctuation non–linear conductivity should have the same temperature dependence as the linear conductivity: both are dominated by $`f\mathrm{exp}(W/T)`$ \[Eq. (4)\]. Neglecting the dependence of $`E_T`$ on $`T`$ we come to the scaling between $`\delta \sigma `$ and $`\sigma _{nl}`$, as the one observed from Fig. 3. For a quantitative comparison of $`\delta \sigma `$ and $`\sigma _{nl}`$ note that $`[d\sigma /d(E^2)]\times 2E_T^2`$ is simply $`\sigma _l`$ ($`\delta \sigma `$) \[Eq. (5)\]. To check this, we show the value $`[d\sigma /d(E^2)]\times 2E_T^2`$ in Fig. 3, where $`d\sigma /d(E^2)`$ is determined from the best fit of $`\sigma (V)`$ ($`VV_T`$) with parabolic dependencies and $`E_T`$ is a fitting parameter. An example of the fit of $`\sigma (V)`$ by Eq. (5) is shown in Fig. 2 with a broken line. For the representative sample we get $`E_T=480`$ V/cm. For different samples, the values obtained from the fit by Eq. (5) agree with the results of direct measurements, though they are somewhat larger . Additional evidence of the correlation between the threshold field and non–linear fluctuation conduction is provided by the measurements of another sample with approximately the same length but larger cross-section area, $`s=1.5\times 10^3`$ $`\mu `$m<sup>2</sup>, and somewhat larger activation energy for $`\delta \sigma (T)`$, $`W=3400`$ K. A similar treatment of $`\sigma _{nl}`$ with the help of Eq. (5) gave us the dependence $`E_T(T)`$. In addition, we were able to measure $`E_T(T)`$ up to $`T=T_P`$ and higher directly, as the onset of sharp non–linear conduction. This was possible after subtracting the part of conductivity $`E^2`$ from each $`\sigma (V)`$ curve . The values of $`E_T`$ determined both ways are presented in Fig. 5 as a function of temperature. Both dependences show mesoscopic-type irregular variations of $`E_T`$ with temperature , though $`E_T`$ obtained from Eq. (5) is about 5 times larger . One can see correlation between the two dependences. From Eq. (4) we obtain reasonable estimates of the frequencies of switching for the fluctuators. To observe distinct excess conductivity (at 140 K for the representative sample, Fig. 3), we should take $`f10^7`$ Hz. This is only 2—3 orders of magnitude higher than we were able to see directly from the time domains of the fluctuations , the latter being restricted by the electric scheme. Thus, both the values and the activation energy for the linear and non–linear fluctuation currents are fairly described by the model proposed. It is worth mentioning that the dependence of $`\|d^3V/dI^3\|`$ vs $`T`$ below $`T_P`$ presented in Ref. fits the Arrhenius law fairly well with $`W4500`$ K, in agreement with the PS measurements at the contacts . The lowering of $`W`$ with the decrease of the sample’s cross-section area also finds natural explanation within the model proposed. In fact, a large threshold field corresponds to high inhomogeneous stress of the CDW in the thin samples : $$<\zeta ^2>^{1/2}2\sqrt{\pi E_T(d\zeta /dq)},$$ (6) The stress lowers the barrier for the PS in certain points . The decrease of the sample cross-section area reveals itself in the growth of $`E_T`$, and thus results in the lowering of the activation energies for $`\delta \sigma (T)`$ and $`\sigma _{nl}(T)`$. Earlier we have explained in a similar way the broadening of the range of the fluctuations and of the Peierls transition along the temperature scale in the thin samples . Note that the model of thermal depinning of the phase-coherent volumes gives a much stronger size dependence of the excess conductivity: the depinning energy is proportional to $`s^{2/3}`$ (Ref. ). So, for the sample with $`s=1.5\times 10^3`$ $`\mu `$m<sup>2</sup> we should expect $`W`$ to be about 7000 K (instead of 3400 K), and the excess conductivity should become negligible in the thick samples. Actually, we found no qualitative difference between the excess conductivity of the thick and thin samples. The activation energy for $`\delta \sigma (T)`$ in thick samples is $`(57)\times 10^3`$ K , in agreement with the activation energy found from the noise probing of the spontaneous PS . The threshold rounding is also noticeable in thick samples , though the study of the non–linear fluctuation conductivity is complicated because of its narrow temperature range and small $`E_T`$. Note that the dependence $`\delta \sigma (T)`$ follows the activation law up to $`T_P`$, and even a little bit above it (Fig. 3), no feature being observed at $`T_P`$. So, the state a little bit above $`T_P`$ could be considered as a CDW saturated with climbing dislocations rather than a single–electron state. The conduction of such a mixture is supplied by the processes of nucleation and motion of the domain boundaries, which exert high internal electric fields to the domains. The fact that the dependencies characterizing the non–linear conductivity deviate from the Arrhenius law at lower temperatures than $`\delta \sigma (T)`$ could be ascribed to the growth of $`E_T`$ near $`T_P`$ (Fig. 5); note also that with increasing $`T`$ the model fails first at finite $`E`$, and then at $`E0`$. In conclusion, we have observed the fluctuation contribution to the conductivity of thin samples of o-TaS<sub>3</sub>, which comprises linear and non–linear parts. We have shown that the spontaneous phase slippage observed in the CDW in the vicinity of $`T_P`$ results to the excess conductivity whose temperature and electric-field dependences match our experimental observations. The simple model proposed requires further development. In particular, the mechanism of PS nucleation and evolution should be considered in terms of nucleation and propagation of dislocation loops in the CDW . A possible contribution of glide of the dislocation lines to the current also requires analysis. In the case of bulk (3D) samples the loops cover only part of the cross-section areas, so transversal interaction of the chains while the local creep proceeds should be considered. A special case is the ribbon–like (2D) samples evidently treated in Refs. . A dislocation loop degenerates into a pair of points interacting logarithmically. Then the system acquires the features of a 2D crystal exhibiting the Kosterlitz-Thouless transition . This approach can explain the lowering of $`T_P`$ in thin crystals and gradual power–law dependences of $`\sigma _{nl}`$ on $`(T)`$. ## IV Acknowledgments We are grateful to P. Monceau for help in the experiment, to Yu. I. Latyshev, Ya. S. Savitskaya, and V. V. Frolov for producing the samples, and to S.N. Artemenko and A.A. Sinchenko for fruitful discussions. This work was supported by the Russian Foundation for Basic Research (grants Nos. 98-02-16667, 99-02-17387), Jumelages 19 (RFBR, grant No. 98-02-22061), and MNTP “Physics of Solid State Nanostructures” (grant No. 97-1052).
warning/0004/astro-ph0004105.html	ar5iv	text	# Cosmic Strings Lens Phenomenology: General Properties of Distortion Fields ## I Introduction It is well known that topological defects may appear whenever, in the thermal history of the universe, a symmetry breaking phase transition occurs , as for instance in grand–unified theories or in some extensions of the standard electro-weak model . Such defects represent spacetime positions where the underlying order parameter cannot relax, because of topological constraints, to its low energy vacuum state . They are expected to interact mainly gravitationally with the ordinary matter so that they can induce (i) deflection and redshifting of massless particles, (ii) accretion of massive non–relativistic particles and (iii) emission of gravitational waves (see e.g. for a review of these effects). One very interesting example, from a high energy physics point of view as well as from a cosmological point of view, corresponds to the case of cosmic strings. In this case, their phenomenological properties are determined by the energy density per unit length of a string, $`U`$. For instance the deflection angle is of order $`4\pi GU`$, $`G`$ being the Newton constant. For defects formed at the grand–unification scale ($`T_{\mathrm{GUT}}10^{16}\mathrm{Gev}`$), we expect effects of a magnitude of $`GU10^6`$. This corresponds for instance to the magnitude of the cosmic microwave background (CMB) anisotropies induced through the Kaiser-Stebbins effects . Although topological defects may have clear observational signature on the CMB sky , observations seem to disfavor such an origin . Nonetheless, it does not mean that topological defects do not exist. Their detection would be of dramatic importance both for astronomy and particle physics since for instance estimation or bounds on their density will help constraining the high energy physics theories predicting their existence. Definitive predictions for string properties are however difficult to obtain, because in particular of the complex evolution equations that may depend on their microscopic structure through their equation of state. For the so-called Goto–Nambu strings (where the energy per unit length $`U`$ and the tension $`T`$ of the string are equal), it was shown that the spacetime around such straight cosmic string was conical . Such a cosmic string formed at GUT scale would therefore induce image pairs of distant objects with angular separation of order $`\theta 5.2^{\prime \prime }\times (GU/10^6)`$. From a pure phenomenological point of view such a string is expected to produce lines of double images . Recognizing the peculiarity of such a system, it was later extended to straight cosmic string with different equation of state and to a string with a lightlike current pulse. Moreover numerical simulations for Goto–Nambu string were also performed in the case of long strings and cosmic loops . More quantitatively the prospects for a direct detection of relic string via gravitational lensing, and in particular the expected number of events, was first discussed by Hindmarsh who estimated the angular length of string per unit area on the sky out to a redshift $`z`$ to be of order $$\theta _{\mathrm{loops}}0.1\nu z^2\mathrm{deg}^1\theta _{\mathrm{long}\mathrm{string}}0.1Az^2\mathrm{deg}^1$$ where $`\nu `$ and $`A`$ are two coefficients of order unity (see also ). In conclusion, it is widely believed that the observation of a cosmic string can be achieved through double image detections, although, in practice, it might be difficult to be positive about such a detection since pairs with the same angular separation appearing by pure coincidence can be very high (as pointed by Cowie and Hu who reported for such a cosmic string lens candidate ). The aim of this article and its companion is to make a systematic investigation of the gravitational lensing effects by cosmic strings. We focus our analysis on the deformation equation of a geodesic bundle in presence of cosmic strings (§ II). Standard approximations of the gravitational lens theory are also discussed in this section. After a description of the cosmic string dynamics in § III, we show in § IV that the deflecting potential of a cosmic string is equivalent to the one by a static distribution of matter on the projection of the string worldsheet onto the observer past light cone. In the course in this calculation we show that the deformation field induced by cosmic strings has a zero convergence (without any approximation). Examples illustrating these results are discussed as well as the validity of the thin lens approximation and the influence of the equation of state. In § V we investigate the phenomenological consequences of the zero convergence property on multiple image systems. In a companion paper, , we propose a phenomenological model of string energy distribution that gives a more quantitative account of these results. ## II Evolution of a light beam In lens systems that are usually considered in cosmology, such as galaxies or galaxy clusters, the metric perturbations correspond to those of scalar perturbations. This is not the case for cosmic string effects where relativistic motions, non trivial equation of state, also induce vector and tensor perturbations. We are thus forced to consider the deformation equations of light beams in their full generality. In the geometric optic approximation, an electromagnetic plane wave propagating in an arbitrary spacetime $``$ without interaction with matter can be described by a null geodesic . The goal of this section is to review the description of the evolution and distortion of a bundle of null geodesics and we start by introducing the standard elements of the gravitational lensing theory and then apply them to a perturbed spacetime. We then discuss the thin lens approximation and finish by some comments on its applicability. ### A Basics of gravitational lensing We consider a bundle of null geodesics $`g`$ propagating in a spacetime $``$. Each geodesic can be described as $$g:x^\mu (\lambda )=\overline{x}^\mu (\lambda )+\xi ^\mu (\lambda )$$ (1) where $`\overline{x}^\mu (\lambda )`$ is the equation of a geodesic $`g_0`$ chosen as reference and $`\xi ^\mu `$ is a displacement vector labeling the other geodesics with respect to $`g_0`$. Greek indices run from 0 to 3 and $`\lambda `$ is an affine parameter along the geodesic $`g_0`$. With these notations, we can define the tangent vector to $`g_0`$ by $$k^\mu \frac{\mathrm{d}\overline{x}^\mu }{\mathrm{d}\lambda }.$$ (2) It is a null vector satisfying the geodesic equation, i.e solution of $$k_\mu k^\mu =0,k^\mu _\mu k^\nu =0,$$ (3) where $`_\mu `$ is the covariant derivative associated to the metric $`g_{\mu \nu }`$ the signature of which will be chosen as $`(,+,+,+)`$. Now, we consider such a bundle converging at a point $`O`$ where we assume that there is an observer with 4–velocity $`u^\mu `$ ($`u^\mu `$ is a timelike vector, i.e. such that $`u^\mu u_\mu =1`$) and we choose the affine parameter $`\lambda `$ to vanish at $`O`$ and to increase toward the past. We then consider at $`O`$ the basis $`(k^\mu ,u^\mu ,n_1^\mu ,n_2^\mu )`$ where $`n_{1,2}^\mu `$ are two spacelike vectors ($`n_a^\mu n_{a\mu }=+1`$, $`a=1,2`$) such that $$n_1^\mu n_{2\mu }=0,k_\mu n_a^\mu =0\text{and}u_\mu n_a^\mu =0$$ (4) and $`k^\mu `$ is the null vector defined in (2). Starting from this basis at $`O`$, we construct such a basis at any point of the light ray worldline $`\overline{x}^\mu `$ by parallelly transporting it as $$k^\mu _\mu X^\nu =0$$ (5) for $`X=u,n_1`$ and $`n_2`$. Note that since $`k^\mu `$ satisfies (3) this implies that (5) is in fact the Fermi–Walker transport and thus $`u`$, $`n_1`$ and $`n_2`$ remain orthonormal and satisfy (4) for all $`\lambda `$. Since the tangent vector $`k_g^\mu k^\mu +\mathrm{d}\xi ^\mu /\mathrm{d}\lambda `$ to each geodesic $`g`$ of the bundle is a null vector, we deduce from $`k_g^\mu k_{g\mu }=g_{\mu \nu }(\overline{x}^\alpha +\xi ^\alpha )k_g^\mu k_g^\nu =0`$ that $`2k_\mu \mathrm{d}\xi ^\mu /\mathrm{d}\lambda +k^\mu k^\nu \xi ^\alpha _\alpha g_{\mu \nu }=0`$ at first order in $`\xi `$. It can then be concluded that, using (5), $`k_\mu \xi ^\mu `$ is constant along the geodesic and vanishes at $`O`$ so that it can be decomposed as $$\xi ^\mu =\xi _0k^\mu +\underset{a=1,2}{}\xi _an_a^\mu .$$ (6) $`\xi _0`$ does not vanish in general, but two such decompositions (6) with different $`\xi _0`$ parameterize the same light ray. We can for instance impose that $`\xi ^\mu `$ is spatial for the observer with 4-velocity $`u^\mu `$ (i.e. $`k_\mu u^\mu =0`$) which then fixes the value of $`\xi _0`$. We also decompose the coordinates of every event of $``$ in the neighborhood of $`g_0`$ as $$x^\mu =\lambda k^\mu +\underset{a=1,2}{}x_an_a^\mu +\tau u^\mu .$$ (7) The equation of evolution of $`\xi ^\mu `$ is obtained by writing the geodesic deviation equation $$\frac{\mathrm{D}^2}{\mathrm{d}^2\lambda }\xi ^\mu =R_{\nu \alpha \beta }^\mu k^\nu k^\alpha \xi ^\beta $$ (8) where $`R_{\mu \nu \alpha \beta }`$ is the Riemann tensor of the metric $`g_{\mu \nu }`$ and where $`\mathrm{D}/\mathrm{d}\lambda k^\nu _\nu `$. Inserting the decomposition (6) in (8) and using the fact that $`\xi _a=n_a^\mu \xi _\mu `$ is a scalar (so that $`\mathrm{D}\xi _a/\mathrm{d}\lambda =\mathrm{d}\xi _a/\mathrm{d}\lambda `$ with $`\mathrm{d}/\mathrm{d}\lambda k^\mu _\mu `$) leads to $$\ddot{\xi }_a=_a^b\xi _b$$ (9) where $`_{ab}R_{\mu \nu \alpha \beta }k^\nu k^\alpha n_a^\mu n_b^\beta `$ and a dot refers to a derivation with respect to $`\lambda `$. Due to the linearity of the geodesic deviation equation (9), $`\xi _a`$ can be related to its initial value $`\dot{\xi }_a(0)`$ through a linear transformation $`𝒟_{ab}`$ as $$\xi _a(\lambda )𝒟_a^b(\lambda )\dot{\xi }_b(0).$$ (10) Since $`\xi (0)=0`$ for a bundle converging at $`O`$, with two derivatives (10) and using this equation again to eliminate $`\dot{\xi }(0)`$, we obtain the equation of evolution for $`𝒟_{ab}`$ $$\ddot{𝒟}_{ab}=_a^c𝒟_{cb}$$ (11) with initial conditions $$𝒟_{ab}(0)=0\text{and}\dot{𝒟}_{ab}(0)=I_{ab},$$ (12) $`I_{ab}`$ being the $`2\times 2`$ identity matrix. This equation has been derived in e.g. . $`𝒟_{ab}`$ characterizes the deformation field while looking in different directions. Quoting that $`\dot{\xi }(0)\theta _\mathrm{I}`$ is the vectorial angle of observation and $`\xi (\lambda _\mathrm{S})\lambda _\mathrm{S}\theta _\mathrm{S}`$ where $`\theta _\mathrm{S}`$ is the vectorial angular position of the source, equation (10) can be rewritten in terms of these angles (see figure 1) as $$\theta _\mathrm{S}^a=\frac{𝒟_b^a(\lambda _\mathrm{S})}{\lambda _\mathrm{S}}\theta _\mathrm{I}^b.$$ (13) The amplification matrix $`𝒜_b^a\mathrm{d}\theta _\mathrm{S}^a/\mathrm{d}\theta _\mathrm{I}^b`$ can be expressed in terms of $`𝒟`$ as $$𝒜_b^a=\frac{𝒟_b^a(\lambda _\mathrm{S})}{\lambda _\mathrm{S}}.$$ (14) In the following, we decompose it in terms of convergence $`\kappa `$ and shear $`\stackrel{}{\gamma }(\gamma _1,\gamma _2)`$ as $$𝒜_{ab}=\left(\begin{array}{cc}1\kappa \gamma _1& \gamma _2\\ \gamma _2& 1\kappa +\gamma _1\end{array}\right).$$ (15) ### B Application to a perturbed spacetime We now restrict our study to a perturbed spacetime of metric $$\mathrm{d}s^2=g_{\mu \nu }\mathrm{d}x^\mu \mathrm{d}x^\nu (\eta _{\mu \nu }+h_{\mu \nu })\mathrm{d}x^\mu \mathrm{d}x^\nu ,$$ (16) with $`\eta _{\mu \nu }`$ being the Minkowski metric $`\eta _{\mu \nu }=\text{diag}(,+,+,+)`$. We work in harmonic gauge $$_\nu h^{\mu \nu }=0,$$ (17) so that the Ricci tensor at linear order in the perturbation $`h_{\mu \nu }`$ reduces to $$R_{\mu \nu }=\frac{1}{2}(_t^2\mathrm{\Delta })h_{\mu \nu },$$ (18) $`\mathrm{\Delta }`$ being the Laplacian $`\mathrm{\Delta }_i^i`$, Latin indices running from 1 to 3. The Einstein equations take the simple form $$(_t^2\mathrm{\Delta })h_{\mu \nu }=16\pi G(T_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }T_\lambda ^\lambda )16\pi G_{\mu \nu },$$ (19) $`T_{\mu \nu }`$ being the stress-energy tensor of the matter perturbation. The solution of this equation can be obtained by means of the Green functions, $`𝒢^{(\pm )}`$, of the d’Alembertian $$\left(\mathrm{\Delta }_\stackrel{}{x}_t^2\right)𝒢^{(\pm )}(\stackrel{}{x}^{},t^{},\stackrel{}{x},t)=\delta ^{(3)}(\stackrel{}{x}\stackrel{}{x}^{})\delta (tt^{})𝒢^{(\pm )}(\stackrel{}{x}^{},t^{},\stackrel{}{x},t)=\frac{1}{4\pi }\frac{\delta (t^{}t\pm \left\|\stackrel{}{x}\stackrel{}{x}^{}\right\|)}{\|\stackrel{}{x}\stackrel{}{x}^{}\|},$$ (20) so that, using the retarded solution, we solve the Einstein equations (18) at linear order as $$h_{\mu \nu }(\stackrel{}{x},t)=4G\frac{\mathrm{d}^3\stackrel{}{x}^{}}{\|\stackrel{}{x}\stackrel{}{x}^{}\|}_{\mu \nu }(\stackrel{}{x}^{},t\left\|\stackrel{}{x}\stackrel{}{x}^{}\right\|).$$ (21) Now, we can solve (11) order by order: at zeroth order, $`_{ab}=0`$ so that it reduces trivially to $$𝒟_{ab}^{(0)}(\lambda )=\lambda I_{ab};$$ (22) at first order, the equation of evolution (11) gives $$\ddot{𝒟}_{ab}^{(1)}(\lambda )=\lambda _{ab}^{(1)}(\lambda )$$ (23) from which we deduce that $`𝒟_{ab}^{(1)}(\lambda )`$ is given by $`𝒟_{ab}^{(1)}(\lambda _\mathrm{S})`$ $`=`$ $`{\displaystyle _0^{\lambda _\mathrm{S}}}\lambda (\lambda _\mathrm{S}\lambda )_{ab}^{(1)}(\lambda )d\lambda .`$ (24) It is interesting to note that while $`𝒟_{ab}`$ is not symmetric in general, it is symmetric at first order in the perturbation. Using the expression of the Riemann tensor as $`2R_{\mu \sigma \nu \rho }=h_{\sigma \nu ,\mu \rho }+h_{\mu \rho ,\sigma \nu }h_{\nu \mu ,\sigma \rho }h_{\sigma \rho ,\nu \mu }`$, we get that $$_{ab}^{(1)}=\frac{1}{2}h_{\nu \sigma ,\mu \rho }k^\nu k^\sigma n_a^\mu n_b^\rho \frac{1}{2}\frac{\mathrm{d}}{\mathrm{d}\lambda }\left(\mathrm{\Gamma }_{\rho \beta }^\alpha \eta _{\alpha \mu }k^\beta n_a^\mu n_b^\rho \right)$$ (25) where $`k^\mu `$, $`n_1^\mu `$ and $`n_2^\mu `$ are evaluated on the unperturbed geodesic (and are thus constant) and where $`\mathrm{\Gamma }_{\rho \beta }^\alpha `$ are the Christoffel symbols at first order in the perturbation. Choosing $$n_a^\mu \delta _a^\mu $$ (26) and defining the deflecting potential $`\mathrm{\Phi }`$ as $$\mathrm{\Phi }(\stackrel{}{x},t)\frac{1}{2}h_{\mu \nu }k^\mu k^\nu ,$$ (27) where $`t(\lambda )=t_0x_{}(\lambda )`$ is the equation of the photon trajectory \[$`t_0`$ is the time of the observation $`t_0=t(\lambda =0)`$\], we obtain that $$𝒟_{ab}^{(1)}(x_a,\lambda _\mathrm{S})=\left(_0^{\lambda _\mathrm{S}}\lambda (\lambda _\mathrm{S}\lambda )_{ab}\mathrm{\Phi }(x_a,x_{}(\lambda ),t(\lambda ))\mathrm{d}\lambda \right)+_0^{\lambda _\mathrm{S}}\lambda (\lambda _\mathrm{S}\lambda )\mathrm{\Psi }_{ab}(x_a,x_{}(\lambda ),t(\lambda ))d\lambda $$ (28) where $$\mathrm{\Psi }_{ab}(\stackrel{}{x},t)\frac{1}{2}\frac{\mathrm{d}}{\mathrm{d}\lambda }\left(\mathrm{\Gamma }_{a\beta }^\alpha \eta _{\alpha b}k^\beta \right)$$ (29) and where $`_a`$ refers to a derivative with respect to the coordinates $`x_a`$ as defined in (7). ### C The thin lens approximation for static distribution of matter In the thin lens approximation, one assumes that the effect of the deflecting body takes place over only a small fraction of the light path. This approximation is usually considered in cases of scalar perturbations. The aim of this paragraph is to recall its derivation in the standard case to see to which extent it applies for cosmic strings. We thus assume that the lens is localised at $`\lambda =\lambda _\mathrm{L}`$ with an extension small compared to $`\lambda _\mathrm{L}`$ and that this matter distribution is static so that $`_{\mu \nu }k^\mu k^\nu \mathrm{\Sigma }(\stackrel{}{x}_{},x_{}(\lambda _\mathrm{L}))\delta (x_{}(\lambda )x_{}(\lambda _\mathrm{L}))`$, where $`\mathrm{\Sigma }`$ is the surface energy density. It follows that the deflecting potential reduces, after integration over the direction of propagation, to $$\mathrm{\Phi }(\stackrel{}{x}_{},x_{})=2G\frac{\mathrm{d}^2\stackrel{}{x}_{}^{}}{\sqrt{\|\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\|^2+(x_{}x_{}(\lambda _\mathrm{L}))^2}}\mathrm{\Sigma }(\stackrel{}{x}_{}^{},x_{}(\lambda _\mathrm{L}))$$ (30) where $`\stackrel{}{x}_{}(x_1,x_2)`$. Since only $`_{ab}\mathrm{\Phi }`$ enters the expression of $`𝒟_{ab}^{(1)}`$ and since this quantity varies as $`(x_{}x_{}(\lambda _\mathrm{L}))^3`$ as soon as we are looking close to the string \[i.e. when $`\|\stackrel{}{x}_{}\stackrel{}{x}_{,\mathrm{L}}\|(x_{}x_{}(\lambda _\mathrm{L}))`$\] and we can approximate the deflecting potential as $$\mathrm{\Phi }(\stackrel{}{x})=\stackrel{~}{\mathrm{\Phi }}(\stackrel{}{x}_{})\delta (x_{}x_{}(\lambda _\mathrm{L}))$$ (31) with $`\stackrel{~}{\mathrm{\Phi }}(\stackrel{}{x}_{})`$ $``$ $`{\displaystyle _0^{\lambda _\mathrm{S}}}\mathrm{\Phi }(\stackrel{}{x}_{},x_{}(\lambda ))d\lambda `$ (32) $`=`$ $`2G{\displaystyle \mathrm{\Sigma }(\stackrel{}{x}_{}^{},x_{}(\lambda _\mathrm{L}))\left[\mathrm{ln}\left(x_{}(\lambda )x_{}(\lambda _\mathrm{L})+\sqrt{(\stackrel{}{x}_{}\stackrel{}{x}_{}^{})^2+(x_{}(\lambda )x_{}(\lambda _\mathrm{L}))^2}\right)\right]_0^{\lambda _\mathrm{S}}\mathrm{d}^2\stackrel{}{x}_{}^{}}.`$ (33) Now, if we assume that the impact parameter is small compared to the two distances lens–object and observer–lens, i.e. $$\|\stackrel{}{x}_{}\stackrel{}{x}_{,\mathrm{L}}\|(x_{}(\lambda _\mathrm{S})x_{}(\lambda _\mathrm{L}),x_{}(\lambda _\mathrm{L})),$$ (34) we deduce that the deflecting potential integrated along the line of sight is given by $$\stackrel{~}{\mathrm{\Phi }}(\stackrel{}{x}_{})=4G\mathrm{ln}\left\|\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right\|\mathrm{\Sigma }(\stackrel{}{x}_{}^{},x_{}(\lambda _\mathrm{L}))\mathrm{d}^2\stackrel{}{x}_{}^{}$$ (35) up to a constant which depends only on $`x_{}(\lambda _\mathrm{L})`$ and $`x_{}(\lambda _\mathrm{S})`$; we forget this constant since $`𝒟_{ab}`$ involving only derivatives of $`\stackrel{~}{\mathrm{\Phi }}`$ and is thus independent of its value. The second contribution of $`𝒟_{ab}^{(1)}`$ involves the computation of the potential $`\mathrm{\Psi }_{ab}`$ and one can show from (29) that if we deal only with scalar perturbations (i.e. such that $`h_{00}=2\varphi `$ and $`h_{ij}=2\psi \delta _{ij}`$) then $`\mathrm{\Psi }_{ab}=0`$. For the vector and tensor perturbations, $`\mathrm{\Psi }_{ab}`$ does not vanish but in the thin lens approximation its contribution corresponds to a boundary term (time dependent but identical for all light rays joining the source and the observer) which can thus be dropped. Then, the amplification matrix, in the thin lens approximation, reduces to $$𝒜_{ab}=I_{ab}4G\frac{\lambda _\mathrm{S}\lambda _\mathrm{L}}{\lambda _\mathrm{S}}\lambda _\mathrm{L}_a_b\mathrm{ln}\left\|\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right\|\mathrm{\Sigma }(\stackrel{}{x}_{}^{})\mathrm{d}^2\stackrel{}{x}_{}^{}$$ (36) which rewrites, after the change of variables $`\stackrel{}{\theta }^{}=\stackrel{}{x}_{}^{}/\lambda _\mathrm{S}`$, as $$𝒜_{ab}(\lambda _\mathrm{S})=I_{ab}4G\frac{\lambda _\mathrm{S}\lambda _\mathrm{L}}{\lambda _\mathrm{S}}_{\theta _\mathrm{I}^a}_{\theta _\mathrm{I}^b}\mathrm{ln}\left\|\stackrel{}{\theta }_\mathrm{I}\stackrel{}{\theta }^{}\right\|\lambda _\mathrm{L}\mathrm{\Sigma }(\stackrel{}{\theta }^{})\mathrm{d}^2\stackrel{}{\theta }^{}.$$ (37) Decomposing $`\mathrm{\Sigma }(\stackrel{}{\theta })`$ as $`\lambda _\mathrm{L}\mathrm{\Sigma }(\stackrel{}{\theta })=\mu (s)\delta \left(\stackrel{}{\theta }\stackrel{}{\theta }_{\mathrm{string}}(s)\right)ds`$ where $`\stackrel{}{\theta }_{\mathrm{string}}`$ represents the locus of the string on the plane $`x_{}=x_{}(\lambda _\mathrm{L})`$, we get that (37) reduces to $$𝒜_{ab}(\lambda _\mathrm{S})=I_{ab}_{\theta _\mathrm{I}^a}_{\theta _\mathrm{I}^b}\phi (\stackrel{}{\theta }_\mathrm{I})\text{with}\phi (\stackrel{}{\theta }_\mathrm{I})4G\frac{\lambda _\mathrm{S}\lambda _\mathrm{L}}{\lambda _\mathrm{S}}\mathrm{ln}\left\|\stackrel{}{\theta }_\mathrm{I}\stackrel{}{\theta }_{\mathrm{string}}(s)\right\|\mu (s)ds$$ (38) where $`\mu (s)`$ is the projected total lineic energy density of the string (which mixes the effect of the lineic energy, the tension and the currents along the string if any) and $`\phi `$ is the effective projected potential. When dealing with topological defects, there are different reasons why such an approximation may not hold. First the strings are extended and move with relativistic speed so that (i) they are a priori not confined in a plane $`\lambda _\mathrm{L}`$constant and (ii) one cannot assume that the distribution of matter of the lens is static so that the time dependence in the line-of-sight integration in (24) has to be taken into account. These issues will be addressed in § IV.B after a description of the general stress–energy tensor of strings (§ III). ### D Comments #### 1 Gravitational potential and deflecting potential As a first comment, let us stress that in general the deflecting potential $`\mathrm{\Phi }`$ is different from the Newtonian gravitational potential. For instance, a general perturbed spacetime has the general post–Newtonian metric $$\mathrm{d}s^2=(12\varphi )\mathrm{d}t^2+2A_i\mathrm{d}x^i\mathrm{d}t+\left[(1+2\psi )\delta _{ij}+2\overline{E}_{ij}\right]\mathrm{d}x^i\mathrm{d}x^j$$ (39) where $`\varphi `$ and $`\psi `$ are the Newtonian potentials, $`A_i`$ and $`\overline{E}_{ij}`$ are the vector (rotational) and tensor (gravitational waves) perturbations satisfying $$\overline{E}_i^i=_i\overline{E}^{ij}=_iA^i=0.$$ (40) It follows that $$\mathrm{\Phi }=\varphi +\psi +A_i\gamma ^i+\overline{E}_{ij}\gamma ^i\gamma ^j$$ (41) where $`\gamma ^i`$ is the direction of observation. This includes effects from the rotation of the deflecting body and of gravitational waves. Indeed, in the case of pure scalar perturbations, we recover that $$\mathrm{\Phi }=2\varphi .$$ In the case of scalar perturbations, one can easily check that $`\mathrm{\Psi }_{ab}=0`$ but topological defects also generate vector and tensor perturbations. In the thin lens approximation, the contribution of these two terms reduces to a boundary term that can be neglected but in the general case of extended object, one has to check that it is still the case for the vector and tensor modes. #### 2 Deflection angle In a general spacetime, the deflection is not straightforward to define. This is for instance the case in a perturbed spacetime with perturbations on all scales (see e.g. for a discussion and a generalization of this concept). If the matter perturbation causing the lensing is localized in space then the metric perturbations generally die away and the spacetime is asymptotically unperturbed. In that case, one can compute the deflection angle simply by solving the geodesic equation (3) at first order in the perturbations. For that purpose, we decompose the tangent vector to the geodesic $`g`$ as $$k_g^\mu =\overline{k}^\mu +\delta k^\mu .$$ (42) $`\overline{k}^\mu `$ is the tangent vector of the light ray in the unperturbed Minkowski spacetime . Note that $`\overline{k}^\mu `$ is different from the vector $`k^\mu `$ defined in (2) which labels the geodesic of reference $`g_0`$. At first order in the perturbation, it is in fact possible to choose the unperturbed geodesic as reference since the displacement $`\xi `$ is first order in the perturbation. Since both $`\eta _{\mu \nu }\overline{k}^\mu \overline{k}^\nu =0`$ and $`g_{\mu \nu }k^\mu k^\nu =0`$, we deduce that $`\overline{k}^\mu \delta k_\mu =0`$ at first order in the perturbations. At linear order, the geodesic equation (3) implies that $$\frac{\mathrm{d}}{\mathrm{d}\lambda }\delta k^\alpha +\delta \mathrm{\Gamma }_{\mu \nu }^\alpha \overline{k}^\mu \overline{k}^\nu =0$$ (43) with $`\delta \mathrm{\Gamma }_{\mu \nu }^\alpha \frac{1}{2}\eta ^{\alpha \beta }(h_{\beta \mu ,\nu }+h_{\beta \nu ,\mu }h_{\mu \nu ,\beta })`$, so that $$\left[\delta k^\alpha \right]_{\lambda _\mathrm{S}}^0=\eta ^{\alpha \nu }\left[h_{\mu \nu }\overline{k}^\mu \right]_{\lambda _\mathrm{S}}^0+\frac{1}{2}\eta ^{\alpha \beta }_{\lambda _\mathrm{S}}^0h_{\mu \nu ,\beta }\overline{k}^\mu \overline{k}^\nu d\lambda ,$$ (44) the integral being performed along the unperturbed geodesic. Forgetting the boundary term (which is a time dependent term but identical for all light rays joining the source and the observer), we can extract the variation of the the photon energy measured by an observer with velocity $`u^\mu `$ $$\delta E=_{\lambda _\mathrm{S}}^0_t\mathrm{\Phi }(\stackrel{}{x}_{},x_{}(\lambda ),t(\lambda ))\mathrm{d}\lambda $$ (45) and the deflection $`\stackrel{}{\alpha }`$ in the plane perpendicular to the line of sight $$\stackrel{}{\alpha }=\stackrel{}{}_{}_{\lambda _\mathrm{S}}^0\mathrm{\Phi }(\stackrel{}{x}_{},x_{}(\lambda ),t(\lambda ))d\lambda .$$ (46) The effect on the photon energy, and thus on its redshift, is nothing else but the well known Sachs–Wolfe effect . Focusing on (46), since $`\mathrm{\Phi }`$ is evaluated along the photon geodesic $`t(\lambda )=t_0x_{}(\lambda )`$, we deduce that $$\left[_t+_x_{}\right]\mathrm{\Phi }(\stackrel{}{x}_{},x_{}(\lambda ),t(\lambda ))=0$$ (47) along the photon path. Hence, rewriting the three dimensional Laplacian as $`\mathrm{\Delta }=_x_{}^2+\mathrm{\Delta }_{}`$ where $`\mathrm{\Delta }_{}`$ is the two dimensional Laplacian and using the Einstein equations (19), we deduce that $$\stackrel{}{}_{}.\stackrel{}{\alpha }=8\pi G_{\lambda _\mathrm{S}}^0(\stackrel{}{x}_{},x_{}(\lambda ),t(\lambda ))d\lambda =8\pi G_0^{t(\lambda _\mathrm{S})t_0}(\stackrel{}{x}_{},x_{}(t),t)dt$$ (48) after choosing the parameterization $`\lambda =t_0t`$ and where we have introduced $`_{\mu \nu }k^\mu k^\nu =T_{\mu \nu }k^\mu k^\nu `$. It follows that the 2–divergence of the deflection depends only on the projection of $`_{\mu \nu }`$ onto the photon trajectory in between the source and the observer and thus vanishes on all directions which do not intersect the string worldsheet. Note that such a result holds for any relativistic and/or extended lens. Now, in the thin lens approximation, one can relate $`\theta _\mathrm{S}`$ and $`\theta _\mathrm{I}`$ (see figure 1) by $$\theta _\mathrm{S}^a=\theta _\mathrm{I}^a\frac{\lambda _\mathrm{S}\lambda _\mathrm{L}}{\lambda _\mathrm{S}}\alpha ^a.$$ (49) The amplification matrix being given by $`𝒜_b^a=\mathrm{d}\theta _\mathrm{S}^a/\mathrm{d}\theta _\mathrm{I}^b`$, we obtain that $$𝒜_{ab}=I_{ab}+\frac{\lambda _\mathrm{S}\lambda _\mathrm{L}}{\lambda _\mathrm{S}}_{\theta ^b}_a_0^{\lambda _\mathrm{S}}\mathrm{\Phi }d\lambda .$$ (50) Since the lens is localised in planes close to $`\lambda =\lambda _\mathrm{L}`$, we have that $`_{\theta ^a}\lambda _\mathrm{L}_a`$ so that the former expression of the amplification matrix reduces to (36) once we use the expression (35) for the integrated potential. In both approaches, we find that the deflection angle is given by the usual expression (see e.g. ) $$\stackrel{}{\alpha }=4G\lambda _\mathrm{L}\frac{\stackrel{}{\theta }^\mathrm{I}\stackrel{}{\theta }^{}}{\|\stackrel{}{\theta }^\mathrm{I}\stackrel{}{\theta }^{}\|^2}\mathrm{\Sigma }(\theta ^{})\mathrm{d}^2\stackrel{}{\theta }^{}$$ and we emphasize that defining the amplification matrix through the deflection angle implicitely assumes that we are in the thin lens approximation (as seen for instance on equation (49)). #### 3 Applicability to cosmology It has probably not escaped a careful reader that we have restricted our calculations to perturbations around a Minkowski spacetime. The justification of such a choice is that the null geodesics of two conformal spacetimes are identical, so that lensing effects are the same. The only difference when considering an expanding Friedman–Lemaître universe will be the computation of the distances, i.e. of $`\lambda _\mathrm{S}`$ and $`\lambda _\mathrm{L}`$ in (37). Note also that the expansion of the universe affects the dynamics of the topological defects network (see e.g. ) but this will not be relevant on time scales of order of the impact parameter which are small compared to the dynamical scales of the universe. ## III Dynamics of cosmic strings The determination of the amplification matrix $`𝒜`$ requires the knowledge of stress–energy tensor of the cosmic strings. In this section, we first present the definition of this tensor and derive the equation of motion of a string. We then focus on the particular case of a non–rotating cosmic loop. ### A Equations of motion of the string As shown by Carter , there is an elegant way to describe the dynamics of a $`(d<n)`$–brane embedded in a $`n`$–dimensional spacetime. In this section, we recall the main steps of this formalism necessary to obtain the dynamical equation of evolution of the string; details can be found in . It requires the introduction of the induced metric on the string worldsheet $$f_{AB}g_{\mu \nu }x_{,A}^\mu x_{,B}^\nu ,$$ (51) where $`A,B\mathrm{}`$ refers to coordinates on the worldsheet and from which one can construct the fundamental tensor $`\overline{\eta }^{\mu \nu }`$ and the orthogonal projector $`^{\mu \nu }`$ as $$\overline{\eta }^{\mu \nu }f^{AB}x_{,A}^\mu x_{,B}^\nu ,\text{and}_\nu ^\mu g_\nu ^\mu \overline{\eta }_\nu ^\mu .$$ (52) The covariant derivative $``$ defines a tangentially projected differentiation operator $$\overline{}_\mu \overline{\eta }_\mu ^\nu _\nu .$$ (53) The second fundamental tensor is defined by $$K_{\mu \nu }^{}{}_{}{}^{\rho }\overline{\eta }_\nu ^\sigma \overline{}_\mu \overline{\eta }_\sigma ^\rho ,$$ (54) and the condition that the worldsheet is integrable, i.e. that all its elements mesh to form a well defined $`d`$–surface, is expressed by $$K_{[\mu \nu ]}^{}{}_{}{}^{\rho }=0,$$ (55) which is a geometric identity for the worldsheet. In the case of a string ($`d=2`$), the Lagrangian density $`\widehat{}`$ can be expressed as $$\widehat{}=\frac{1}{\sqrt{g}}\sqrt{f}\mathrm{d}^2\zeta \overline{}\delta ^{(4)}\left[x^\mu x^\mu (\zeta ^A)\right]$$ (56) where $`f`$ is the determinant of the metric $`f_{AB}`$ defined in (51). $`\widehat{}`$ is distributional and not confined to the string worldsheet whereas $`\overline{}`$ is locally regular but confined on the string worldsheet. The string action can then be written either in terms of $`\widehat{}`$ or of $`\overline{}`$ as $$S_{\mathrm{string}}=\mathrm{d}^2\zeta \sqrt{f}\overline{}=\mathrm{d}^4x\sqrt{g}\widehat{}.$$ (57) Now, making an infinitesimal variation $`\delta g_{\mu \nu }`$ of the action (57) provides a definition of the “surface” stress energy tensor density $`\stackrel{~}{T}^{\mu \nu }`$ (confined and regular) and of the stress energy tensor density $`\widehat{T}^{\mu \nu }`$ (distributional and unconfined) by $$2\delta S_{\mathrm{string}}=\mathrm{d}^2\zeta \sqrt{f}\stackrel{~}{T}^{\mu \nu }\delta g_{\mu \nu }=\mathrm{d}^4x\sqrt{g}\widehat{T}^{\mu \nu }\delta g_{\mu \nu }.$$ These two stress–energy tensor are related by $$\widehat{T}^{\mu \nu }=\frac{1}{\sqrt{g}}\mathrm{d}^2\zeta \sqrt{f}\stackrel{~}{T}^{\mu \nu }\delta ^{(4)}\left[x^\mu x^\mu (\zeta ^A)\right].$$ (58) The internal coordinate stress energy tensor, $`\stackrel{~}{T}^{AB}`$, has been projected onto a corresponding background stress energy tensor, $`\stackrel{~}{T}^{\mu \nu }`$, as in (52). One can then show that the general form of the dynamical equation of the string is $$\overline{}_\mu \stackrel{~}{T}^{\mu \nu }=f^\nu $$ (59) $`f^\nu `$ being the force exerted on the string by any background field such as e.g. an electromagnetic field. This dynamical equation (59) is equivalent to the more natural equation of conservation $`_\mu \widehat{T}^{\mu \nu }=\widehat{f}^\nu `$ with $`\widehat{f}^\nu `$ related to $`f^\nu `$ as in (58). For a string of energy per unit length $`U`$ and of tension $`T`$, the surface stress energy tensor density is of the form $$\stackrel{~}{T}^{\mu \nu }=Uu^\mu u^\nu Tv^\mu v^\nu ,$$ (60) where $`u^\mu `$ and $`v^\mu `$ are respectively a timelike ($`u^\mu u_\mu =1`$) and a spacelike ($`v^\mu v_\mu =+1`$) unit vector tangent to the string worldsheet (i.e. $`_\nu ^\mu u^\nu =_\nu ^\mu v^\nu =0`$) so that $`\overline{\eta }^{\mu \nu }=u^\mu u^\nu +v^\mu v^\nu `$. The dynamical equation governing the evolution of the string is given by the tangential projection of (59) which, in the free case we are considering, leads to $$\eta _\nu ^\rho \overline{}_\mu \stackrel{~}{T}^{\mu \nu }=0.$$ (61) This equation of evolution (61) can then be solved once we are given an equation of state, i.e. $`U(T)`$. Such equations are provided once the microscopic structure of the string is described. The most well known are the Goto–Nambu strings ($`U=T`$) and the Nielsen–Olesen strings ($`U+T=const.`$) and some have been obtained for superconducting strings . ### B Application to a non rotating cosmic string loop In the case of a non rotating circular loop of radius $`R`$, we work in cylindrical coordinates $`(t,r,\theta ,z)`$ and assume that it is lying in the plane $`z=z_s`$. Parametrising the loop worldsheet as $$t_{\mathrm{loop}}=t,\stackrel{}{r}_{\mathrm{loop}}\stackrel{}{r}_{}=R(t)(\mathrm{cos}\theta ,\mathrm{sin}\theta ),x_3z_{\mathrm{loop}}z_s=0,$$ (62) one can show that the unit spacelike vector tangent to the string worldsheet is $`v^\mu =\theta ^\mu `$ and we have $$u_\mu =\gamma \delta _\mu ^t+\gamma \dot{R}\delta _\mu ^r,\theta _\mu =R\delta _\mu ^\theta ,$$ (63) with $`\gamma (1\dot{R}^2)^{1/2}`$ being the Lorenz factor associated with the radial contraction/expansion of the string. They satisfy $`_\nu ^\mu u^\nu =_\nu ^\mu \theta ^\nu =0`$. From (60) and (58), the stress–energy tensor entering the Einstein equations (19) is given by $$T^{\mu \nu }(\stackrel{}{x},t)=\left(\begin{array}{cccc}\gamma ^2U& \gamma ^2\dot{R}U& 0& 0\\ \gamma ^2\dot{R}U& \gamma ^2\dot{R}^2U& 0& 0\\ 0& 0& R^2T& 0\\ 0& 0& 0& 0\end{array}\right)\delta (zz_s)\delta (rR).$$ (64) We now need to write the dynamical evolution equation (61) to get the evolution of the radius of the loop as a function of time. Using that $`\overline{\eta }_{\mu \nu }=u_\mu u_\nu +\theta _\mu \theta _\nu `$ and the expression (60) of the stress–energy tensor of the string, the equation (61) can be rewritten as $$(u_\mu u^\nu +\theta _\mu \theta ^\nu )_\nu \left(Uu^\mu u^\rho T\theta ^\mu \theta ^\rho \right)=0,$$ (65) which reduces to $`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}(\gamma UR)`$ $`=`$ $`0`$ (66) $`{\displaystyle \frac{\mathrm{d}^2}{\mathrm{d}t^2}}R`$ $`=`$ $`{\displaystyle \frac{1}{\gamma ^2R}}{\displaystyle \frac{T}{U}}.`$ (67) This system of equations for ($`U,T,R`$) is not closed and can be solved when one specifies an equation of state $`U(T)`$. Equation (66) shows that the total energy $`\gamma RU`$ of the loop is a constant of motion. Note that we have not used the identity (55) which is identically satisfied in our present example. ## IV Lensing by a cosmic string ### A A first example As a first application, let us consider a static straight cosmic string lying along the axis $`x_2`$ in a plane perpendicular to the line of sight (direction $`x_3`$ on figure 4 with $`\phi =0`$) so that $$_{\mu \nu }=\frac{1}{2}\left(\begin{array}{cccc}UT& & & \\ & U+T& & \\ & & UT& \\ & & & U+T\end{array}\right)\delta (\lambda _\mathrm{L}\theta _1)\delta (x_{}x_{}(\lambda _\mathrm{L})),$$ (68) with $`U`$ and $`T`$ depending on $`x_2=\lambda _\mathrm{L}\theta _2`$ only. Since $`=U(x_2)\delta (x_1)\delta (x_{}x_{}(\lambda _\mathrm{L}))`$, the first integral of (28), after integration over $`x_1^{}`$, reduces to $$𝒜_{ab}=I_{ab}+2\frac{G}{\lambda _\mathrm{S}}U(x_2^{})dx_2^{}_{ab}J((x_2x_2^{})^2+x_1^2)$$ with $$J((x_2x_2^{})^2+x_1^2)=_0^{\lambda _\mathrm{S}}\frac{\lambda (\lambda _\mathrm{S}\lambda )}{\sqrt{(x_2x_2^{})^2+x_1^2+(\lambda \lambda _\mathrm{L})^2}}d\lambda $$ where we have chosen a parametrisation such that $`x_{}(\lambda )=\lambda `$. This latter integral can be computed and gives $`J(A)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\lambda _\mathrm{S}3\lambda _\mathrm{L})\sqrt{A+(\lambda _\mathrm{S}\lambda _\mathrm{L})^2}\left(\lambda _\mathrm{S}{\displaystyle \frac{3}{2}}\lambda _\mathrm{L}\right)\sqrt{A+\lambda _\mathrm{L}^2}`$ (70) $`+(\lambda _\mathrm{S}\lambda _\mathrm{L}\lambda _\mathrm{L}^2{\displaystyle \frac{1}{2}}A)\mathrm{ln}{\displaystyle \frac{\lambda _\mathrm{S}\lambda _\mathrm{L}+\sqrt{A+(\lambda _\mathrm{S}\lambda _\mathrm{L})^2}}{\lambda _\mathrm{L}+\sqrt{A+\lambda _\mathrm{L}^2}}}`$ with $`A(x_2x_2^{})^2+x_1^2`$. It follows that the amplification matrix is given by $$𝒜_{ab}=I_{ab}+2\frac{G}{\lambda _\mathrm{S}}U(x_2^{})dx_2^{}\left(2J^{}(A)\delta _{ab}+(x_ax_a^{})(x_bx_b^{})J^{\prime \prime }(A)\right)$$ (71) where $`J^{}\mathrm{d}J/\mathrm{d}A`$ and $`x_1^{}=0`$. Now, one can estimate the dominant term in the integral of expression (71) when we are looking close to the string (i.e. when $`x_1,x_2\lambda _\mathrm{S}\lambda _\mathrm{L},\lambda _\mathrm{L}`$). For that purpose we assume that $`\lambda _\mathrm{S}\lambda _\mathrm{L}\lambda _\mathrm{L}`$ and set $`\lambda \lambda _\mathrm{S}\lambda _\mathrm{L}\lambda _\mathrm{L}`$. We then split the integral of expression (71) in a contribution where $`Ax_2^2<\lambda ^2`$ and another where $`Ax_2^2>\lambda ^2`$. Using the expansion of $`J(A)`$ in these two regimes as $$J(A)=\left(\lambda _\mathrm{S}\lambda _\mathrm{L}\right)\{\begin{array}{cc}\lambda _\mathrm{L}\mathrm{ln}(A)+constant+𝒪\left(\frac{A}{\lambda ^2}\right)\hfill & A/\lambda ^2<1\hfill \\ \sqrt{A}\left(1+𝒪\left(\frac{\lambda }{\sqrt{A}}\right)\right)\hfill & A/\lambda ^2>1\hfill \end{array}$$ one can show that, as long as we are looking close to the string, the dominant contribution to the amplification matrix is given by $$𝒜_{ab}I_{ab}4G\frac{\lambda _\mathrm{S}\lambda _\mathrm{L}}{\lambda _\mathrm{S}}_{\theta ^a}_{\theta ^b}\mathrm{ln}\left\|\stackrel{}{\theta }^I\stackrel{}{\theta }^{}\right\|U(\theta _2^{})d\theta _2^{}.$$ (72) In the particular case where $`U`$ is constant, this approximate give the general result of the deflection by a straight cosmic string and one can thus thing that this is a good approximation when $`U`$ is fluctuating around a mean value. Note also that for such an infinite string lying in a plane perpendicular to the line of sight we recover the general form (38) of the deformation matrix in the thin lens approximation. In the following of this article, we investigate more general results concerning the deformation fiels by a cosmic string which do not assume that we are in the thin lens regime. ### B General Results In the general case, the source term generated by a cosmic string will be localized on the string worldsheet so that \[see equation (60)\] $$(\stackrel{}{x},t)=d\zeta \stackrel{~}{}(\stackrel{}{x},t)\delta \left(\stackrel{}{x}\stackrel{}{r}(\zeta ,t)\right)$$ (73) where $`(t,\stackrel{}{r}(\zeta ,t))`$ is a parameterization of the string worldsheet; $`\stackrel{}{r}(\zeta ,t)`$ represents the locus of the string on each constant time hypersurface and $`\stackrel{~}{}`$ is the energy density per unit length \[note that we have chosen a parametrisation such that $`t`$ is both the coordinate time and an intrinsic coordinate of the string worldsheet which implies that we have a one dimensional integration on the spatial internal coordinate $`\zeta `$ and not a two dimensional integration as in (60)\]. Inserting this decomposition in (21), we deduce that the deflecting potential (27) is given, after integration over space, by $$\mathrm{\Phi }(\stackrel{}{x},t)=2G\frac{\mathrm{d}t^{}\mathrm{d}\zeta }{\left\|\stackrel{}{x}\stackrel{}{r}(\zeta ,t^{})\right\|}\stackrel{~}{}[\stackrel{}{r}(\zeta ,t^{}),t^{}]\delta \left(t^{}t+\left\|\stackrel{}{x}\stackrel{}{r}(\zeta ,t^{})\right\|\right).$$ (74) Following , the integration over $`t^{}`$ can be performed by introducing $`t_{\mathrm{string}}(\stackrel{}{x},\zeta ,t)`$ solution of $$tt_{\mathrm{string}}(\stackrel{}{x},\zeta ,t)=\left\|\stackrel{}{x}\stackrel{}{r}(\zeta ,t_{\mathrm{string}}(\stackrel{}{x},\zeta ,t))\right\|$$ (75) so that $$\delta \left(t^{}t+\left\|\stackrel{}{x}\stackrel{}{r}(\zeta ,t^{})\right\|\right)=\frac{\left\|\stackrel{}{x}\stackrel{}{r}(\zeta ,t^{})\right\|}{\left\|\stackrel{}{x}\stackrel{}{r}(\zeta ,t^{})\right\|_t\stackrel{}{r}(\zeta ,t^{}).(\stackrel{}{x}\stackrel{}{r}(\zeta ,t^{}))}\delta \left(t^{}t_{\mathrm{string}}\right)$$ from which we deduce that $$\mathrm{\Phi }(\stackrel{}{x},t)=2G\frac{\stackrel{~}{}[\stackrel{}{r}(\zeta ,t_{\mathrm{string}}(\stackrel{}{x},\zeta ,t)),t_{\mathrm{string}}(\stackrel{}{x},\zeta ,t)]}{\left\|\stackrel{}{x}\stackrel{}{r}(\zeta ,t_{\mathrm{string}}(\stackrel{}{x},\zeta ,t))\right\|_t\stackrel{}{r}(\zeta ,t_{\mathrm{string}}(\stackrel{}{x},\zeta ,t)).(\stackrel{}{x}\stackrel{}{r}(\zeta ,t_{\mathrm{string}}(\stackrel{}{x},\zeta ,t)))}d\zeta .$$ (76) $`\mathrm{\Phi }`$ on the point $`(\stackrel{}{x},t)`$ is then given by the projection of the string energy on the past light cone of this point, i.e. on the curve $`\{\stackrel{}{r}(\zeta ,t_{\mathrm{string}}(\stackrel{}{x},\zeta ,t)),t_{\mathrm{string}}(\stackrel{}{x},\zeta ,t)\}`$ which is the intersection of the string worldsheet with the past light cone of the event $`(\stackrel{}{x},t)`$ \[see figure 2\]. Now, focusing on $`\kappa `$, the deformation matrix is explicitly given by (27) and (28) and, proceeding as for the deflection angle, one can easily sort out that $`\delta ^{ab}𝒟_{ab}^{(1)}`$ $`=`$ $`{\displaystyle _0^{\lambda _\mathrm{S}}}\lambda (\lambda _\mathrm{S}\lambda )\mathrm{\Delta }_{}\mathrm{\Phi }(\stackrel{}{x}_{},x_{}(\lambda ),t(\lambda ))d\lambda +{\displaystyle _0^{\lambda _\mathrm{S}}}\lambda (\lambda _\mathrm{S}\lambda )\delta ^{ab}\mathrm{\Psi }_{ab}d\lambda `$ (77) $`=`$ $`{\displaystyle _0^{\lambda _\mathrm{S}}}\lambda (\lambda _\mathrm{S}\lambda )\left[\left(_t^2_{}^2\right)\mathrm{\Phi }(\stackrel{}{x}_{},x_{}(\lambda ),t(\lambda ))8\pi G(\stackrel{}{x}_{},x_{}(\lambda ),t(\lambda ))\right]d\lambda `$ (78) $`+`$ $`{\displaystyle _0^{\lambda _\mathrm{S}}}\lambda (\lambda _\mathrm{S}\lambda )\delta ^{ab}\mathrm{\Psi }_{ab}d\lambda `$ (79) The first term of the first integral vanishes when evaluated on the photon geodesic so that the contribution of the first integral to the convergence $`\kappa `$ defined in (15) reduces to $`\kappa (\stackrel{}{x}_{},t_0)=4\pi G{\displaystyle _0^{\lambda _\mathrm{S}}}{\displaystyle \frac{\lambda (\lambda _\mathrm{S}\lambda )}{\lambda _\mathrm{S}}}(\stackrel{}{x}_{},x_{}(\lambda ),t(\lambda ))d\lambda =4\pi G{\displaystyle _{t_{\mathrm{emission}}}^{t_0}}{\displaystyle \frac{(tt_{\mathrm{emission}})(t_0t)}{t_0t_{\mathrm{emission}}}}(\stackrel{}{x}_{},x_{}(t),t)dt.`$ (80) where we recall that $`t_0t(\lambda =0)`$ is the time of reception and where $`t_{\mathrm{emission}}t(\lambda _\mathrm{S})`$ is the time of emission. The contribution of the second integral of (79) vanishes since (i) $`\mathrm{\Psi }_{ab}\delta ^{ab}=0`$ both for scalar perturbations \[$`\varphi `$ and $`\psi `$ in (39)\] and for vector perturbations \[$`A_i`$ in (39)\] and (ii) for tensor modes \[$`\overline{E}_{ij}`$ in (39)\] $`\mathrm{\Psi }_{ab}\delta ^{ab}\frac{\mathrm{d}}{\mathrm{d}\lambda }\left[(_t+_{})(\overline{E}_1^1+\overline{E}_2^2)\right]=0`$ when evaluated on the photon trajectory. We conclude that the convergence $`\kappa `$ is given by (80). It is then given by the distribution of matter evaluated on the photon trajectory, up to a geometrical factor. The lensing by a cosmic string is thus equivalent to the lensing by a linear distribution of matter. As a consequence, $`\kappa =0`$ everywhere but on directions such that the observer past light cone intersects the string worldsheet; this result, valid whatever the equation of state, is one of the main results of this article. It holds for any relativistic lens and does not rely on the thin lens approximation. For instance if the string is lying in a plane perpendicular to the line of sight then (80) reduces to $`\kappa (\stackrel{}{x}_{},t_0)=4\pi G{\displaystyle \frac{\lambda _\mathrm{S}\lambda _\mathrm{L}}{\lambda _\mathrm{S}}}\lambda _\mathrm{L}\mathrm{\Sigma }(\stackrel{}{x}_{},x_{}(\lambda _\mathrm{L}),t(\lambda _\mathrm{L}))_{}.\stackrel{}{\alpha }.`$ (81) Under this form, again we see that $`\kappa =0`$ everywhere but on directions intersecting the string worldsheet. Note that the drivation of (80) relies strongly on the fact that we took the trace of the deformation matrix in (79). Therefore, similar expression to (80) cannot be obtained for the other components of the amplification matrix and one has to rely on (76). It has also to be noted that the expression of the amplification matrix in (50) (which relies on the thin lens approximation) together with the general result (80) implies that the phenomenological description of $`𝒜_b^a`$ in (38) is very general. It holds for any string dynamics provided the extension of the string is small enough for the thin lens approximation to hold. Different aspects of these results are illustrated in the next paragraphs. A more elaborate phenomenological investigation based on Eq. (38) is proposed in a companion paper . ### C Lensing by a non–rotating cosmic string loop perpendicular to the line of sight We now consider as an application the case of a non–rotating circular loop oscillating in a plane perpendicular to the line of sight. Its dynamics is described by the set of equations (6667) and from its parameterization (62) we deduce that $$(\stackrel{}{x},t)=_0^{2\pi }R(t)\gamma (t)U(t)\delta (x_3)\delta (\stackrel{}{x}_{}\stackrel{}{r}_{})d\theta .$$ (82) The deflecting potential (27) integrated along the line of sight is then given by $`I[\mathrm{\Phi }]`$ $``$ $`{\displaystyle _0^{\lambda _\mathrm{S}}}\mathrm{\Phi }(\stackrel{}{x}(\lambda ),t(\lambda ))d\lambda `$ (83) $`=`$ $`2G{\displaystyle _{x_\mathrm{L}}^{x_\mathrm{S}x_\mathrm{L}}}dx_3{\displaystyle _0^{2\pi }}d\theta {\displaystyle _{R^2}}\mathrm{d}^2\stackrel{}{x}_{}^{}{\displaystyle _{\mathrm{}}^+\mathrm{}}dt^{}\gamma RU{\displaystyle \frac{\delta \left(t^{}+x_3\tau _0+\sqrt{\left(\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right)^2+x_3^2}\right)}{\sqrt{\left(\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right)^2+x_3^2}}}\delta \left(\stackrel{}{x}_{}^{}\stackrel{}{r}_{}(t^{},\theta )\right)`$ (84) where we have parameterized the geodesic as $`t=\tau _0x_3`$ with $`\tau _0t_0x_\mathrm{L}`$. Now, defining the new variable $`z`$ as $$zx_3+\sqrt{\left(\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right)^2+x_3^2}$$ (85) which satisfies $$\frac{\mathrm{d}z}{z}=\frac{\mathrm{d}x_3}{\sqrt{\left(\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right)^2+x_3^2}}.$$ (86) The integrated potential (83) reduces to $$I[\mathrm{\Phi }]=2G_0^{2\pi }d\theta _{R^2}\mathrm{d}^2\stackrel{}{x}_{}^{}_A^B\frac{\mathrm{d}z}{z}_{\mathrm{}}^+\mathrm{}𝑑t^{}\gamma RU\delta (t^{}+z\tau _0)\delta \left(\stackrel{}{x}_{}^{}\stackrel{}{r}_{}(t^{},\theta )\right)$$ (87) where the limits of integration are $`A`$ $`=`$ $`x_\mathrm{L}+\sqrt{\left(\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right)^2+x_\mathrm{L}^2}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\left(\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right)^2}{x_\mathrm{L}}}`$ (88) $`B`$ $`=`$ $`x_\mathrm{S}x_\mathrm{L}+\sqrt{\left(\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right)^2+(x_\mathrm{S}x_\mathrm{L})^2}2(x_\mathrm{S}x_\mathrm{L}).`$ (89) The approximate values of $`A`$ and $`B`$ are obtained at lowest order when we consider zones close to the string in comparison with the distance string–observer and string–source. After integration over $`t^{}`$ and using the loop equation of motion (6667) which state that $`\gamma RU`$ is constant, we get that $$I[\mathrm{\Phi }]=2G\gamma RUJ(\stackrel{}{x}_{},t_0),$$ (90) where $`J(\stackrel{}{x}_{},t_0)`$ is a dimensionless geometrical integral given by $$J(\stackrel{}{x}_{},t_0)=_0^{2\pi }d\theta _{R^2}\mathrm{d}^2\stackrel{}{x}_{}^{}_{\frac{1}{2}\frac{\left(\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right)^2}{x_\mathrm{L}}}^{2(x_\mathrm{S}x_\mathrm{L})}\frac{\mathrm{d}z}{z}\delta \left(\stackrel{}{x}_{}^{}\stackrel{}{r}_{}(z\tau _0,\theta )\right).$$ (91) This integral can be rewritten as $$J(\stackrel{}{x}_{},t_0)=_0^{2(x_\mathrm{S}x_\mathrm{L})}\frac{\mathrm{d}z}{z}_0^{2\pi }d\theta _{\left\|\stackrel{}{x}_{}\stackrel{}{x}_{}^{}\right\|^2<2x_\mathrm{L}z}\mathrm{d}^2\stackrel{}{x}_{}^{}\delta \left(\stackrel{}{x}_{}^{}\stackrel{}{r}_{}(z\tau _0,\theta )\right).$$ (92) We deduce that, at $`z`$ constant, the integral over $`x_{}^{}`$ reduces to the computation of the angle $`\beta `$ (see figure 3) of string within the disk of radius $`\sqrt{2x_\mathrm{L}z}`$ which can be computed as followed: Setting $`u\sqrt{2x_\mathrm{L}z}`$ and $`hOH`$, we deduce from $$u=u_1+u_2,u_1^2+h^2=\stackrel{}{x}_{}^2,u_2^2+h^2=R^2$$ (93) where $`u_1CH`$ and $`u_2HD`$ on figure 3 and $`R`$ now stands for $`R(t=\tau _0)`$, that $$u_1=\frac{\stackrel{}{x}_{}^2R^2+u^2}{2u}u_2=\frac{\stackrel{}{x}_{}^2+R^2+u^2}{2u}$$ (94) and thus that $$\mathrm{cos}\beta =\frac{\stackrel{}{x}_{}^2+R^2u^2}{2\|\stackrel{}{x}_{}\|R}.$$ (95) The integral over $`\theta `$ and $`\stackrel{}{x}_{}^{}`$ obviously vanishes if $`u<\left\|\|\stackrel{}{x}_{}\|R\right\|`$, reduces to $`2\pi `$ when $`u>\|\stackrel{}{x}_{}\|+R`$ and gives $`2\beta (u)`$ otherwise. Then, after splitting the integral over $`u`$ in (92) in three pieces ($`[0,\left\|\|\stackrel{}{x}_{}\|R\right\|]`$, $`[\left\|\|\stackrel{}{x}_{}\|R\right\|,\|\stackrel{}{x}_{}\|+R]`$ and $`[\|\stackrel{}{x}_{}\|+R,2\sqrt{x_\mathrm{L}(x_\mathrm{S}x_\mathrm{L})}]`$) we get $`J(\stackrel{}{x}_{},t_0)`$ $`=`$ $`{\displaystyle _{\left\|\|\stackrel{}{x}_{}\|R\right\|}^{\|\stackrel{}{x}_{}\|+R}}2\beta (u){\displaystyle \frac{\mathrm{d}u}{u}}+2\pi {\displaystyle _{\|\stackrel{}{x}_{}\|+R}^{2\sqrt{x_\mathrm{L}(x_\mathrm{S}x_\mathrm{L})}}}{\displaystyle \frac{\mathrm{d}u}{u}},`$ (96) with $`\beta `$ given by (95). We can compute this integral in the two following regimes 1. If $`\|\stackrel{}{x}_{}\|<R(\tau _0)`$, (96) can be rewritten as $$J=2\|\stackrel{}{x}_{}\|R_1^1\frac{\text{Arccos}v}{\|\stackrel{}{x}_{}\|^2+R^22\|\stackrel{}{x}_{}\|Rv}dv+2\pi \mathrm{ln}\left[\frac{2\sqrt{x_\mathrm{L}(x_\mathrm{S}x_\mathrm{L})}}{\|\stackrel{}{x}_{}\|+R}\right]$$ (97) which can be computed to give $$J=C_1$$ (98) where $`C_1`$ is a constant depending on $`x_\mathrm{L}`$, $`x_\mathrm{S}`$ and $`R(\tau _0)`$. Then, there is no deflection of a light ray passing inside a large loop, as first pointed out in and as expected from the Gauss theorem. 2. If $`\|\stackrel{}{x}_{}\|>R(\tau _0)`$, (97) now gives after integration $$J=C_12\pi \mathrm{ln}\frac{\|\stackrel{}{x}_{}\|}{R}$$ (99) and we conclude that a small loop perpendicular to the line of sight acts as a point mass $`M=2\pi \gamma RU`$ whatever its equation of state. We checked that this is also valid for a tilted circular loop and it is natural to expect that the fact that a loop acts as a point mass at a distance larger than its caracteristic size is valid whatever the geometry of the loop. One can also check from (9899) that $`\mathrm{\Delta }_{}J(\stackrel{}{x}_{},t_0)=0`$ if $`\stackrel{}{x}_{}R(\tau _0)`$. Since $`\kappa \mathrm{\Delta }_{}J`$, we recover the result from (80) that the convergence vanishes if $`\stackrel{}{x}_{}\stackrel{}{r}_{}`$. ### D Tilted static straight cosmic string To finish, let us consider a tilted static straight cosmic string aligned along an axis making an angle $`\phi `$ with the direction $`x_2`$ \[see figure 4\] and for which, from (68), $$=\left[U(\mathrm{})T(\mathrm{})\mathrm{sin}^2\phi \right]d\mathrm{}\delta (\stackrel{}{x}\stackrel{}{r}(\mathrm{}))$$ (100) where $`\stackrel{}{r}(\mathrm{})`$ is a parameterization of the string. If we choose $`\mathrm{}`$ such that $$\stackrel{}{r}(\mathrm{}):r_1=0,r_2=\mathrm{}\mathrm{cos}\phi ,r_3=x_{}(\lambda _\mathrm{L})\mathrm{}\mathrm{sin}\phi ,$$ (101) the deflecting potential (76) is given by $$\mathrm{\Phi }(\stackrel{}{x})=2G\frac{\mathrm{d}\mathrm{}}{\left\|\stackrel{}{x}\stackrel{}{r}(\mathrm{})\right\|}\left[UT\mathrm{sin}^2\phi \right](\mathrm{}).$$ (102) When $`U`$ and $`T`$ are constant, it can be integrated to get $$\mathrm{\Phi }(\stackrel{}{x})=2G\left[UT\mathrm{sin}^2\phi \right]\left(C_{\mathrm{}}\mathrm{ln}\left[\stackrel{}{x}^2(x_2\mathrm{cos}\phi x_3\mathrm{sin}\phi )^2\right]\right)$$ (103) where $`C_{\mathrm{}}`$ is an infinite constant and where we have introduced $`x_3x_{}x_\mathrm{L}`$. The infinite constant $`C_{\mathrm{}}`$ can then be forgotten because only the derivatives of $`\mathrm{\Phi }`$ enters the computation of the deflection angle and of the amplification matrix which are the observable quantities. The deflection angle is then given by $$\stackrel{}{\alpha }=_{}_{x_\mathrm{L}}^{x_\mathrm{S}x_\mathrm{L}}\mathrm{\Phi }(x_1,x_2,x_3)dx_3.$$ (104) After integration over $`x_3`$, we get $`\alpha _1`$ $`=`$ $`{\displaystyle \frac{4G\left[UT\mathrm{sin}^2\phi \right]}{\mathrm{cos}\phi }}\left[\mathrm{arctan}\left({\displaystyle \frac{x_2\mathrm{sin}\phi +x_3\mathrm{cos}\phi }{x_1}}\right)\right]_{x_3=x_\mathrm{L}}^{x_3=x_\mathrm{S}x_\mathrm{L}}`$ (105) $`\alpha _2`$ $`=`$ $`4G\left[UT\mathrm{sin}^2\phi \right]\mathrm{tan}\phi \left[\mathrm{ln}\sqrt{x_1^2+(x_2\mathrm{sin}\phi +x_3\mathrm{cos}\phi )^2}\right]_{x_3=x_\mathrm{L}}^{x_3=x_\mathrm{S}x_\mathrm{L}}.`$ (106) In the limit where $`(x_\mathrm{S}x_\mathrm{L})\left\|\mathrm{cos}\phi \right\|`$ and $`x_\mathrm{L}\left\|\mathrm{cos}\phi \right\|`$ are much larger than $`\|x_1\|`$ and $`\left\|x_2\mathrm{sin}\phi \right\|`$, we get $`\alpha _1`$ $``$ $`4\pi G\left[U\mathrm{cos}\phi +(UT)\mathrm{sin}\phi \mathrm{tan}\phi \right]`$ (107) $`\alpha _2`$ $``$ $`4GU\mathrm{tan}\phi \mathrm{ln}{\displaystyle \frac{x_\mathrm{S}x_\mathrm{L}}{x_\mathrm{L}}}.`$ (108) This has to be compared with the standard result for a Goto–Nambu string for which $`\alpha =4\pi GU\mathrm{cos}\phi `$ . Now, as pointed out by Peter in the case of a string perpendicular to the line of sight, there are two origins to the deflection: the deficit angle (term proportional to $`U+T`$) and a contribution from the curvature (proportional to $`UT`$). One can understand such a result by decomposing the stress–energy tensor (68) as $`2_{\mu \nu }=2U\times \mathrm{diag}(0,1,0,1)+(UT)\times \mathrm{diag}(1,1,1,1)`$, i.e. as the superposition of a Goto–Nambu string and a linear distribution of non–relativistic matter of density $`\rho UT`$ per unit length. Then it is straightforward to see that the bending angle of the second contribution depends only on the projected mass per unit length and thus becomes larger by a factor $`1/\mathrm{cos}\phi `$ as found in (107). A consequence of this result is that, for general cosmic strings not perpendicular to the line of sight, one expects to have larger deflection than for a Goto–Nambu string. In this case, the accuracy of the thin lens approximation can be investigated. For that purpose, we remind that the shear is given by $$\left(\begin{array}{c}\gamma _1\\ \gamma _2\end{array}\right)=2G_{x_\mathrm{L}}^{x_\mathrm{S}x_\mathrm{L}}\frac{(x_\mathrm{L}+x_3)(x_\mathrm{S}x_\mathrm{L}x_3)}{x_\mathrm{S}}\left(\begin{array}{c}\frac{1}{2}\left[_2^2_1^2\right]\\ _1_2\end{array}\right)_{\mathrm{}}^{\mathrm{}}\frac{\left[U(\mathrm{})T(\mathrm{})\mathrm{sin}^2\phi \right]\mathrm{d}\mathrm{}}{\sqrt{x_1^2+(x_2\mathrm{}\mathrm{cos}\phi )^2+(x_3\mathrm{}\mathrm{sin}\phi )^2}}.$$ (109) Due to the derivatives with respect to $`x_1`$ and $`x_2`$, it is easy to see that the integral over $`\mathrm{}`$ is peaked around $`x_3\mathrm{}\mathrm{sin}\phi x_2\mathrm{tan}\phi `$. Thus, on a field of width $`x_2=x_\mathrm{S}\theta _2`$, the variation of geometric factor is bounded by $$\frac{\delta \left\|(x_\mathrm{L}x_3)(x_\mathrm{S}x_\mathrm{L}x_3)\right\|}{x_\mathrm{L}(x_\mathrm{S}x_\mathrm{L})}\frac{\left(1+2\frac{x_\mathrm{L}}{x_\mathrm{S}}\right)}{x_\mathrm{L}/x_\mathrm{S}\left(1x_\mathrm{L}/x_\mathrm{S}\right)}\left\|\mathrm{tan}\phi \right\|\theta _2<\frac{3}{x_\mathrm{L}/x_\mathrm{S}\left(1x_\mathrm{L}/x_\mathrm{S}\right)}\|\mathrm{tan}\phi \|\theta _2,$$ from which we deduce that since $`\theta _21`$, the thin lens approximation is still very good for tilted string with a tilt not larger than $`\phi =\pi /4`$ say (see figure 5 for a numerical estimation of the relative error). ### E Discussion In this section, we have shown that the deflecting potential of any extended lens with relativistic motion is obtained by considering the projected energy density on the photon past light cone. This implies, in the case of cosmic strings, that the convergence $`\kappa `$ vanishes everywhere but on the string projection. We then studied the case of a loop oscillating in a plane perpendicular to the line of sight and show that the equation of motion of the loop can be used to integrate the deflecting potential. We found that a photon propagating inside such a circular loop was not deflected and those propagating outside were deflected as if the loop was a point mass object. This generalizes the result by de Laix and Vachaspati to strings with any equation of state and shows how the equation of motion of the loop enables to factorize the integrated deflecting potential. This lets us conjecture that this result will be valid whatever the geometry of the loop, the geometric factor $`J`$ being different for each loop geometry. We finished by discussing the case of static tilted cosmic strings to emphasize that, for non–Goto–Nambu strings, there can be larger deflections and we also discussed on this example the validity of the thin lens approximation for strings. ## V Phenomenology of a deformation field with $`\kappa =0`$ As seen in the previous section, we expect the deformation field of a cosmic string to be such that $`\kappa =0`$. Indeed, in (80) we only showed that $`\kappa `$ vanishes in all directions such that the light ray arriving with this direction does not intersect the string worldsheet. In this section, we are mainly interested in objects which do not overlap the string such that their deformation is the one with a $`\kappa =0`$ field. Another restriction of this study is that we assume that the caracteristic size of the object is smaller than the caracteristic size of the variaiton of the shear $`\gamma `$; thus we donnot consider galaxies lying in the immediate neighborhood of the string (see the companion article for an illustration). With these restrictions, we can consider that the deformation field has a zero convergence and the goal of this section is to study the main phenomenological properties of such a field. The two kinds of questions we would like to answer are: * Given a source of surface $`S^\mathrm{S}`$ and ellipticity $`e^\mathrm{S}`$, what can we say about the morphologies of all its possible images? * Given two objects, how can we know that they are the images of the same source? This question reduces to study the allowed morphologies of the sources that have the same images. We start by setting the general framework and then study respectively (Q1) and (Q2). This study is a first step toward the discrimination between pairs of lensed sources by a cosmic string and fake lenses . This study is made within the assumption that the shear variations over observed background images is small. It may actually be a severe limitation for such an approach if the string energy density is small. ### A Describing the morphology of a cosmic object To any object of elliptic shape such as a galaxy or a cluster, we can associate a positive definite symmetric matrix $`M_{ab}`$ describing its shape as $$X^tMX1$$ (110) where $`X^t=(x_1,x_2)`$. This matrix can be diagonalised as $$M=P^t(\theta )\left(\begin{array}{cc}\lambda _+& 0\\ 0& \lambda _{}\end{array}\right)P(\theta )$$ (111) where $`P`$ is a rotation matrix defined by $$P(\theta )\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)$$ (112) and a subscript $`t`$ denotes the transposition. $`\lambda _{}\lambda _+`$ are the two positive eigenvalues of $`M`$ and $`\theta `$ is an angle describing the orientation of its principal axis with respect to the basis $`n_a^\mu `$. Thus any object can be characterized by the set ($`\theta `$, $`\lambda _{}`$, $`\lambda _+`$) from which we can define the surface $`S`$ and ellipticity $`e`$ of the object respectively as $`S(M)`$ $``$ $`\text{det}(M)=\lambda _+\lambda _{},`$ (113) $`e(M)`$ $``$ $`{\displaystyle \frac{\lambda _+\lambda _{}}{\lambda _++\lambda _{}}}=\sqrt{14{\displaystyle \frac{\text{det}(M)}{[\text{tr}(M)]^2}}}`$ (114) and we also define $`ϵ`$ as $$ϵ1e^2=4\frac{\text{det}(M)}{[\text{tr}(M)]^2}.$$ (115) These definitions can indeed be inverted to get the two eigenvalues in terms of $`e`$ and $`S`$ as $$\lambda _\pm ^2=S\left(\frac{1+e}{1e}\right)^{\pm 1}.$$ (116) Following Mellier , the ellipticity must be bounded by $$ϵ0.5e0.71.$$ (117) In the following, we will not be interested in the orientation of the object and we then define the shape as being the set $`(S,e)`$. The shape matrix $`M^\mathrm{I}`$ of any image can be related to the shape matrix $`M^\mathrm{S}`$ of its associated source as (see e.g. ) $$M^\mathrm{I}=𝒜^1M^\mathrm{S}𝒜^1$$ (118) \[this is valid only if we consider that the carateristic size of the source is smaller than the characteristic size associated with the variation of the shear\]. Decomposing $`M^\mathrm{S}`$ as in (111) and introducing $`\stackrel{~}{M}^\mathrm{I}PM^\mathrm{I}P^t`$ which represents the same shape as $`M^\mathrm{I}`$ after a rotation of $`\theta `$, we obtain that $$\stackrel{~}{M}^\mathrm{I}=(P𝒜^1P^t)\left(\begin{array}{cc}\lambda _+& 0\\ 0& \lambda _{}\end{array}\right)(P𝒜^1P^t).$$ (119) Thus, $`\stackrel{~}{M}^\mathrm{I}`$ is the image of the source $`\stackrel{~}{M}^\mathrm{S}\left(\begin{array}{cc}\lambda _+& 0\\ 0& \lambda _{}\end{array}\right)`$ by the transformation, $$\stackrel{~}{𝒜}^1P𝒜^1P^t=\frac{1}{1\stackrel{~}{\gamma }^2}\left(\begin{array}{cc}1+\stackrel{~}{\gamma }_1& \stackrel{~}{\gamma }_2\\ \stackrel{~}{\gamma }_2& 1\stackrel{~}{\gamma }_1\end{array}\right)\text{with}\stackrel{}{\stackrel{~}{\gamma }}=P(2\theta )\stackrel{}{\gamma }.$$ (120) As long as we are interested only on the shape (i.e. surface and ellipticity) of the sources and/or images, we always can choose one of them to be diagonal. ### B Morphology of the images of a given source ¿From the previous analysis, we can conclude that, if we are not interested in the relative orientation of the source and of the image, we can just restrict the problem by considering the source and the transformation matrix to be given by $$M^\mathrm{S}=\left(\begin{array}{cc}\lambda _+& 0\\ 0& \lambda _{}\end{array}\right),𝒜^1=\frac{1}{1\gamma ^2}\left(\begin{array}{cc}1+\gamma _1& \gamma _2\\ \gamma _2& 1\gamma _1\end{array}\right).$$ (121) Setting $$\stackrel{}{\gamma }\gamma \left(\begin{array}{c}\mathrm{cos}\alpha \\ \mathrm{sin}\alpha \end{array}\right)\text{with}\gamma 0\text{and}\alpha [0,2\pi ],$$ (122) we can easily show that the shape $`(S^\mathrm{I},e^\mathrm{I})`$ of the image is related to the one of the source $`(S^\mathrm{S},e^\mathrm{S}`$) by $`S^\mathrm{I}`$ $`=`$ $`{\displaystyle \frac{1}{(1\gamma ^2)^2}}S^\mathrm{S}`$ (123) $`ϵ^\mathrm{I}`$ $`=`$ $`{\displaystyle \frac{(1\gamma ^2)^2}{(1+\gamma ^2+2\gamma e^\mathrm{S}\mathrm{cos}\alpha )^2}}ϵ^\mathrm{S},`$ (124) where $`ϵ`$ is defined in (115). In the case of a circular source ($`e^\mathrm{S}=0`$) we deduce that, since $`(S^\mathrm{I}/S^\mathrm{S},e^\mathrm{I})`$ depends only on $`\gamma `$, (i) two images of same surface have same ellipticity and that (ii) all the images lie on a curve in the plane $`(S^\mathrm{I}/S^\mathrm{S},e^\mathrm{I})`$. In figure 6, we depict the variation of $`S^\mathrm{I}`$ and $`e^\mathrm{I}`$ in function of $`\gamma `$ and the ensemble of the images of a circular source. In the general case ($`e^\mathrm{S}0`$) we can determine all the morphologies of the images of a given source in the plane $`(S^\mathrm{I}/S^\mathrm{S},e^\mathrm{I})`$. In figures 7, we depict these sets respectively for $`e^\mathrm{S}=0.25`$ and $`e^\mathrm{S}=0.5`$. We note that all the curves have the same asymptot $`(S^\mathrm{I}/S^\mathrm{S})=\mathrm{},e^\mathrm{I}=1`$ which is reached when $`\gamma =1`$, i.e. on the critical line; on these points, the amplification $`\mu [\mathrm{det}(𝒜)]^1`$ is infinite. Whatever $`e^\mathrm{S}`$, there exist always two circular images (i.e. such that $`ϵ^\mathrm{I}=1`$) obtained for $$\gamma _\pm =\frac{e^\mathrm{S}}{1\sqrt{ϵ^\mathrm{S}}}\alpha \pi [2\pi ]$$ (125) (the condition on $`\alpha `$ is obtained from the requirement that $`\gamma >0`$; there are also two solutions for $`\alpha 0[2\pi ]`$ but they lead to negative values of $`\gamma `$). Now, since from (123) $`S_{}^\mathrm{I}/S_+^\mathrm{I}=(1\gamma _{}^2)^2/(1\gamma _+^2)^2`$, the measure of the area of two such circular images enables us (i) to determine the ellipticity $`e^\mathrm{S}`$ of their common source and (ii) the shears $`\gamma _\pm `$ of the two deformations. Indeed the bound on the ellipticity (117) has to be fulfilled by $`e^\mathrm{S}`$ and this can be a test to reject fake pairs of images. In the more general case of a pair of non circular images, one cannot reconstruct the ellipticity of their source but we can still conclude from the ratio of their surfaces if they correspond to transformations, $`\stackrel{}{\gamma }_1`$ and $`\stackrel{}{\gamma }_2`$, that are both subcritical ($`\gamma _1<1`$ and $`\gamma _2<1`$) or where one is critical and the other subcritical ($`\gamma _1<1`$ and $`\gamma _2>1`$). For small deformations ($`\gamma 1`$), we have that $`S^\mathrm{I}/S^\mathrm{S}`$ $``$ $`1+2\gamma ^2+𝒪(\gamma ^4)`$ (126) $`ϵ^\mathrm{I}/ϵ^\mathrm{S}`$ $``$ $`14e^S\mathrm{cos}\alpha \gamma +4\left(3(e^S)^2\mathrm{cos}^2\alpha 1\right)\gamma ^2+𝒪(\gamma ^3)`$ (127) so that the images almost lie on a parabola centered on $`(S^\mathrm{I},e^\mathrm{I})=(S^\mathrm{S},e^\mathrm{S})`$. In such a weak field, two images of the same object will have almost the same surface but can have very different ellipticities. ### C Morphology of possible sources of a given image We now ask the reverse question: given an image $`(S^\mathrm{I},e^\mathrm{I})`$, from which set of sources can it be the image? Following the description and notations of the two previous sections, we now set $$M^\mathrm{I}=\left(\begin{array}{cc}\lambda _+& 0\\ 0& \lambda _{}\end{array}\right),𝒜=\left(\begin{array}{cc}1\gamma _1& \gamma _2\\ \gamma _2& 1+\gamma _1\end{array}\right)$$ (128) from which we deduce that, since $`M^\mathrm{S}=𝒜M^\mathrm{I}𝒜`$, the shape of the source is related to the one of its images by $`S^\mathrm{S}`$ $`=`$ $`(1\gamma ^2)^2S^\mathrm{I}`$ (129) $`ϵ^\mathrm{S}`$ $`=`$ $`{\displaystyle \frac{(1\gamma ^2)^2}{(1+\gamma ^22\gamma e^\mathrm{I}\mathrm{cos}\alpha )^2}}ϵ^\mathrm{I}.`$ (130) As a first exercise, we depict on figure 9 the sources $`(S^\mathrm{S},e^\mathrm{S})`$ of a circular image ($`e^\mathrm{I}=0`$). A priories on the ellipticity of the sources (117) and on the strength of the deformation $`\gamma `$ may enable us to extract from such a plot informations about the source of a circular image. In the general case where $`e^\mathrm{I}0`$, we can reconstruct all the morphologies of the source that can lead to the observed image. As an example we depict such sets in the two cases where $`e^\mathrm{I}=0.5`$ and $`e^\mathrm{I}=0.25`$ respectively on figures 10 and 11. We note that all the curves pass through the point $`(S^\mathrm{S},e^\mathrm{S})=(0,1)`$ which is reaches when $`\gamma =1`$, i.e. on the critical line. Again, one can sort out that any object can be the image of two circular sources obtained by the two transformations $$\gamma _\pm =\frac{e^\mathrm{I}}{1\sqrt{ϵ^\mathrm{I}}}\alpha 0[2\pi ]$$ (131) and the measure of $`(S^\mathrm{I},e^\mathrm{I})`$ enables to determine $`\gamma _\pm `$ and the two surfaces $`S_\pm ^\mathrm{S}`$. If one measures the shape of two images $`(S_1^\mathrm{I},e_1^\mathrm{I})`$ and $`(S_2^\mathrm{I},e_2^\mathrm{I})`$ one can reconstruct, as in figures 10 and 11, the set of morphologies of their possible sources. Given bounds on $`e^\mathrm{S}`$, as in (117), and on $`\gamma `$ one can figure out graphically if these two observations are likely to be images of the same object. Indeed for very weak deformations ($`\gamma 1`$) one gets that two images of the same objet must have almost the same surface but can have very different ellipticity (see figure 12 where we have plotted the source shapes of objects of different shape for small deformation). ## VI Conclusion We have considered the general lensing properties of objects such as cosmic strings where relativistic motions and non-trivial equation of state induce metric perturbations of all sorts. We demonstrated that the deformation field of a string system on the image plane is the same as the one of a static linear distribution of matter projected on the photon trajectory. A consequence of this result is that the deformation field has a vanishing convergence ($`\kappa =0`$) everywhere but on the projection of the intersection of the observer past light cone and the string worldsheet. We explicitly illustrate this result with the case of a circular cosmic string loop in a plane perpendicular to the line-of-sight for which we generalize the results found in to string with any equation of state. We also showed that (i) the fact that the deformation field outside the loop is equivalent to the one obtained by a massive point and (ii) that a light ray passing inside the loop is not deflected are due to the energy conservation of the string. We also paid attention to the validity of thin lens approximation for this unusual lens system. This approximation is discussed in detail through the case of a static tilted straight cosmic string. It lead us to point out that for string with a general equation of state, the deflection may be more important than for a Goto–Nambu string, this being understood by the fact that the stress–energy tensor of a general cosmic string can always be decomposed as the superposition of a Goto–Nambu string and a lineic distribution of non–relativistic matter. The deflection is then due to the combined effect of the deficit angle of the Goto–Nambu string and of the curvature induced by the lineic distribution of matter. We also studied general phenomenological consequences of deformation fields with zero convergence on multiple image systems, the main goal being to be able to assess if two images are likely to form an image pair of the same source. For that purpose, we derived all the image shapes of a given source as well as all the source shapes of a given image. These results may serve as a groundwork for the elaboration of string detection strategies. We are aware that this latter part is limited by the fact that we only took advantage of the zero-convergence property. In practical more intricate distortion properties of the images are likely to be useful. This is one of the aim of the companion paper where more quantitative phenomenological properties are presented that take into account local string energy fluctuations. ## Acknowledgements We would like to thank Ruth Durrer, Yannick Mellier, Patrick Peter and Albert Stebbins for discussions on the subject and the anonymous referee for his numerous comments and corrections.
warning/0004/hep-ph0004052.html	ar5iv	text	# 1 Introduction ## 1 Introduction This paper is devoted to further study of theoretical aspects of the deconfining temperature phase transition in nonabelian gauge theories. It is an immediate continuation of our earlier work . In we showed that the deconfining phase transition in the pure Yang Mills theory is characterised by the change of behaviour of the ’t Hooft loop operator $`V(C)`$. In the “cold” phase the ’t Hooft loop has a perimeter law behaviour $`<V(C)>\mathrm{exp}\{aP(C)\}`$, while in the “hot” phase it has an area law behaviour $`<V(C)>\mathrm{exp}\{\alpha S(C)\}`$. In the present paper we want to sharpen somewhat this observation and further discuss related questions. We wish to point out that $`V`$ is in fact an order parameter which probes the breaking of a physical symmetry of the Yang Mills theory. The symmetry in question is the magnetic $`Z_N`$ symmetry discussed by ’t Hooft . The deconfining phase transition is therefore characterized by the change in the mode of realization of a global $`Z_N`$ symmetry: the symmetry is broken spontaneously in the “cold” phase while it is restored in the “hot“ phase. The previous two paragraphs may sound at first like a red herring. After all an order parameter for the deconfining phase transition as well as a related $`Z_N`$ symmetry have been discussed for many years. The order parameter in question is the free energy of an external static colour source in the fundamental representation: the Polyakov line $`P=\mathrm{Tr}P\mathrm{exp}\{ig_0^\beta 𝑑tA_0\}`$. The $`Z_N`$ symmetry is the transformation $`P\mathrm{exp}\{i\frac{2\pi }{N}\}P`$. We will refer to this transformation as the electric $`Z_N`$. There is however a great difference between the physical nature of $`P`$ and $`V`$ and the associated $`Z_N`$ symmetries. The operator $`V`$ is a canonical operator in the physical Hilbert space of the Yang Mills theory. The magnetic $`Z_N`$ symmetry similarly is a transformation that acts on quantum states in the physical Hilbert space. On the other hand $`P`$ has a very different status. It is not an operator in the Hilbert space and as such not a canonical order parameter. It appears as an auxilliary object when projecting onto gauge invariant physical subspace of the Hilbert space. The “electric” $`Z_N`$ \- the operation that transforms $`P`$ by multiplying it by a phase - similarly is not a canonical symmetry. These issues were discussed in detail in . There is no transformation of states in the physical Hilbert space that is related to this “symmetry”, although it is indeed a symmetry of the Euclidean path integral representing the statistical sum. This is not to say of course that $`P`$ and electric $`Z_N`$ are useless concepts. The standard effective action, defined by the constrained path integral $$\mathrm{exp}S_{eff}(P)=DA_0\delta (PP(A_0))\mathrm{exp}S(A)$$ (1) is gauge invariant. It is instrumental in computing the vortex expectation value. The way the electric $`Z(N)`$ symmetry is realized in $`S_{eff}`$ is also related to the behaviour of the order parameter of the magnetic $`Z(N)`$. We will discuss this in detail in section 3 . However if one wants to describe the deconfinement phase transition in terms of a canonical order parameter in the same way as the Ising transition is described in terms of magnetisation, one should zero in on $`V`$ rather than on $`P`$ and should study the magnetic $`Z_N`$ symmetry rather than electric $`Z_N`$. This is what we intend to do in this paper. The action of the magnetic $`Z_N`$ symmetry is very different in 2+1 and 3+1 dimensional cases. In 2+1 dimensions it acts very much like usual global symmetry in a scalar theory with the order parameter being a scalar vortex field. In 3+1 dimensions the symmetry acts not like a standard global symetry - its “charge” is an integral over a two dimensional spacelike surface rather than over the whole of the three dimensional space<sup>1</sup><sup>1</sup>1These type of symmetries nowadays are frequently discussed in the context of “M - theory”.. As a consequence its order parameter is not a local field but rather a magnetic vortex stretching over macroscopic distances. It is therefore convenient to start the discussion with the three dimensional gauge theories and to present all the arguments in this case. The generalization of appropriate aspects of this discussion to 3+1 dimensions will be given in the last part of every section. The plan of this paper is the following. In Section 2 we recap the definition of the ’t Hooft loop and its 2+1 dimensional analog - the magnetic vortex operator. We formulate the arguments for the existence of the magnetic $`Z_N`$ symmetry in theories without fundamental matter fields. We also show by explicit construction that the generator of this symmetry in the pure gluodynamics is none other than the spatial Wilson loop. In Section 3 we discuss the relation between the behaviour of the ’t Hooft loop and the realization of the magnetic $`Z_N`$ in the ground state of the theory. We demonstrate that the mode of the realization of the symmetry changes at the deconfining phase transition, while spontaneously broken at low temperature the symmetry is restored above the phase transition. In Section 4 we present in a toy model a simple physical picture explaining how the behaviour of spatial Wilson loop discriminates at zero temperature between the phases with broken and unbroken magnetic $`Z_N`$. In the phase where the $`Z_N`$ symmetry is broken, $`W`$ must have an area law while in the case of unbroken $`Z_N`$ it must have perimeter law. We explain why this argument does not generalize to the high temperature phase and thus why the area law behaviour of the Wilson loop in the hot phase is consistent with restoration of the magnetic $`Z_N`$ symmetry. Finally in Section 5 we conclude with a short discussion. ## 2 The magnetic $`Z_N`$ symmetry and the ’t Hooft loop operator. In this section we discuss the notion of the magnetic $`Z_N`$ symmetry and its order parameter - ’t Hooft loop, or magnetic vortex operator. Most of the material contained here is not new and, perhaps with the exception of explicit identification of the $`Z_N`$ generator with the spatial Wilson loop, is contained in . At the risk of being repetitive we have decided nevertheless to include this extended introductory part, since we feel that the concept of magnetic $`Z_N`$ symmetry is not widely appreciated in the community. The $`Z_N`$ symmetry structure is the basis for our discussion of the deconfining phase transition in the following sections. Let us start by recalling the argument due to ’t Hooft that a nonabelian $`SU(N)`$ gauge theory with charged fields in adjoint representation posesses a global $`Z_N`$ symmetry . We discuss the 2+1 dimensional case first. Consider a theory with several adjoint Higgs fields so that varying parameters in the Higgs sector the $`SU(N)`$ gauge symmetry can be broken completely. In this phase the perturbative spectrum will contain the usual massive “gluons” and Higgs particles. However in addition to that there will be heavy stable magnetic vortices. Those are the analogs of Abrikosov-Nielsen-Olesen vortices in the superconductours and they must be stable by virtue of the following topological argument. The vortex configuration away from the vortex core has all the fields in the pure gauge configuration $$H^\alpha (x)=U(x)h^\alpha ,A^\mu =iU^\mu U^{}$$ (2) Here the index $`\alpha `$ labels the scalar fields in the theory, $`h^\alpha `$ are the constant vacuum expectation values of these fields, and $`U(x)`$ is a unitary matrix. As one goes around the location of the vortex in space, the matrix $`U`$ winds nontrivially in the gauge group. This is possible, since the gauge group in the theory without fundamental fields is $`SU(N)/Z_N`$ and it has a nonvanishing first homotopy group $`\mathrm{\Pi }_1(SU(N)/Z_N)=Z_N`$. Practically it means that when going around the vortex location full circle, $`U`$ does not return to the same $`SU(N)`$ group element $`U_0`$, but rather ends up at $`\mathrm{exp}\{i\frac{2\pi }{N}\}U_0`$. Adjoint fields do not feel this type of discontinuity in $`U`$ and therefore the energy of such a configuration is finite. Since such a configuration can not be smoothly deformed into a trivial one, a single vortex is stable. Processes involving annihilation of N such vortices into vacuum are allowed since N-vortex configurations are topologically trivial. One can of course find explicit vortex solutions once the Higgs potential is specified. As any other semiclassical solution in the weak coupling limit the energy of such a vortex is inversely proportional to the gauge coupling constant and therefore very large. One is therefore in a situation where the spectrum of the theory contains a stable particle even though its mass is much higher than masses of many other particles (gauge and Higgs bosons) and the phase space for its decay into these particles is enormous. The only possible reason for the existence of such a heavy stable particle is that it must carry a conserved quantum number. The theory therefore must possess a global symmetry which is unbroken in the completely higgsed phase. The symmetry group must be $`Z_N`$ since the number of vortices is only conserved modulo $`N`$. Now imagine changing smoothly the parameters in the Higgs sector so that the expectation values of the Higgs fields become smaller and smaller, and finally the theory undergoes a phase transition into the confining phase. One can further change the parameters so that the adjoint scalars become heavy and eventually decouple completely from the glue. This limiting process does not change the topology of the gauge group and therefore does not change the symmetry content of the theory. We conclude that the pure Yang-Mills theory also posesses a $`Z_N`$ symmetry. Of course since the confining phase is separated from the completely Higgsed phase by a phase transition one may expect that the $`Z_N`$ symmetry in the confining phase is represented differently. In fact the original paper of ’t Hooft as well as subsequent work convincingly argued that in the confining phase the $`Z_N`$ symmetry is spontaneously broken and this breaking is related to the confinement phenomenon. The physical considerations given above can be put on firmer formal basis. In particular one can construct explicitly the generator of the $`Z_N`$ as well as the order parameter associated with it \- the operator that creates the magnetic vortex . We will now describe this construction. ### 2.1 The Abelian case Consider first an Abelian gauge theory. In this case the homotopy group is $`Z`$ and therefore we expect the $`U(1)`$ rather than $`Z_N`$ magnetic symmetry. It is in fact absolutely straightforward to identify the relevant charge. It is none other than the magnetic flux through the equal time plane, with the associated conserved current being the dual of the electromagnetic field strength $$\mathrm{\Phi }=d^2xB(x),^\mu \stackrel{~}{F}_\mu =0$$ (3) The current conservation is insured by the Bianchi identity. A group element of the $`U(1)`$ magnetic symmetry group is $`\mathrm{exp}\{i\alpha \mathrm{\Phi }\}`$ for any value of $`\alpha `$. A local order parameter - a local field $`V(x)`$ which carries the magnetic charge - is also readily constructed. It has a form of the singular gauge transformation operator with the singularity at the point $`x`$ $$V(x)=\mathrm{exp}\frac{i}{g}d^2y\left[ϵ_{ij}\frac{(xy)_j}{(xy)^2}E_i(y)+\mathrm{\Theta }(xy)J_0(y)\right]$$ (4) where $`\mathrm{\Theta }(xy)`$ is the polar angle function and $`J_0`$ is the electric charge density of whatever matter fields are present in the theory. The cut discontinuity in the function $`\mathrm{\Theta }`$ is not physical and can be chosen parallel to the horizontal axis. Using the Gauss’ law constraint this can be cast in a different form, which we will find more convenient for our discussion $$V(x)=\mathrm{exp}\frac{2\pi i}{g}_C𝑑y^iϵ_{ij}E_i(y)$$ (5) where the integration goes along the cut of the function $`\mathrm{\Theta }`$ which starts at the point $`x`$ and goes to spatial infinity. The operator does not depend on where precisely one chooses the cut to lie. To see this, note that changing the position of the cut $`C`$ to $`C^{}`$ adds to the phase $`\frac{2\pi }{g}_Sd^2x_iE^i`$ where $`S`$ is the area bounded by $`CC^{}`$. In the theory we consider only charged particles with charges multiples of $`g`$ are present. Therefore the charge within any closed area is a multiple integer of the gauge coupling $`_Sd^2x_iE^i=gn`$ and the extra phase factor is always unity. The meaning of the operator $`V`$ is very simple. From the commutation relation $$V(x)B(y)V^{}(x)=B(y)+\frac{2\pi }{g}\delta ^2(xy)$$ (6) it is obvious that $`V`$ creates a pointlike magnetic vortex of flux $`2\pi /g`$. Despite its nonlocal appearance the operator $`V`$ can be proven to be a local Lorentz scalar field. The locality is the consequence of the fact that $`V(x)`$ commutes with any local gauge invariant operator in the theory $`O(y)`$ except when $`x=y`$. This is due to the coefficient $`2\pi /g`$ in the exponential which ensures that the Aharonov-Bohm phase of the vortex created by $`V`$ and any dynamical charged particle present in the theory vanishes. Eqs.(3,5) formalize the physical arguments of ’t Hooft in the abelian case. ### 2.2 The non-Abelian case at weak coupling. Let us now move onto the analogous construction for nonabelian theories. Ultimately we are interested in the pure Yang - Mills theory. It is however illuminating to start with the theory with an adjoint Higgs field and take the decoupling limit explicitly later. For simplicity we discuss the $`SU(2)`$ gauge theory. Consider the Georgi-Glashow model - $`SU(2)`$ gauge theory with an adjoint Higgs field. $$=\frac{1}{4}F_{\mu \nu }^aF^{a\mu \nu }+\frac{1}{2}(𝒟_\mu ^{ab}H^b)^2+\stackrel{~}{\mu }^2H^2\stackrel{~}{\lambda }(H^2)^2$$ (7) where $$𝒟_\mu ^{ab}H^b=_\mu H^agf^{abc}A_\mu ^bH^c$$ (8) At large and positive $`\stackrel{~}{\mu }^2`$ the model is weakly coupled. The $`SU(2)`$ gauge symmetry is broken down to $`U(1)`$ and the Higgs mechanism takes place. Two gauge bosons, $`W^\pm `$, acquire a mass, while the third one, the “photon”, remains massless to all orders in perturbation theory. The theory in this region of parameter space resembles very much electrodynamics with vector charged fields. The Abelian construction can therefore be repeated. The $`SU(2)`$ gauge invariant analog of the conserved dual field strength is $$\stackrel{~}{F}^\mu =\frac{1}{2}[ϵ_{\mu \nu \lambda }F_{\nu \lambda }^an^a\frac{1}{g}ϵ^{\mu \nu \lambda }ϵ^{abc}n_a(𝒟_\nu n)^b(𝒟_\lambda n)^c]$$ (9) where $`n^a\frac{H^a}{\|H\|}`$ is the unit vector in the direction of the Higgs field. Classically this current satisfies the conservation equation $$^\mu \stackrel{~}{F}_\mu =0$$ (10) The easiest way to see this is to choose a unitary gauge of the form $`H^a(x)=H(x)\delta ^{a3}`$. In this gauge $`\stackrel{~}{F}`$ is equal to the abelian part of the dual field strength in the third direction in colour space. Its conservation then follows by the Bianchi identity. Thus classically the theory has a conserved $`U(1)`$ magnetic charge $`\mathrm{\Phi }=d^2x\stackrel{~}{F}_0`$ just like QED. However the unitary gauge can not be imposed at the points where $`H`$ vanishes, which necessarilly happens in the core of an ’t Hooft-Polyakov monopole. It is well known of course that the monopoles are the most important nonperturbative configurations in this model. Their presence leads to a nonvanishing small mass for the photon as well as to confinement of the charged gauge bosons with a tiny nonperturbative string tension. As far as the monopole effects on the magnetic flux, their presence leads to a quantum anomaly in the conservation equation (10). As a result only the discrete $`Z_2`$ subgroup of the transformation group generated by $`\mathrm{\Phi }`$ remains unbroken in the quantum theory. The detailed discussion of this anomaly, the residual $`Z_2`$ symmetry and their relation to monopoles is given in . The order parameter for the magnetic $`Z_2`$ symmetry is constructed analogously to QED as a singular gauge transformation generated by the gauge invariant electric charge operator $$J^\mu =ϵ^{\mu \nu \lambda }_\nu (\stackrel{~}{F}_\lambda ^an^a),Q=d^2xJ_0(x)$$ (11) Explicitly $`V(x)`$ $`=`$ $`\mathrm{exp}{\displaystyle \frac{i}{g}}{\displaystyle d^2y\left[ϵ_{ij}\frac{(xy)_j}{(xy)^2}n^a(y)E_i^a(y)+\mathrm{\Theta }(xy)J_0(y)\right]}`$ (12) $`=`$ $`\mathrm{exp}{\displaystyle \frac{2\pi i}{g}}{\displaystyle _C}𝑑y^iϵ_{ij}n^aE_j^a(y)`$ One can think of it as a singular $`SU(2)`$ gauge transformation with the field dependent gauge function $$\lambda ^a(y)=\frac{1}{g}\mathrm{\Theta }(xy)n^a(\stackrel{}{y})$$ (13) This field dependence of the gauge function ensures the gauge invariance of the operator $`V`$. Just like in QED it can be shown, that the operator $`V`$ is a local scalar field. Again like in QED, the vortex operator $`V`$ is a local eigenoperator of the abelian magnetic field $`B(x)=\stackrel{~}{F}_0`$. $$[V(x),B(y)]=\frac{2\pi }{g}V(x)\delta ^2(xy)$$ (14) That is to say, when acting on a state it creates a pointlike magnetic vortex which carries a quantized unit of magnetic flux. The $`Z_2`$ magnetic symmetry transformation is generated by the operator $$U=\mathrm{exp}\{i\frac{g}{2}\mathrm{\Phi }\}$$ (15) and acts on the vortex field $`V`$ as a phase rotation by $`\pi `$ $$e^{i\frac{g}{2}\mathrm{\Phi }}V(x)e^{i\frac{g}{2}\mathrm{\Phi }}=V(x)$$ (16) An operator closely related to $`U`$ and which will be of interest to us in the following, is the generator of the magnetic $`Z_2`$ transformation only inside some closed contour $`C`$ $$U(C)=\mathrm{exp}\{i\frac{g}{2}_Sd^2xB(x)\}$$ (17) where the integration is over the area $`S`$ bounded by $`C`$. The analog of the commutator eq.(16) for this operator is $`U_CV(x)U_C^{}=`$ $`V(x)`$ $`,xS`$ (18) $`V(x)`$ $`,xS`$ Taking the contour $`C`$ to run at infinity $`U_C`$ becomes the generator of $`Z_2`$. We now have the explicit realization of the magnetic $`Z_2`$ symmetry in the Georgi-Glashow model. ### 2.3 The pure gauge theory. Our next step is to move on to the pure Yang Mills theory. This is achieved by smoothly varying the $`\stackrel{~}{\mu }^2`$ coefficient in the Lagrangian so that the coeficient of the mass term of the Higgs field becomes positive and eventually arbitrarily large. It is well known that in this model the weakly coupled Higgs regime and strongly coupled confining regime are not separated by a phase transition. The pure Yang Mills limit in this model is therefore smooth. In the pure Yang Mills limit the expressions eq.(9,12,17) have to be taken with care. When the mass of the Higgs field is very large, the configurations that dominate the path integral are those with very small value of the modulus of the Higgs field $`\|H\|1/M`$. The modulus of the Higgs field in turn controls the fluctuations of the unit vector $`n^a`$, since the kinetic term for $`n`$ in the Lagrangian is $`\|H\|^2(D_\mu n)^2`$. Thus as the mass of the Higgs field increases the fluctuations of $`n`$ grow in both, amplitude and frequency and the magnetic field operator $`B`$ as defined in eq.(9) fluctuates wildly. This situation is of course not unusual. It happens whenever one wants to consider in the effective low energy theory an operator which explicitly depends on fast, high energy variables. The standard way to deal with it is to integrate over the fast variables. There could be two possible outcomes of this procedure. Either the operator in question becomes trivial (if it depends strongly on the fast variables), or its reduced version is well defined and regular on the low energy Hilbert space. The “magnetic field” operator $`B`$ in eq.(9) is obviously of the first type. Since in the pure Yang Mills limit all the orientations of $`n^a`$ are equally probable, integrating over the Higgs field at fixed $`A_\mu `$ will lead to vanishing of $`B`$. However what interests us is not so much the magnetic field but rather the generator of the magnetic $`Z_2`$ transformation $`U_C`$ of eq.(17). In the pure Yang-Mills limit we are thus lead to consider the operator $`U_C=lim_{H0}{\displaystyle Dn^a}`$ $`\mathrm{exp}\left\{\|H\|^2(\stackrel{}{D}n_a)^2\right\}`$ $`\mathrm{exp}\left\{i{\displaystyle \frac{g}{4}}{\displaystyle _C}d^2x\left(ϵ_{ij}F_{ij}^an^a{\displaystyle \frac{1}{g}}ϵ^{ij}ϵ^{abc}n_a(𝒟_in)^b(𝒟_jn)^c\right)\right\}`$ The weight for the integration over $`n`$ is the kinetic term for the isovector $`n_a`$. As was noted before the action does not depend on $`n^a`$ in the YM limit. This term however regulates the integral and we keep it for this reason. This operator may look somewhat unfamiliar at first sight. However in a remarkable paper Diakonov and Petrov showed that eq.(2.3) is equal to the trace of the fundamental Wilson loop along the contour $`C`$<sup>2</sup><sup>2</sup>2We note that Dyakonov and Petrov had to introduce a regulator to define the path integral over $`n`$. The regulator they required was precisely of the same form as in eq.(2.3). It is pleasing to see that this regulator appears naturally in our approach as the remnant of the kinetic term of the Higgs field.. $$U_C=W_C\mathrm{Tr}𝒫\mathrm{exp}\left\{\mathrm{ig}_\mathrm{C}\mathrm{dl}^\mathrm{i}\mathrm{A}^\mathrm{i}\right\}$$ (20) We conclude, that in the pure Yang-Mills theory the generator of the magnetic $`Z_2`$ symmetry is the fundamental spatial Wilson loop along the boundary of the spatial plane. There is a slight subtlety here that may be worth mentioning. The generator of a unitary transformation should be a unitary operator. The trace of the fundamental Wilson loop on the other hand is not unitary. One should therefore strictly speaking consider instead a unitarized Wilson loop $`\stackrel{~}{W}=\frac{W}{\sqrt{WW^{}}}`$. However the factor between the two operators $`\sqrt{WW^{}}`$ is an operator that is only sensitive to behaviour of the fields at infinity. It commutes with all physical local operators $`O(x)`$ unless $`x\mathrm{}`$. In this it is very different from the Wilson loop itself, which has a nontrivial commutator with vortex operators $`V(x)`$ at all values of $`x`$. Since the correlators of all gauge invariant local fields in the pure Yang Mills theory are massive and therefore short range, the operator $`\sqrt{WW^{}}`$ must be a constant operator on all finite energy states. The difference between $`W`$ and $`\stackrel{~}{W}`$ is therefore a trivial constant factor and we will not bother with it in the following. Perhaps of more concern is the difference between $`W_C`$ and $`\stackrel{~}{W}_C`$ when the contour $`C`$ is not at infinity. However here again the factor between the two operators $`\sqrt{W_CW_C^{}}`$ is only sensitive to physical degrees of freedom on the contour $`C`$ and not inside it. Due to its presence the vacuum averages of $`W_C`$ and $`\stackrel{~}{W}_C`$ may differ at most by a factor which has a perimeter behaviour $`<W_C>=\mathrm{exp}\{mP(C)\}<\stackrel{~}{W}_C>`$ where $`P(C)`$ is a perimeter of $`C`$. The question we will be interested in is whether $`<W_C>`$ has a perimeter or area behaviour. As far as the answer to this question is concerned $`W_C`$ and $`\stackrel{~}{W}_C`$ are completely equivalent, and we will not make distinction between them. In the rest of this paper we will therefore refer to $`W`$ as the generator of $`Z_2`$ keeping this little caveat in mind. Next we consider the vortex operator eq.(12). Again we have to integrate it over the orientations of the unit vector $`n^a`$. This integration in fact is equivalent to averaging over the gauge group. Following one can write $`n_a`$ in terms of the SU(2) gauge transformation matrix $`\mathrm{\Omega }`$. $$\stackrel{}{n}=\frac{1}{2}\mathrm{Tr}\mathrm{\Omega }\tau \mathrm{\Omega }^{}\tau _3$$ (21) The vortex operator in the pure gluodynamics limit then becomes $$\stackrel{~}{V}(x)=D\mathrm{\Omega }\mathrm{exp}\frac{2\pi i}{g}_C𝑑y_iϵ_{ij}\mathrm{Tr}\mathrm{\Omega }E_j\mathrm{\Omega }^{}\tau _3$$ (22) This form makes it explicit that $`\stackrel{~}{V}(x)`$ is defined as the gauge singlet part of the following, apparently non gauge invariant operator $$V(x)=\mathrm{exp}\frac{2\pi i}{g}_C𝑑y^iϵ_{ij}E_i^3(y)$$ (23) The integration over $`\mathrm{\Omega }`$ obviously projects out the gauge singlet part of $`V`$. In the present case however this projection is redundant. This is because even though $`V`$ itself is not gauge invariant, when acting on a physical state it transforms it into another physical state<sup>3</sup><sup>3</sup>3This is not a trivial statement, since a generic nongauge invariant operator has nonvanishing matrix elements between the physical and an unphysical sectors.. By physical states we mean the states which satisfy the Gauss’ constraint in the pure Yang-Mills theory. This property of $`V`$ was noticed by ’t Hooft . To show this let us consider $`V(x)`$ as defined in eq.(23) and its gauge transform $`V_\mathrm{\Omega }=\mathrm{\Omega }^{}V\mathrm{\Omega }`$ where $`\mathrm{\Omega }`$ is an arbitrary nonsingular gauge transformation operator. The wave functional of any physical state depends only on gauge invariant characteristics of the vector potential, i.e. only on the values of Wilson loops over all possible contours. $$\mathrm{\Psi }[A_i]=\mathrm{\Psi }[\{W(C)\}]$$ (24) Acting on this state by the operators $`V`$ and $`V_\mathrm{\Omega }`$ respectively we obtain $`V\|\mathrm{\Psi }>`$ $`=`$ $`\mathrm{\Psi }_V[A_i]=\mathrm{\Psi }[\{VW(C)V^{}\}]`$ $`V_\mathrm{\Omega }\|\mathrm{\Psi }>`$ $`=`$ $`\mathrm{\Psi }_V^\mathrm{\Omega }[A_i]=\mathrm{\Psi }[\{V_\mathrm{\Omega }W(C)V_\mathrm{\Omega }^{}\}]`$ (25) It is however easy to see that the action of $`V(x)`$ and $`V_\mathrm{\Omega }(x)`$ on the Wilson loop is identical - they both multiply it by the centergroup phase (which stays unaffected by $`\mathrm{\Omega }`$) if $`x`$ is inside $`C`$ and do nothing otherwise. Therefore we see that $$V\|\mathrm{\Psi }>=V_\mathrm{\Omega }\|\mathrm{\Psi }>$$ (26) for any physical state $`\mathrm{\Psi }`$. Thus we have $$\mathrm{\Omega }V\|\mathrm{\Psi }>=\mathrm{\Omega }V\mathrm{\Omega }^{}\|\mathrm{\Psi }>=V\|\mathrm{\Psi }>$$ (27) where the first equality follows from the fact that a physical state is invariant under action of any gauge transformation $`\mathrm{\Omega }`$ and the second equality follows from eq.(26). But this equation is nothing but the statement that the state $`V\|\mathrm{\Psi }>`$ is physical, i.e. invariant under any nonsingular gauge transformation. We have therefore proved that when acting on a physical state the vortex operator creates another physical state. For an operator of this type the gauge invariant projection only affects its matrix elements between unphysical states. Since we are only interested in calculating correlators of $`V`$ between physical states, the gauge projection is redundant and we can freely use $`V`$ rather than $`\stackrel{~}{V}`$ to represent the vortex operator. It is instructive to note that this propery is not shared by the Wilson loop. One can in fact represent the Wilson loop as a singlet gauge projection of a simple Abelian loop operator. The second exponential in eq.(2.3) can be written as $$\mathrm{exp}\left\{i\frac{g}{2}_C𝑑l^iA_a^in^a\frac{i}{2}d^2xϵ_{ij}ϵ^{abc}n_a_in_b_jn_c\right\}$$ (28) Using eq.(21) we can rewrite the integral in eq.(2.3) -omitting the regulating kinetic piece- as: $$W_C=D\mathrm{\Omega }\mathrm{exp}\left\{i\frac{g}{2}\mathrm{Tr}\tau _3\left(\mathrm{\Omega }A^i\mathrm{\Omega }^{}+i\mathrm{\Omega }^i\mathrm{\Omega }^{}\right)𝑑l^i\right\}$$ (29) The Wilson loop is therefore the gauge singlet part of the Abelian loop $$U_C^A=\mathrm{exp}i\frac{g}{2}\mathrm{Tr}A^i\tau _3𝑑l^i$$ (30) The matrix elements of $`W_C`$ and $`U_C^A`$ on physical subspace therefore are the same. However $`U_C^A`$ as opposed to $`V`$ does have nonvanishing nondiagonal matrix elements, that is matrix elements between the physical and the unphysical sectors. It therefore can not be used instead of $`W_C`$ in gauge theory calculations. For example non gauge invariant states will contribute as intermediate states in the calculation of quantities like the correlation function $`<U_{C1}^AU_{C_2}^A>`$, while their contribution vanishes in similar correlators which involve the Wilson loop. The generalization of the preceding discussion to $`SU(N)`$ gauge theories is straightforward. Once again one can start with the Georgi-Glashow like model, where the $`SU(N)`$ is higgsed to $`U(1)^{(N1)}`$<sup>4</sup><sup>4</sup>4In $`SU(N)`$ theories with $`N>2`$ there in principle can be phases separated from each other due to spontaneous breaking of some global symmetries. For instance the $`SU(3)`$ gauge theory with adjoint matter has a phase with spontaneously broken charge conjugation invariance . Still even in this phase the confining properties are the same as in the strongly coupled pure Yang-Mills theory, with the Wilson loop having an area law.. The construction of the vortex operator and the generator of $`Z_N`$ in this case is very similar and the details are given in . Taking the mass of the Higgs field to infinity again projects the generator onto the trace of the fundamental Wilson loop. The vortex operator can be taken as $$V(x)=\mathrm{exp}\{\frac{4\pi i}{gN}_Cdy^iϵ_{ij}\mathrm{Tr}(\mathrm{YE}_\mathrm{i}(\mathrm{y}))$$ (31) where the hypercharge generator $`Y`$ is defined as $$Y=\mathrm{diag}(1,1,\mathrm{},(N1))$$ (32) and the electric field is taken in the matrix notation $`E_i=\lambda ^aE_i^a`$ with $`\lambda ^a`$ \- the $`SU(N)`$ generator matrices in the fundamental representation. ### 2.4 Generalization to 3+1 dimensions To conclude this section we discuss how the magnetic symmetry structure generalizes to four dimensions. The conserved $`Z_N`$ generator in the Georgi-Glashow model is defined through $$U_S=\mathrm{exp}\left\{i\frac{g}{2}_Sd^2S^i\left(B_i^an^a\frac{1}{g}ϵ^{ijk}ϵ^{abc}n_a(𝒟_jn)^b(𝒟_kn)^c\right)\right\}$$ (33) Although the definition of $`U`$ contains explicitly the surface $`S`$ through which the abelian magnetic flux is integrated, the operator in fact does not depend on $`S`$ but is specified completely by its boundary. This is because changing $`S`$ changes the phase of $`U`$ by the magnetic flux through the closed surface. The only dynamical objects that carry magnetic flux in the theory are ’t Hooft-Polyakov monopoles. Since their flux is quantized in units of $`4\pi /g`$ the change in the phase is always a multiple integer of $`2\pi `$. In the pure Yang-Mills limit the operator $`U_S`$ again reduces to the trace of the fundamental Wilson loop along the boundary of $`S`$. Taking the contour to infinity defines the generator of magnetic $`Z_N`$. As we have already noted, this charge is a little unusual in that it is defined as a surface integral rather than a volume integral. As a result the order parameter for this symmetry transformation is not a local but rather a stringy field. This is of course just a restatement of the fact that magnetic vortices in 3+1 dimensions are stringlike objects. The operator that creates a vortex can still be defined in a way similar to 2+1 dimensions. Skipping the intermediate steps which we went through in the previous discussion we give the final result for the pure Yang Mills $`SU(N)`$ gauge theory. The magnetic vortex along the curve $`C`$ is created by the following operator of the ”singular gauge transformation”<sup>5</sup><sup>5</sup>5The derivative term $`^i\omega `$ in this expression should be understood to contain only the smooth part of the derivative and to exclude the contribution due to the discontinuity of $`\omega `$ on the surface $`S`$. $$V(C)=\mathrm{exp}\{\frac{i}{gN}d^3x\mathrm{Tr}(\mathrm{D}^\mathrm{i}\omega _\mathrm{C}\mathrm{Y})\mathrm{E}^\mathrm{i}\}=\mathrm{exp}\{\frac{4\pi \mathrm{i}}{\mathrm{gN}}_\mathrm{S}\mathrm{d}^2\mathrm{S}^\mathrm{i}\mathrm{Tr}(\mathrm{YE}^\mathrm{i})\}$$ (34) with $`\omega _C(x)`$, the singular gauge function which is equal to the solid angle subtended by $`C`$ as seen from the point $`x`$. The function $`\omega `$ is continuos everywhere, except on a surface $`S`$ bounded by $`C`$, where it jumps by $`4\pi `$. Other than the fact that $`S`$ is bounded by $`C`$, its location is arbitrary. The vortex loop and the spatial Wilson loop satisfy the ’t Hooft algebra $$V^{}(C)W(C^{})V(C)=e^{\frac{2\pi i}{N}n(C,C^{})}W(C^{})$$ (35) where $`n(C,C^{})`$ is the linking number of the curves $`C`$ and $`C^{}`$. One can consider closed contours $`C`$ or infinite contours that run through the whole system. For an infinite contour $`C`$ and the Wilson loop along the spatial boundary of the system the linking number is always unity. The $`V(C)`$ for an infinite loop is therefore an eigenoperator of the $`Z_N`$ magnetic symmetry and is the analog of the vortex operator $`V(x)`$ in 2+1 dimensions. Any closed vortex loop of fixed size commutes with the Wilson loop if the contour $`C^{}`$ is very large. Such a closed loop is thus an analog of the vortex-antivortex correlator $`V(x)V^{}(y)`$, which also commutes with the global symmetry generator, but has a nontrivial commutator with $`U_C`$ if $`C`$ encloses only one of the points $`x`$ or $`y`$. To summarize this section, we have shown that pure Yang Mills theory in 2+1 and 3+1 dimensions has a global $`Z_N`$ magnetic symmetry. The generator of the symmetry group in both cases is the trace of the fundamental Wilson loop along the spatial boundary of the system. The order parameter for this symmetry in 2+1 dimensions is a local scalar field $`V(x)`$, while in 3+1 dimensions a stringlike field $`V(C)`$. In both cases the field $`V`$ is gauge invariant on physical states and is a bona fide canonical order parameter which distinguish in gauge invariant way the phases of the theory. In the next section we discuss the realization of the magnetic symmetry in the confined and the deconfined phases. ## 3 Hot and cold realizations of the magnetic $`Z_N`$. As with any global symmetry, it is important to understand what is the mode of realization of magnetic $`Z_N`$ in the ground state of the theory. This mode of realization depends of course on the parameters of the theory as well as on the temperature. The situation at zero temperature is well understood. ### 3.1 2+1 dimensions. Again we start with three dimensions. There is a very general argument due to ’t Hooft<sup>6</sup><sup>6</sup>6The original argument as stated in is formulated for 3+1 dimensional theories, however its generalization to 2+1 dimensions requires only linguistic changes. stating that if the theory does not have zero mass excitations the area law of the Wilson loop implies the nonvanishing expectation value of the vortex operator $`V(x)`$. Conversely if the Wilson loop has a perimeter law the expectation value of $`V(x)`$ must vanish and the correlation function $`V(x)V^{}(y)`$ must have an exponential falloff with $`\|xy\|`$. It follows that in the pure Yang Mills theory the vacuum expectation value of the vortex operator does not vanish and therefore the $`Z_N`$ magnetic symmetry is spontaneously broken. The same is true in the partially broken Higgs phase of the Georgi-Glashow model. As mentioned in the last section the confining and the Higgs regimes in this model are analytically connected and therefore the realization of all global symmetries in the two regimes is the same. In fact in the weakly coupled Higgs phase this can be verified by the direct calculation of the expectation value of $`V`$ . This calculation maps very simply into the classic monopole plasma calculation of Polyakov and was discussed in detail in . One can also explicitly construct the low energy effective Lagrangian in terms of the field $`V`$ which realizes the spontaneously broken $`Z_N`$ symmetry and describes the low energy spectrum of the Georgi - Glashow vacuum. $$=_\mu V^{}^\mu V\lambda (V^{}V\mu ^2)^2\zeta (V^2+V^2)$$ (36) A similar effective Lagrangian with some quantitative differences was argued to be valid also for the pure Yang-Mills theory in . The application of the ’t Hooft argument at finite temperature is somewhat less straightforward. Since at finite temperature the Lorentz invariance is broken, the temporal and spatial Wilson loops do not necessarily have the same behaviour and one has to be more careful. The original argument relates the behaviour of the vortex operator and the temporal Wilson loop. At finite temperature in the Euclidean formalism the extent of the system in the temporal direction is finite. As a result it is not possible to distinguish between the area and perimeter law for ”asymptotically” large temporal loops. Instead the role of the temporal Wilson loop is taken over by the Polyakov line - the loop that winds around the total volume of the system in the temporal direction. Thus one expects that in the deconfining phase where the Polyakov line has a nonvanishing vacuum average, the vortex operator should have vanishing expectation value. Indeed this can be easily confirmed by the explicit calculation of the VEV of the vortex operator using the method of . In the calculation was performed in 3+1 dimensions, but adapting it to 2+1 dimensional case is trivial. We give below a brief outline. Consider the equal time vortex-antivortex correlation function. At finite temperature is it given by the following expression $$<V(x)V^{}(y)>=\mathrm{Tr}e^{\frac{\beta }{2}(E^2+B^2)}e^{i\frac{2\pi i}{g}_x^yϵ_{ij}𝑑l^iE_3^j}$$ (37) The line integral can be taken along the straight line $`L`$ connecting the points $`x`$ and $`y`$. For definiteness we take $`x`$ and $`y`$ to be separated in the direction of the first axis. Introducing the imaginary time axis and the Lagrange multiplier field $`A_0`$ in the standard way this expression can be transformed to $$<V(x)V^{}(y)>=DA_iDA_0\mathrm{exp}\{\frac{1}{2}_0^\beta 𝑑td^3x\left(_0A_i^a(D_iA_0)^a\delta (t)a_i^a\right)^2+(B^a)^2\}$$ (38) where the “external field “ $`a^i`$ is given by $$a_i^a(𝐱)=\delta ^{a3}\delta _{i2}\frac{2\pi }{g}\delta (𝐱L)$$ (39) The delta function in time in front of the external field $`a^i`$ in eq.(38) appears for the following reason. As we saw in section 2.3 the product of the vortex and an antivortex operator is gauge invariant. This is because it induces a singular gauge transformation which is continuous up to the center element. However if we split it up in imaginary time into infinitesimal bits $`V_{dt}=e^{idtT\frac{2\pi i}{g}_x^yϵ_{ij}𝑑l^iE_3^j}`$ then any single such bit separately is not gauge invariant, since the transformation it induces is genuinely discontinuous across the line connecting the points $`x`$ and $`y`$. The operators $`V_{dt}`$ therefore do not commute with the projection operator on physical states. To obtain the path integral representation for the expectation value eq.(37) we should introduce the projection operator only at the last point in imaginary time and not at the intermediate points. In the path integral language this corresponds to the gauge fixing $`A_0=0`$ everywhere except at one time, say $`t=0`$. In this gauge it is straightforward to see that the Gaussian integration over the electric field leads to the usual path integral with $`_iA_0`$ shifted by $`a_i`$. This can then be rotated to an arbitrary gauge with the result eq.(38)<sup>7</sup><sup>7</sup>7In this derivation we dropped commutator terms between the Hamiltonian and the exponent in the vortex-antivortex operator. These commutator terms only exist at $`t=0`$ and therefore drop out in the continuum limit (i.e are $`O(dt)`$). In the lattice realization they are indeed present and complete the expression eq.(38) to the twisted plaquette representation. To evaluate the path integral eq.(38) we follow the standard procedure and integrate out all modes except for the Polyakov loop in a saddle point approximation,. This leads to the effective action $`S_{eff}(q,a_i)`$, where $`1/2TrP=\mathrm{cos}q`$ $$<V(x)V^{}(y)>=Dq\mathrm{exp}S_{eff}(q,a_i)$$ (40) To one loop order the effective action is given by $$S_{eff}=d^2x\left(\frac{2T}{g^2}(_iq+\frac{g}{2}a_i)^2+U(q)\right)$$ (41) The matrices $`\tau ^a`$ in $`A_\mu `$ are the generators of $`SU(2)`$ in the fundamental representation and are normalized according to $`\mathrm{tr}\tau ^a\tau ^b=\frac{1}{2}\delta ^{ab}`$. The one loop effective potential $`U`$ is related to a Bernoulli polynomial and can be read off the expressions in . The only property of $`U`$ which is important to us is that it has two degenerate minima at $`q=0,\pi `$. To calculate the correlator we have to find the configuration of $`q`$ which minimizes the action eq.(41). Qualitatively the form of the solution is clear. The considerations identical to those in tell us that it must be the “broken” electric $`Z_2`$ domain wall: half a wall ( $`q_{x_2\mathrm{}}0`$, $`q(x_2=0)=\frac{\pi }{2}`$) above the line $`L`$ and half a wall ( $`q(x_2=0)=\frac{\pi }{2}`$, $`q_{x_2\mathrm{}}0`$,) below the line $`L`$ separated by a discontinuity $`\delta q=\pi `$. The action of such a configuration is $`S_{eff}=\stackrel{~}{\alpha }\|xy\|`$ where $`\stackrel{~}{\alpha }`$ is the “$`Z_2`$ domain wall tension”. The vortex correlator is thus given by $$<V(x)V^{}(y)>=\mathrm{exp}\{\stackrel{~}{\alpha }\|xy\|\}$$ (42) As $`\|xy\|`$ become large the correlation function decreases exponentially, and thus the expectation value of the vortex operator vanishes. For the $`SU(N)`$ group this calculation trivially generalizes and gives the same result. The exponential decay is also obtained for the correlator of $`V^m`$ with any power $`m<N`$. Recall that the vortex operator is the order parameter for the magnetic $`Z_N`$ symmetry. Moreover the powers of $`V`$ exhaust all possible local order parameters<sup>8</sup><sup>8</sup>8The latter statement is correct modulo multiplication of $`V^m`$ by local gauge invariant and $`Z_N`$ invariant operators. These possible factors do not change the fact of the exponential decay of the corelators and are therefore unimportant for our discussion.. Their vanishing is therefore an unambiguous indication that the magnetic $`Z_N`$ is restored in the high temperature deconfined phase. In hindsight this is not very surprising. Indeed, we are dealing with physical discrete symmetry which is spontaneously broken at zero temperature. When the system is heated it is unavoidable that entropy effects take over and at some sufficiently high temperature the symmetry must be restored. A good qualitative guide here is the effective Lagrangian eq.(36). It describes a simple $`Z_2`$ invariant scalar theory. There is very little doubt that a system described by this Lagrangian indeed undergoes a symmetry restoring phase transition at some $`T_c`$. Moreover the effective Lagrangian approach also suggests that this phase transition has deconfining character. As shown in the charged states in the effective theory eq.(36) are represented by solitonic configurations of the vortex field $`V`$ with unit winding number. The energy of any such state is linearly divergent in the infrared. The reason is that due to finite degeneracy of vacuum states, the minimum energy configuration looks like a quasi onedimensional strip across which the phase of $`V`$ winds. The energy density inside this ”electric flux tube” is proportional to the vacuum expectation value of $`V`$. When the VEV vanishes, so does the string tension. Stated in other words, when $`V`$ vanishes, the phase fluctuations are large and the winding number is not a sharp observable. An external charge is thus screened easily by regions of space around it with vanishing $`V`$ . The phase with $`<V>=0`$ is therefore not confining. In the theory with several Higgs fields this phase exists even at zero temperature and corresponds to a completely Higgsed phase - where the gauge group is broken completely. In such a Higgs phase indeed the colour is screened rather than confined. In the pure Yang-Mills theory this phase is absent at zero temperature, but is realized as the deconfined phase at $`T>T_c`$. We thus see that the behaviour of the vortex operator at high temperature does indeed match the simple intuition coming from a $`Z_N`$ invariant effective Lagrangian very well. ### 3.2 Extension to 3+1 dimensions The ’t Hooft argument now states that the vanishing vacuum average of the Polyakov line is incompatible with the area law behaviour of the spatial ’t Hooft loop and vice versa. This means that in the confining phase the ’t Hooft loop has perimeter law. In the high temperature deconfined phase the behaviour of the spatial ’t Hooft loop must become area since the average of the Polyakov line is finite. Again this is confirmed by explicit calculation in . A more subtle question is how the behaviour of the ’t Hooft loop relates to the realization of the magnetic $`Z_N`$ symmetry. The $`Z_N`$ symmetry does not have an order parameter which is a local field defined at a point. The only order parameters in the strict sense (an eigenoperator with a nonvanishing eigenvalue) is a ’t Hooft line $`V(C)`$ which runs through the whole system . In a system which is finite in the direction of the loop, but is infinite in the perpendicular directions everything is clear cut. In this case there are two possibilities: a) $`<V>0`$ and the magnetic $`Z_N`$ broken, or b). $`<V>=0`$ and the nagnetic $`Z_N`$ restored. In the system infinite in all directions $`C`$ is necessarily an infinite line, and the expectation value $`<V(C)>`$ clearly vanishes irrespective of whether $`Z_N`$ is broken or not. The ’t Hooft loop along a closed contour on the other hand is never zero, since it is globally invariant under the $`Z_N`$ transformation. It is therfore impossible to find an operator whose VEV distinguishes between the two possible realizations of the magnetic symmetry by vanishing in one phase and not vanishing in the other. Nevertheless the behaviour of the closed loop does indeed reflect the realization mode of the symmetry, since it is qualitatively different in the two possible phases. Namely vacuum expectation value of a large closed ’t Hooft loop (by large, as usual we mean that the linear dimensions of the loop are much larger than the correlation length in the theory) has an area law decay if the magnetic symmetry is spontaneously broken, and perimeter law decay if the vacuum state is invariant. To understand the physics of this behaviour it is useful to think of the ’t Hooft line as built of “local” operators - little “magnetic dipoles”. Consider eq.(34) with the contour $`C`$ running along the x- axis and the surface $`S`$ chosen as the $`(x,y)`$ plane. Let us mentally divide the line into (short) segments of length $`2\mathrm{\Delta }`$ centered at $`x_i`$. Each one of these segments is a little magnetic dipole and the ’t Hooft loop is a product of the operators that create these dipoles. The definition of these little dipole operators is somewhat ambiguous but since we only intend to use them here for the purpose of an intuitive argument any reasonable definition will do. It is convenient to define a single dipole operator in the following way $$D_\mathrm{\Delta }(x)=\mathrm{exp}\{id^3y[a_i^+(x+\mathrm{\Delta }y)+a_i^{}(x\mathrm{\Delta }y)]\mathrm{Tr}(\mathrm{YE}^\mathrm{i}(\mathrm{y}))\}$$ (43) where $`a_i^\pm (xy)`$ is the c-number vector potential of the abelian magnetic monopole (antimonopole) of strength $`4\pi /gN`$. The monopole field correpsonding to $`a_i`$ contains both, the smooth $`x_i/x^3`$ part as well as the Dirac string contribution. The Dirac string of the monopole - antimonopole pair in eq.(43) is chosen so that is connects the points $`x\mathrm{\Delta }`$ and $`x+\mathrm{\Delta }`$ along the straight line. The dipole operators obviously have the property $$D_\mathrm{\Delta }(x)D_\mathrm{\Delta }(x+2\mathrm{\Delta })=D_{2\mathrm{\Delta }}(x+\mathrm{\Delta })$$ (44) This is because in the product the smooth field contribution of the monopole in $`D_\mathrm{\Delta }(x)`$ cancels the antimonopole contribution in $`D_\mathrm{\Delta }(x+2\mathrm{\Delta })`$, while the Dirac string now stretches between the points $`(x\mathrm{\Delta })`$ and $`(x+3\mathrm{\Delta })`$. When multiplied over the closed contour, the smooth fields cancel out completely, while the surviving Dirac string is precisely the magnetic vortex created by a closed ’t Hooft loop operator. The ’t Hooft loop can therefore be written as $$V(C)=\mathrm{\Pi }_{x_i}D_\mathrm{\Delta }(x_i)$$ (45) The dipole operator $`D(x_i)`$ is an eigenoperator of the magnetic flux defined on a surface that crosses the segment $`[x_i\mathrm{\Delta },x_i+\mathrm{\Delta }]`$. Suppose the magnetic symmetry is broken. Then we expect the dipole operator to have a nonvanishing expectation value<sup>9</sup><sup>9</sup>9The magnetic dipole operators defined above are strictly speaking not local, since they carry the long range magnetic field of a dipole. However, the dipole field falls off with distance very fast. Therefore even though this fall off is not exponential the slight nonlocality of $`D`$ should not affect the following qualitative discussion. $`<D>=d(\mathrm{\Delta })`$. If there are no massless excitations in the theory, the operators $`D(x_i)`$ and $`D(x_j)`$ should be decorrelated if the distance $`x_ix_j`$ is greater than the correlation length $`l`$. Due to eq.(45), the expectation value of the ’t Hooft loop should therefore roughly behave as $$<V(C)>=d(l)^{L/l}=\mathrm{exp}\{\mathrm{ln}\left(\frac{1}{d(l)}\right)\frac{L}{l}\}$$ (46) where $`L`$ is the perimeter of the loop. In the system of finite length $`L_x`$, the vacuum expectation value of the vortex line which winds around the system in $`x`$-direction is therefore finite as in eq.(46) with $`LL_x`$. On the other hand in the unbroken phase the VEV of the dipole operator depends on the size of the system in the perpendicular plane $`L_y`$. For large $`L_y`$ it must vanish exponentially as $`d=\mathrm{exp}\{aL_y\}`$. So the expectation value of $`V`$ behaves at finite $`L_y`$ in the unbroken phase as: $$<V(C)>=\mathrm{exp}\{aL_yL_x\}$$ (47) and vanishes as $`L_y\mathrm{}`$. Thus in a system which is finite in $`x`$ direction, but infinite in $`y`$ direction, the ’t Hooft line in the $`x`$ direction has a finite VEV in the broken phase and vanishing VEV in the unbroken phase. In the limit of the infinite system size $`L_x\mathrm{}`$ the VEV obviously vanishes in both phases. This is of course due to the fact that $`V`$ is a product of infinite number of dipole operators, and this product vanishes even if individual dipole operators have finite VEV<sup>10</sup><sup>10</sup>10The VEV of the dipole $`D`$ must be smaller than one since $`D`$ is defined as a unitary operator.. However one can avoid any reference to finite size system and infinite vortex lines by considering closed ’t Hooft loops. For a closed loop with long sides along $`x`$ axis at $`y=0`$ and $`y=R`$ the above argument leads to the conclusion that in the broken phase $`V`$ must have a perimeter law, eq.(46). In the unbroken phase the correlation between the dipoles at $`y=0`$ and dipoles at $`y=R`$ should decay exponentially $`<D(0)D(R)>\mathrm{exp}\{\alpha \frac{R}{l}\}`$ and thus $$<V(C)>=\mathrm{exp}\{\alpha \frac{LR}{l^2}\}=\mathrm{exp}\{\alpha \frac{S}{l^2}\}$$ (48) Thus the perimeter behaviour of $`<V(C)>`$ indicates a vacuum state which breaks spontaneously the magnetic $`Z_N`$ while the area behaviour means that the magnetic $`Z_N`$ is unbroken. The results of then mean that in 3+1 dimensions as well as in 2+1 dimension the magnetic symmetry is restored above the deconfining phase transition, in the sense of eq.( 47). In the next section we discuss what is the implication of this conclusion on the behaviour of the spatial Wilson loop. ## 4 Spatial Wilson loop at low and high temperature. As we have shown in Section 2 the spatial Wilson loop is the generator of the magnetic $`Z_N`$ symmetry. We expect therefore that the mode of realization of the magnetic $`Z_N`$ is strongly linked to the behaviour of $`W`$. The argument is simplest to state for a toy model which exemplifies the basic physics in a very simple setting. Rather than talk about nonabelian gauge theory, consider a scalar theory of a complex field $`\varphi `$ with global $`Z_N`$ theory in 2+1 dimensions. $$=_\mu \varphi _\mu \varphi ^{}+\lambda (\varphi ^{}\varphi \mu ^2)^2+\zeta \left(\varphi ^N+(\varphi ^{})^N\right)$$ (49) The generator of the $`Z_N`$ symmetry is given by $$U=\mathrm{exp}\left\{i\frac{2\pi }{N}d^2xj_0(x)\right\}=\mathrm{exp}\left\{\frac{2\pi }{N}d^2x(\pi \varphi \pi ^{}\varphi ^{})\right\}$$ (50) where $`\pi =_0\varphi ^{}`$ is the momentum conjugate to the field $`\varphi `$. Obviously with the canonical commutation relations between $`\pi `$ and $`\varphi `$ one has $$U\varphi (x)U^{}=e^{i\frac{2\pi }{N}}\varphi (x)$$ (51) We will be interested in the behaviour of the operator which generates the $`Z_N`$ transformation only inside some region $`S`$ of the two dimensional plane. $`U(S)`$ $`=`$ $`\mathrm{exp}\left\{{\displaystyle \frac{2\pi }{N}}{\displaystyle _S}d^2x(\pi \varphi \pi ^{}\varphi ^{})\right\}`$ (52) $`U(S)\varphi (x)U^{}(S)`$ $`=`$ $`e^{i\frac{2\pi }{N}}\varphi (x)xS`$ $`=`$ $`\varphi (x),xS`$ We will refer to this operator as the U-loop. Throughout this discussion we assume that there are no massless excitations in the spectrum of the theory and that the linear dimensions of the area $`S`$ are much larger than the correlation length. The statement we are aiming at is that at zero temperature in the phase with broken $`Z_N`$ the U-loop has an area law behaviour while in the phase with unbroken $`Z_N`$ this changes into the perimeter law behaviour. ### 4.1 U-loop in the broken phase Consider the broken phase first. We are interested in the vacuum expectation value of $`U(S)`$. This is nothing but the overlap of the vacuum state $`<0\|`$ and the state which is obtained by acting with $`U(S)`$ on the vacuum state $`\|S>=U\|0>`$. If the symmetry is broken, the field $`\varphi `$ in the vacuum state is pointing in some fixed direction in the internal space. In the state $`\|S>`$ on the other hand its direction in the internal space is different - rotated by $`2\pi /N`$ \- at points inside the area $`S`$. In the local theory with finite correlation length the overlap between the two states approximatelly factorises into the product of the overlaps taken over the region of space of linear dimension of order of the correlation length $`l`$ $$<0\|S>=\mathrm{\Pi }_x<0_x\|S_x>$$ (53) where the label $`x`$ is the coordinate of the point in the center of a given small region of space. For $`x`$ outside the area $`S`$ the two states $`\|0_x>`$ and $`\|S_x>`$ are identical and therefore the overlap is unity. However for $`x`$ inside $`S`$ the states are different and the overlap is therefore some number $`e^\gamma `$ smaller than unity. The number of such regions inside the area is obviously of order $`S/l^2`$ and we thus $$<U(S)>=\mathrm{exp}\{\gamma \frac{S}{l^2}\}$$ (54) In a weakly coupled theory this argument is confirmed by explicit calculation. The expectation value of the U-loop in the theory eq.(49) is given by the following path integral $$<U(C)>=𝑑\varphi 𝑑\varphi ^{}\mathrm{exp}\left\{(_\mu \varphi +i\varphi \chi _\mu )(_\mu \varphi ^{}i\varphi ^{}\chi _\mu )+\lambda (\varphi ^{}\varphi \mu ^2)^2+\zeta \left(\varphi ^N+(\varphi ^{})^N\right)\right\}$$ (55) with $`\chi _\mu (x)`$ $`=`$ $`{\displaystyle \frac{2\pi }{N}}\delta ^{\mu 0}\delta (x_0),xS`$ (56) $`=`$ $`0,xS`$ (57) This expression directly follows from eq.(52) and integration over the canonical momentum in the phase space path integral. At weak coupling this path integral is dominated by a simple classical configuration. First, it is clear that the solution must be such that the phase of the field $`\varphi `$ has a discontinuity of $`2\pi /N`$ when crossing the surface $`S`$ since otherwise the action is UV divergent due to singular $`\chi `$. Asymtotically at large distance from the surface the field should approach its vacuum expectation value. Since the source term $`\chi `$ vanishes outside $`S`$, eveywhere where $`\varphi `$ is continuous it has to solve classical equations of motion. Also, for values of $`x_1`$ and $`x_2`$ which are well inside $`S`$ the profile $`\varphi `$ should not depend on these coordinates, but should only depend on $`x_0`$. It is easy to see that a solution with these properties exists: it is given by the “broken” domain wall solution. Recall that the vacuum is degenerate and so there certainly exists a classical solution of the equations of motion which interpolates between two adjacent vacuum states $`\varphi _{x_0\mathrm{}}\varphi _0`$ and $`\varphi _{x_0\mathrm{}}e^{i\frac{2\pi }{N}}\varphi _0`$. Breaking this classical solution along the plane $`x_0=0`$ and rotating the piece $`x_0<0`$ by $`2\pi /N`$ produces precisely the configuration with the correct boundary conditions and the discontinuity structure. The path integral in eq.(55) is therefore dominated by this classical configuration. Its action (up to corrections associated with the boundary effects of $`S`$) is $`\alpha S`$ where $`\alpha `$ is the classical wall tension of the domain wall which separates two adjacent $`Z_N`$ vacua. Thus we find that the expectation value of the U-loop is related to the domain wall tension of the $`Z_N`$ domain wall by $$<U(S)>=\mathrm{exp}\{\alpha S\}$$ (58) ### 4.2 U-loop in the unbroken phase Now consider the unbroken phase. Again the U-loop average has the form of the overlap of two states which factorizes as in eq.(53). Now however all observables noninvariant under $`Z_N`$ vanish in the vacuum. The action of the symmetry generator does not affect the state $`\|0>`$. The state $`\|S>`$ is therefore locally exactly the same as the state $`\|0>`$ except along the boundary of the area $`S`$. Therefore the only regions of space which contribute to the overlap are those which lay within one correlation length from the boundary. Thus $$<U(S)>=\mathrm{exp}\{\gamma P(S)\}$$ (59) where $`P(S)`$ is the perimeter of the boundary of $`S`$. The absence of the area law is again easily verified by a perturbative calculation. In the unbroken phase the fluctuations of the field $`\varphi `$ as well as the current density $`j_0=i(\pi \varphi \pi ^{}\varphi ^{})`$ are small. To leading order in the coupling constant $$<U(S)>=\mathrm{exp}\left\{\frac{1}{2}_{x,yS}d^2xd^2y<j_0(x)j_0(y)>\right\}$$ (60) The possible area law contribution in the exponent is $$Sd^2x<j_0(0)j_0(x)>=S\underset{p0}{lim}G(p)$$ (61) where $`G(p)`$ is the Fourier transform of the charge density correlation function. The correlator of the charge densities however vanishes at zero momentum. This is because in the leading perturbative order the symmetry of the theory is actually $`U(1)`$ and not just $`Z_N`$ as seen in eq. 55. Since the vacuum state is invariant it follows that the total charge $`Q=d^2xj_0(x)=j_0(p=0)`$ on this state vanishes, and so does any correlation function that involves zero momentum component of the charge density. So the area contribution in eq.(61) is zero. Strictly speaking in the leading order in perturbation theory eq.(60) is not the complete result. The exact expression contains in the exponential also higher point correlators of the current density. Again however the possible area law contribution contains correlators of the total charge $`Q`$ with powers of $`j_0`$ and therefore vanishes. ### 4.3 U-loop at high temperature Let us see now how the argument changes at high temperature. The important difference is that the vacuum is not a pure state but rather a statistical ensemble. The average of the U-loop is therefore not given by a single matrix element but rather by $$<U>=\underset{i}{}e^{E_iT}<i\|U\|i>$$ (62) Let us consider the theory in which the $`Z_N`$ symmetry is broken at zero temperature. For concreteness we will think about $`Z_2`$ symmetric theory, although qualitatively the discussion does not change for any $`N`$. The two degenerate vacuum states are characterized by the value of the condensate $`<\varphi >=\pm \mu `$. In order to understand the behaviour of the $`U`$-loop we have to figure out what types of states contribute to the thermal ensemble. At zero temperature the only states that are of interst are those with finite energy. There are two towers of such states $`\|n>_\mu `$ and $`\|n>_\mu `$ \- constructed above each one of the degenerate vacua. These two towers of states do not talk to each other, not only because their overlap is zero, but also because they can not be connected to each other by action of any local (or semilocal) operator $`{}_{\mu }{}^{}<n\|O\|n^{}>_\mu =0`$. An immediate corollary of this is that a superposition of the type $`\|\alpha ,\beta >=\alpha \|n_\mu >+\beta \|n>_\mu `$ violates clustering property of the correlators of local operators. For this reason at zero temperature in a spontaneously broken theory we are never interested in states which carry sharp quantum numbers of the broken symmetry. At finite temperature however we are also asking after states with finite energy density, and therefore infinite energy. This part of the spectrum looks rather different if the energy density involved is high enough. The two vacuum configurations of the potential in eq.(49) are separated by a finite barrier. Let us call the hight of this barrier $`H`$. The states with energy density lower than $`H`$ still separate into two towers. We will denote these states by $`\|l>`$. However higher energy density states, with $`ϵ>H`$ have different nature. Their wave function is not localized in the field space to the vicinity of one of the vacua, but rather is spread over distances larger than the distance between the two vacuum states $`2\mu `$. These states therefore naturally carry sharp quantum numbers with respect to the broken $`Z_2`$ symmetry. These states we will denote by $`\|h>`$. In fact one expects that the higher the energy density the more these states look like the multiparticle states of a symmetric phase. That is to say as long as $`ϵ>>H`$ it does not matter whether the potential has the double well structure or a single vacuum. These highly excited states should look like states with finite density of particles which carry the $`Z_2`$ charge. At low temperatures, when the entropy effects are not important the contribution to the thermal ensemble comes only from the $`\|l>`$ \- sector since the Bolzman factor for any of the $`\|h>`$ states vanishes exponentially in the infinite volume limit. As we have argued earlier, the average of $`U`$ in each one of these states has an area law behaviour and so does the whole temperature average of $`U`$. When the temperature reaches $`T_C`$ the phase transition occurs. The reason for the onset of the phase transition is that when the equilibrium energy density reaches critical threshold value, the $`\|h>`$ sector states start contributing to the thermal ensemble. The sudden change in the entropy due to these new channels drives the phase transition. Above the phase transition therefore there are two kinds of states that contribute to thermal averages. One can then write $$<U>_{T>T_c}=\underset{n}{}e^{E_n^l/T}<n,l\|U\|n,l>+\underset{s}{}e^{E_n^h/T}<n,h\|U\|n,h>$$ (63) where $`n`$ stands for all other quantum numbers. In fact once the entropy effects become important enough to excite the $`\|h>`$-sector, the contribution of $`\|l>`$ states to any physical observable becomes negligible. As discussed earlier each state in the second term gives a perimeter contribution $`\mathrm{exp}\{\gamma _sP\}`$ to the average of the $`U`$-loop, so one could be tempted to conclude that the loop must have a perimeter law just like in zero temperature vacuum of a symmetric phase. This however is not necessarily the case. The reason is that $`<n,h\|U\|n,h>`$ is not positive definite. In fact the number of states in which it is positive is roughly equal to the number of states on which it is negative. It is therefore very likely that the leading perimeter behaviour will cancel and the net result will again be an area law for $`<U>`$. Indeed if the ensemble can be thought of as an ensemble of $`Z_2`$ charged free particles, the area law for $`U`$-loop follows immediately<sup>11</sup><sup>11</sup>11This argument is borrowed from .We thank Mike Teper and Biaggio Lucini for discussions of this point.. The $`U`$-loop in such an ensemble is $$<U(S)>=\underset{x_i,n}{}\frac{1}{n!}\mu ^n(1)^n$$ (64) where $`\mu `$ is the fugacity of a single particle, and the summation goes over all coordinates $`x_i`$ of the particles inside the area $`S`$ and over all possible numbers of particles $`n`$. Assuming that particles have a finite size $`\mathrm{\Delta }`$, so that there are $`S/\mathrm{\Delta }`$ possibilities to place one particle inside the area $`S`$ in the dilute gas approximation the sum gives $$<U(S)>=\mathrm{exp}\{\frac{\mu }{\mathrm{\Delta }}S\}$$ (65) We stress that the thermal ensemble of particles is $`Z_2`$ invariant. The density matrix for such an ensemble can be written in the particle basis as $$\rho =\underset{n(x)}{}\frac{1}{n!}\mu ^n\|n><n\|$$ (66) The operator of the $`Z_2`$ transformation acts on the $`n`$-particle states as $$U\|n>=(1)^n\|n>$$ (67) and so $$U\rho U^{}=\rho $$ (68) The explicit simple formula eq.(65) is derived in the dilute gas approximation. We expect that the physics will be similar as long as the interaction between the particles is short range. Whenever the interaction is long range the behaviour of $`<U(S)>`$ can be different. For instance one does not expect area law behaviour if the particles are bound into pairs since in this case only $`Z_2`$ invariant states contribute to the thermal ensemble. Our conclusion is that at finite temperature the behaviour of the $`U`$-loop is not strongly related to the mode of the realization of the $`Z_N`$ symmetry. It is rather more likely to have an area behaviour. To reiterate, the physics involved is very simple. At zero temperature when acting on a state, the $`U`$-loop performs the $`Z_N`$ transformation inside the loop. The only degrees of freedom that are changed by this operation inside the loop, are the $`Z_N`$ \- noninvariant fields. If the vacuum wavefunction depends on the configuration of the noninvariant degrees of freedom (the state in question is not $`Z_N`$ invariant) the action of $`U`$-loop affects the state everywhere inside the loop. The VEV of $`U`$-loop then falls off as an area. If the vacuum is $`Z_N`$ invariant, the wavefunction does not depend on the configuration of the noninvariant degrees of freedom. The action of $`U`$-loop then perturbs the state only along the perimeter, hence the perimeter law in the unbroken phase. At finite temperature however the thermal ensemble may even in the symmetric phase contain significant contributions from states with nonvanishing $`Z_N`$ charges. The $`U`$-loop therefore perturbs the thermal ensemble very significantly everywhere inside the area, and the natural outcome is an area law. The argument is quite general and does not depend on the exact form of the $`Z_N`$ invariant potential and more generally on the field content of the theory - we could have added any number of extra fields to the theory eq.(49) without changing the conclusions. The same relation must exist between the mode of realization of the magnetic $`Z_N`$ symmetry and the behaviour of the Wilson loop in the pure Yang-Mills theory. The direct analogs of the scalar field $`\varphi (x)`$ in eq.(49) and the U-loop of the scalar theory are correspondingly the vortex field $`V(x)`$, and the spatial Wilson loop $`W(C)`$. As we have shown in the previous section, the magnetic $`Z_N`$ is restored at high temperature. The Wilson loop is nevertheless likely to have an area law as is indeed indicated by all existing lattice data. In this context we note that analytic strong coupling results also give area behaviour. ### 4.4 Wilson loop in 3+1 dimensions The previous considerations generalize to the 3+1 dimensons. At zero temperature in the broken phase when acting with the Wilson loop $`W(C)`$ on the vacuum one changes the state of those magnetic vortices which loop through $`C`$. The number of such vortices which are present in a generic configuration in the broken phase is proportional to the minimal area subtending $`C`$. The number of the degrees of freedom that is changed by the action of $`W`$ is thus proportional to the area $`S`$. Each of these degrees of freedom contributes a factor smaller than unity to the overlap with the vacuum state and so the VEV of $`W`$ scales with the exponential of the area. In the unbroken phase the vacuum does not contain vortices of arbitrarily large size. The size of the vortices present in the vacuum is cutoff by the relevant correlation length. This is the case if the gauge group is completely broken by the Higgs mechanism. Therefore for contours $`C`$ of linear dimension much larger than this length, the action of $`W(C)`$ only disturbs degrees of freedom close to the contour $`C`$ itself and the VEV must have the perimeter behaviour. At high temperature even the symmetric thermal ensemble is populated by vortices. These vortices are not free since apart from them the ensemble also contains “free” charges. However unless this background of charges induces long range interactions between the vortices, the most probable result for the Wilson loop is the area law. A more detailed knowledge of vortex dynamics is necessary to draw a firm conclusion. To close this section we note that the present considerations do not apply to Abelian theories. The magnetic symmetry does exist in this case too, but here it is the continuos $`U(1)`$ group and the spectrum is massless. In this case there is no reason to expect the local factorization of the overlap and generically therefore the arguments of this section do not hold. In particular in the presence of long range correlations it is perfectly possible that the Wilson loop has a perimeter law even though the state is perturbed everywhere inside the area boundede by the loop<sup>12</sup><sup>12</sup>12In 2+1 dimensions it is actualy only the noncompact Abelian theories that are excluded from the consideration. Compact theories are massive and therefore should behave in the same way as the nonabelian Yang-Mills.. ## 5 Discussion. In this paper our aim was to point out two facts. First that the calculation of the VEV of the ’t Hooft loop implies the restoration of the magnetic $`Z_N`$ symmetry above the deconfinement transition. Second, that the mode of realization of magnetic symmetry is closely related to the behaviour of the spatial Wilson loop. At zero temperature this relation is very rigid: spontaneous breaking of $`Z_N`$ implies the area law behaviour for $`W`$ while unbroken $`Z_N`$ leads to perimeter law behaviour. At high temperature however even though the $`Z_N`$ symmetry is restored, the Wilson loop may have an area law. This is the consequence of the fact that even a $`Z_N`$ \- invariant the thermal ensemble can contain a significant contribution of $`Z_N`$ nonsinglet states. The area law is particularly simply understood if the thermal ensemble at high temperature is well appproximated by an ensemble of weakly interacting magnetic vortices. The vortex gas argument has been previously brought up in favour of the area behaviour of the spatial Wilson loop in . This behaviour is also confirmed by several lattice gauge theory calculations . An alternative possibility is that due to the as yet unknown vortex dynamics, there is vortex - antivortex binding. To explore such a possibility it would be very interesting to measure on the lattice the free energy of a magnetic vortex. In ref.() the behaviour of the free energy of magnetic and electric fluxes has been discussed in the low temperature phase. To be able to do it in the lattice framework one has to define the theory in a finite volume. As discussed by ’t Hooft this can be achieved by imposing on the potentials periodic boundary conditions modulo a gauge transformation. As discussed in ref. this admits the presence of vortices in 2+1 and of the vortex lines in 3+1 dimensions. ’t Hooft’s discussion was based on a Euclidean rotation identity for the twisted 4d path integrals valid for any temperature, and a factorization property of magnetic and electric fluxes. In the notation of ref.: $$F(\stackrel{}{e},\stackrel{}{m})=F_e(\stackrel{}{e})+F_m(\stackrel{}{m})$$ (69) Its validity at low T is very reasonable, but is inconsistent with the Euclidean rotation identity at high T. Based on this ’t Hooft could prove ($`N3)`$ that in the confining phase, where the free energy of an electric flux is linear with the length (with the string tension $`\rho `$), the free energy of magnetic flux vanishes exponentially in the infinite volume limit. For a magnetic flux in, say the z-direction it is $`\mathrm{exp}\rho L_xL_y`$. Thus the free energy of a magnetic flux is related to the behaviour of the Wilson loop. The free energy of an electric flux in the z direction in the hot phase vanishes exponentially like $`\mathrm{exp}\alpha L_xL_y`$ where $`\alpha `$ is the surface tension found in ref.. So the next obvious question is how the magnetic flux free energy behaves in the hot phase. In a vortex gas picture this free energy should vanish exponentially in the infinite volume limit. Such a calculation in 2+1 dimensions has been performed and (modulo some uncertainty related to imperfect measurement of global vorticity) results are consistent with this expectation . It would also be instructive to see how in the hot phase the additivity of electric and magnetic fluxes is broken. We note that a recent lattice calculation measures the monopole-antimonopole correlation. The results of point to the screened behaviour of this correlation function for all temperatures. So in the hot phase it behaves like its electric partner, the correlator of Polyakov loops<sup>13</sup><sup>13</sup>13We note however that this simulation also points to the Coulomb behaviour for the spatial ’t Hooft loop in the hot phase, in contradiction to analytic results and early lattice results . We feel that here more work should be done to clarify the situation.. This via ’t Hoofts argument, is consistent with the measured area behaviour of Wilson loops and would imply that the magnetic flux free energy would fall off with an area law for all temperatures. It is interesting to note that the vortex gas picture is equally applicable in high temperature confining and nonconfining gauge theory. In particular one can consider an $`SU(N)`$ gauge theory with sufficient number of adjoint Higgs fields, so that the gauge theory at zero temperature is broken completely. In this situation the magnetic $`Z_N`$ symmetry is unbroken in the vacuum and the Wilson loop has a perimeter law. Magnetic vortices are finite energy excitations with the mass of order $`M=M_v^2/g^2`$, where $`M_v`$ is the vector boson mass. When the system is heated one expects that the thermal ensemble will contain a dilute gas of these vortices at any temperature. Therefore at any finite temperature the spatial string tension should be nonzero, although at low temperatures it will be exponentially suppressed if the theory is weakly coupled: $`\sigma M_v^2\mathrm{exp}\{M_v^2/g^2T\}`$. Acknowledgements The work of A.K. is supported by PPARC. The work of CPKA and AK is supported by a joint CNRS-Royal Society project. We are indebted to Frithjof Karsch for bringing ref. to our attention. We thank Ian Kogan, Biaggio Lucini and Mike Teper for very interesting and useful dicsussions, which in particular helped to clarify a fundamental flaw in the arguments in the first version of this paper.
warning/0004/cond-mat0004102.html	ar5iv	text	# Subdiffusive fluctuations of “pulled” fronts with multiplicative noise ## Abstract We study the propagation of a “pulled” front with multiplicative noise that is created by a local perturbation of an unstable state. Unlike a front propagating into a metastable state, where a separation of time scales for sufficiently large $`t`$ creates a diffusive wandering of the front position about its mean, we predict that for so-called pulled fronts, the fluctuations are subdiffusive with root mean square wandering $`\mathrm{\Delta }(t)t^{1/4}`$, not $`t^{1/2}`$. The subdiffusive behavior is confirmed by numerical simulations: For $`t600`$, these yield an effective exponent slightly larger than $`1/4`$. Since the late 1930s, when the concept of front propagation emerged in the field of population dynamics , interest in this type of problems has been growing steadily in chemistry , physics and mathematics . In physics, the importance of the problem has become more and more clear since it plays a role in a large variety of situations, ranging from reaction-diffusion systems to pattern forming systems in general . Front propagation into unstable states is an interesting dynamical problem by itself. For a front evolving from a local perturbation there are but two possible propagation mechanisms that are determined by the nonlinearities in the equation of motion: Either the nonlinearities determine the velocity of the front that then is called “pushed”; or the nonlinearities simply cause saturation and the velocity is determined by a linearisation about the unstable state. Fronts of this type are called “pulled” because they are “pulled along” by the spreading and growth of small perturbations about the unstable state . Hence pulled front propagation can occur only if the penetrated state is linearly unstable. The pushed and pulled regimes are also known as nonlinear and linear marginal stability . For the discussion below, it is important to realize that pushed fronts relax exponentially in time to their long time asymptotes, but that pulled fronts relax algebraically without characteristic time scale . Hence an adiabatic decoupling of some outer dynamics from the internal relaxation of a pulled front is not possible , and stochastic pulled fronts may show anomalous scaling . Generally, noise can affect the phenomenological description of a reaction-diffusion system in various ways. A first possibility is intrinsic noise modelled typically by additive thermal noise in a Langevin type equation. A second possibility, on which the present paper is focused, is at the external level, e.g. due to fluctuations of some control parameter. An example are the fluctuations of the luminosity intensity in the photosensitive Belousov-Zhabotinsky reaction . Such fluctuations enter the dynamical equation as multiplicative noise. The multiplicative noise of the control parameter usually results in a modification of the mean propagation velocity of the front and in a stochastic wandering of the front position around its mean propagation. This means that the noisy front can be thought of as a coherent structure whose motion can be decomposed into drift plus Brownian motion, very much like a particle sedimenting in a fluid. The drift component corresponds to an average front, with the average taken over the ensemble of all the realizations of the noise. It propagates according to a deterministic equation of motion, whose dynamical parameters are in the simplest case just renormalised by the noise. Theoretically, the important question then arises whether the effects of the fluctuations of the front can be understood in terms of a diffusive or subdiffusive wandering of some suitably defined front position. The renormalisation of the front velocity has been studied in the pushed and pulled regime , while the wandering process is understood only in the pushed case , where it has been shown to be diffusive: the root mean square position of the front $`\mathrm{\Delta }`$ grows with time as $`\sqrt{2D_ft}`$. Actually, the expression for the effective front diffusion coefficient $`D_f`$ derived by Armero et al. was found to break down for pulled fronts, and it was suggested that the wandering of pulled fronts is subdiffusive. In this paper we take up the issue of the stochastic wandering of pulled fronts about their mean position, and predict that in the presence of multiplicative noise pulled fronts behave subdiffusively, with $`\mathrm{\Delta }t^{1/4}`$. This prediction is based on two different arguments. First of all, we heuristically insert the leading edge asymptotics of the relaxing pulled front into the expression for the diffusion coefficient $`D_f`$ of pushed fronts, and immediately find $`\mathrm{\Delta }t^{1/4}`$. Our second argument for the subdiffusive $`\mathrm{\Delta }t^{1/4}`$ behaviour comes from mapping the dynamically important region onto the KPZ equation. We finally also present data of extensive numerical simulations that support our analytical prediction that the wandering is subdiffusive with exponent close to 1/4. The qualitative difference between pushed and pulled fronts results from the fact that the dynamically important region for pushed fronts is the interior front region, whose extent is finite, while that of pulled fronts is the leading edge ahead of the front . Starting from a local initial perturbation, the leading edge region grows without bound and as we shall see, this causes the subdiffusive behaviour. The power law relaxation of deterministic pulled fronts is another manifestation of the leading edge dominated dynamics of pulled fronts . For concreteness, we derive our results by including noise in the one dimensional prototype front equation $`{\displaystyle \frac{\varphi }{t}}=D{\displaystyle \frac{^2\varphi }{x^2}}+f(\varphi ),f(\varphi )=\varphi (1\varphi )(a+\varphi ).`$ (1) Here $`a`$ is a parameter which plays the role of the control parameter. Equation (1) has a stable state $`\varphi =1`$ and a stationary state $`\varphi =0`$ whose relative stability can be tuned by changing the value of the parameter $`a`$. The case $`\frac{1}{2}<a<\frac{1}{2}`$ leads to pushed dynamics, while $`\frac{1}{2}<a<1`$ produces pulled fronts . For the case $`a=1`$, which we will study, the so-called Fisher-Kolmogoroff-Petrovsky-Piscounoff (F-KPP) equation is recovered. Let us assume now that the parameter $`a`$ is replaced by a new fluctuating parameter $`a(x,t)`$ with average $`\overline{a}`$, $`aa(x,t)=\overline{a}+\mu (x,t)`$, where $`\mu (x,t)`$ is a Gaussian noise with the moments: $`\mu (x,t)_\mu `$ $`=`$ $`0,`$ (2) $`\mu (x,t)\mu (x^{},t^{})_\mu `$ $`=`$ $`2\epsilon C(\lambda _\mu \|xx^{}\|)\delta (tt^{}),`$ (3) with $`𝑑xC(\lambda _\mu ,\|x\|)=1`$. We interpret the stochastic $`p.d.e.`$ defined by (1) – (3) in the Stratonovich sense . Notice that if $`1/\lambda _\mu `$ is much smaller than any other length scale in the system, the noise defined by the correlator (3) is effectively white in both time and space. Since according to (1) $`\varphi `$ converges to 1 and is noiseless behind the front, we can suitably define the position $`\mathrm{x}_f(t)`$ of a noisy front propagating to the right into the unstable state $`\varphi =0`$ by $`\mathrm{x}_f(t)={\displaystyle _0^{\mathrm{}}}𝑑x\varphi (x,t),`$ (4) The displacement $`\mathrm{\Delta }\mathrm{x}_f(t)=\mathrm{x}_f(t)\mathrm{x}_f(0)`$ on average grows with the noise renormalized mean velocity $`\overline{v}_R=\dot{\mathrm{x}}_f_\mu `$. The fluctuations about the mean displacement $`\mathrm{\Delta }\mathrm{x}_f(t)_\mu =\overline{v}_Rt`$ are measured by $`\mathrm{\Delta }(t)=\sqrt{\left(\mathrm{\Delta }\mathrm{x}_f(t)\mathrm{\Delta }\mathrm{x}_f(t)_\mu \right)^2_\mu }.`$ (5) If we relate $`\mathrm{\Delta }(t)`$ to a diffusion coefficient $`D_f`$ by writing $`\mathrm{\Delta }^2(t)={\displaystyle _0^t}𝑑t^{}\mathrm{\hspace{0.33em}2}D_f(t^{}).`$ (6) then for pushed fronts the following expression for the diffusion coefficient $`D_f`$ can be derived : $`D_f=\epsilon {\displaystyle \frac{_{\mathrm{}}^{\mathrm{}}𝑑\xi e^{2\overline{v}_R\xi }(d\overline{\varphi }/d\xi )^2g^2(\overline{\varphi })}{\left[_{\mathrm{}}^{\mathrm{}}𝑑\xi e^{\overline{v}_R\xi }(d\overline{\varphi }/d\xi )^2\right]^2}}`$ (7) In this formula, $`\overline{\varphi }`$ is the deterministic field associated with the front moving with the renormalised pushed velocity $`\overline{v}_R`$, $`g(\overline{\varphi })=\frac{f}{a}\|_{\overline{a}}`$ is the derivative of the reaction term with respect to the control parameter, and $`\xi =x\overline{v}_Rt`$ is the comoving coordinate. For pushed fronts, $`D_f`$ given by (7) is finite and time-independent, and hence this gives the diffusive behavior $`\mathrm{\Delta }^2(t)=2D_ft`$. This means that on sufficiently long time scales the random displacement is approximately Markovian, i.e., the sum of uncorrelated and equally distributed random displacements on shorter time scales. As an example of a pulled front with multiplicative noise, we now study the case $`\overline{a}=1`$: $`{\displaystyle \frac{\varphi }{t}}=D{\displaystyle \frac{^2\varphi }{x^2}}+\varphi +\mu \varphi \mu \varphi ^2\varphi ^3.`$ (8) The noise renormalized mean velocity $`\overline{v}_R^{}`$ of the pulled front can be calculated explicitly : $`\overline{v}_R^{}=\dot{\mathrm{x}}(t)_\mu =2\sqrt{D(1+\epsilon C(0))}.`$ (9) However, it is immediately clear that the fluctuation formula (7) cannot naively be extended to the pulled regime. First of all, for a pulled front the expression (7) simply diverges. The divergence of solvability-type expressions actually holds more generally for perturbative expansions about a pulled front . For a pulled front, the dynamically important region is the leading edge defined as the region where linearisation about the unstable state is a valid approximation; the fact that solvability-type integrals like (7) diverge there reflects that the dynamically important region becomes semi-infinite. Second, a pulled front has no characteristic relaxation time , so there is no reason for the Markovian approximation underlying diffusive wandering. Rather the leading edge relaxes asymptotically as $`\varphi `$ $``$ $`\alpha \xi _Re^{\lambda _R^{}\xi _R}e^{\xi _R^2/4Dt}/t^{3/2},\lambda _R^{}=\overline{v}_R^{}/2,`$ (11) $`\text{ for }\xi _R=x\overline{v}_R^{}t1\text{ and }t1.`$ The presence of the $`\alpha \xi _R/t^{3/2}`$ term in front of the exponentials is actually the fingerprint of the full equation being nonlinear. The expression (11) defines a time dependent Gaussian cutoff $`\xi _c\sqrt{4Dt}`$, which regularizes the integrals in (7). In fact, the evaluation of (7) with (11) yields $`D_f(t){\displaystyle \frac{3\epsilon }{(\overline{v}_R^{})^2\sqrt{\pi D}}}{\displaystyle \frac{1}{\sqrt{t}}}(t1).`$ (12) Notice that for large times $`D_f(t)`$ vanishes, marking the nondiffusive wandering of pulled fronts. Insertion into (6) yields $`\mathrm{\Delta }(t)=\sqrt{2{\displaystyle _0^t}𝑑t^{}D_f(t^{})}\left({\displaystyle \frac{12\epsilon }{(\overline{v}_R^{})^2\sqrt{\pi D}}}\right)^{1/2}t^{1/4},`$ (13) so the fluctuations are subdiffusive with exponent $`1/4`$ rather than $`1/2`$. Although the above argument does capture the essential features of fluctuating pulled fronts, it is not entirely systematic, as it is based on the extrapolation of the solvability condition (7) to the pulled regime. In order to substantiate the scaling $`\mathrm{\Delta }(t)t^{1/4}`$ for a relaxing pulled front with a time-dependent analysis, let’s go back to Eq. (8). The leading edge region can be studied by means of the leading edge transformation, $`\varphi (x,t)=\psi (\xi ,t)e^{\lambda ^{}\xi },`$ (14) $`\xi =xv^{}t,v^{}=2,\lambda ^{}=1.`$ (15) Eq. (8) can then be written as $`{\displaystyle \frac{\psi }{t}}`$ $`=`$ $`D{\displaystyle \frac{^2\psi }{\xi ^2}}\psi `$ (16) $`+`$ $`e^\xi \left[(1+\mu )\psi e^\xi \mu \psi ^2e^{2\xi }\psi ^3e^{3\xi }\right].`$ (17) For $`\xi 1`$, the nonlinearities can be neglected $`{\displaystyle \frac{\psi }{t}}=D{\displaystyle \frac{^2\psi }{\xi ^2}}+\mu \psi ,\mathrm{for}\xi 1.`$ (18) Notice that the noise in this “directed polymer” equation still is multiplicative. The Cole-Hopf transformation $`\psi (\xi ,t)=e^{h(\xi ,t)},`$ (19) converts (18) into an equation with additive noise: $`{\displaystyle \frac{h}{t}}=D{\displaystyle \frac{^2h}{\xi ^2}}+D\left({\displaystyle \frac{h}{\xi }}\right)^2+\mu ,\mathrm{for}\xi 1.`$ (20) Eq. (20) is the celebrated 1-dimensional Kardar Parisi Zhang (KPZ) interface equation . The essential difference between our problem and previous studies of the KPZ equation are the initial and boundary conditions. After some temporal evolution, the nonlinearities in the original $`\varphi `$ equation will lead to the fluctuationless saturation of $`\varphi `$ at the value of unity for $`\xi 1`$, which corresponds to the fluctationless slope $`h\lambda ^{}\xi `$ behind the front: It is as if the KPZ equation has to be solved in the positive half-space with (roughly) a fixed boundary. On the other hand, by translating (11) back into $`h`$, we see that for large $`\xi `$ and $`t`$, the average interface shape $`h_{av}`$ should be given by $`h_{av}\mathrm{ln}(\alpha \xi _R/t^{3/2})+\lambda ^{}\xi \lambda _R^{}\xi _R\xi _R^2/4Dt.`$ (21) Thus, apart from the logarithmic term the average interface is essentially tilted but flat up to the time-dependent cross-over $`\xi _c\sqrt{4Dt}`$ , and beyond $`\xi _c`$ it has the shape of a downward curved parabola with time dependent curvature. Together with the fact that the nonlinear term in (20) gives an average nonzero growth velocity, this makes the problem into a nonstandard fluctuating interface problem. Our central approximation is now to consider the relaxing front in the essentially straight but fluctuating section between 0 and $`\sqrt{4Dt}`$ as a KPZ interface with time dependent length $`L=𝒪(\xi _c)`$. As the scaling exponents of the KPZ equation are robust with respect to a geometric change of the fluctuating surface , we use the KPZ scaling functions for the root mean square width $`W`$ of the interface $`h`$, $`W(L,t)=t^\beta Y\left({\displaystyle \frac{t}{L^z}}\right),\beta =1/3,z=3/2,`$ (22) where $`W=\sqrt{\overline{h(x,t)\overline{h}(x,t))^2}_\mu }`$, with the bar denoting a spatial average. The scaling function $`Y(s)`$ will depend on the shape of the roughening surface, but always has the limits $`Y(s)s^\beta `$ for $`s\mathrm{}`$, $`Y(0)\mathrm{const}`$. Inserting our approximation $`L\sqrt{t}`$, we get: $`W(L,t)L^{z\beta }(\sqrt{t})^{z\beta }=t^{1/4}.`$ (23) The final step of our argument is to convert this result in a prediction for the fluctuations of the front position. If we measure the position of the front by tracking a certain height $`c`$, $`\varphi (x_c,t)=\mathrm{const}=c`$, and use the relations (14) and (19), we find: $`\varphi (x_c,t)=e^{\lambda _R^{}(x_c\overline{v}_R^{}t)+h}=\mathrm{const}=c.`$ (24) This implies that fluctuations in $`h`$ are just identical with fluctuations in $`x_c`$. Therefore we get $`\mathrm{\Delta }(t)t^{1/4}`$ (25) which reproduces the scaling of our previous result (13). We have also performed numerical simulations of the noisy front equation (1) with $`a=0.3`$ (pushed) and $`a=1`$ (pulled, F-KPP Equation (8)) following the lines of . The initial condition was taken as a step function $`\varphi (x,0)=\theta (x_0x)`$. The numerical integration has been performed using a standard explicit Euler algorithm, in both cases the value of the noise was set to $`\epsilon =0.5`$, and the zero value of the spatial noise correlator $`C(0)`$ was chosen as the inverse spatial integration mesh, $`C(0)=1/\mathrm{\Delta }x`$ . The result is shown in Fig. 1, where the function $`\mathrm{\Delta }(t)`$ is plotted in both the pushed and the pulled case. The specific features of the pulled regime make the problem quite delicate from the numerical point of view. In order to minimize finite size effects, which are particularly worrisome in this regime , we have worked with a large system size ($`L=3000`$) and gridsize $`\mathrm{\Delta }x=1`$ (the change in $`v^{}`$ and $`D`$ due to the finite gridsize effect was taken into account following the prescription of ). This made sure that even at time $`t=600`$, the leading edge of the front never reached the boundary of the system. We have also checked our program and system size extensively both for deterministic and noisy fronts, taking into account grid and time step effects according to . FIG. 1. Diffusive and subdiffusive spreading of the front position. The dot-dashed curve correponds to the pushed case ($`a=0.3`$) and the solid one corresponds to the pulled case ($`a=1`$). The dashed straight line is the prediction (13), while the dotted line indicates a slope $`1/2`$. Our final result, based on averaging over 10000 front realizations, is shown in Fig. 1; it clearly confirms the subdiffusive behaviour predicted by our analytical arguments. Quantitatively, when we associate a single effective exponent with the late time slope in the log-log plot of Fig. 1, we get an effective exponent of about 0.29 rather then $`1/4`$. Over the time interval we have studied, the actual value of $`\mathrm{\Delta }(t)`$ is somewhat larger than an asymptotic prediction (13), which is indicated with a dashed line. This may be due to the fact that (13) only gives the behavior for such long times that the time integral is dominated by its large $`t`$ behavior. The fact that $`\mathrm{\Delta }`$ is only of the order of $`4`$ at our latest times suggests that this asymptotic regime is only reached at very late times. Indeed, assuming that finite size effects are negligible, we attribute the fact that the effective exponent is slightly larger than $`1/4`$ to the presence of slow crossovers, which surely are present in the system. Some of these can be estimated, while others are more difficult to trace. (i) We already noticed previously that we are actually dealing with a slightly curved KPZ interface, for which the crossover scaling functions are not known, and that the way in which the cutoff $`\xi _c=𝒪(\sqrt{t})`$ enters the KPZ analysis requires further study. (ii) The corrections to our asymptotic estimates for the integrals in (7) are all of order $`1/\sqrt{t}`$, with possible logarithmic corrections . This indicates that the corrections to the scaling $`\mathrm{\Delta }t^{1/4}`$ are of order $`t^{1/4}`$, possibly with logarithmic corrections. (iii) If initially $`\varphi `$ falls off as $`\mathrm{exp}(\lambda _R^{}x)`$, then the associated KPZ interface remains straight towards $`\xi =\mathrm{}`$. For this case the KPZ scaling predicts $`\mathrm{\Delta }t^{1/3}`$. Presumably a crossover between exponent $`1/3`$ and $`1/4`$ could be present when starting with an initial condition slightly faster decaying than $`\mathrm{exp}(\lambda _R^{}x)`$. The identification of such a crossover and the modification of the global exponent due to these special initial conditions is an issue that will be addressed elsewhere. We finally stress, that our results apply to a much larger class of equations than nonlinear diffusion equations (1). The methods of generalization are analogous to those of ; a closely related result is the general argument put forward in that noisy pulled fronts in more than one dimension should not obey KPZ scaling. We thank J. Casademunt and L. Schäfer for useful discussions. A.R. thanks the Instituut-Lorentz for kind hospitality. He was supported by the European Commission project ERBFMRX-CT96-0085 and U.E. by the Dutch Science Foundation NWO.
warning/0004/cond-mat0004171.html	ar5iv	text	# 𝑡-𝑈-𝑊 Model of a 𝑑_{𝑥²-𝑦²} Superconductor in the Proximity of an AF Mott Insulator: Diagrammatic Studies vs. Quantum-Monte-Carlo Simulations ## I Introduction One salient aspect of the high-$`T_c`$ materials is the vicinity of two, at first sight rather different, states of matter, superconductivity (SC) and antiferromagnetism (AF) in their phase diagram. The transition between the undoped (AF) system at half-filling and the SC phase is driven by doping with mobile holes. In most of the materials, this transition is not direct, and a disordered “spin-glass” phase occurs in between. However, it has been argued that the “clean” material would display a direct transition from AF to SC phases, and that the spin-glass phase occurs due to the high sensitivity to impurity disorder in the vicinity of the phase transition. A direct transition from an insulating into a SC phase in a quasi-two-dimensional system (such as the high-$`T_c`$ materials) is a very interesting, yet insufficiently understood issue. In fact, it is not clear whether this transition is second order down to zero temperature, and thus is related to a quantum critical point, or whether there is a finite-temperature classical bicritical point. In the framework of the projected SO(5)-theory of high-$`T_c`$ superconductivity , it has been suggested that the AF and the SC phases may indeed coexist in some portion of the temperature versus doping phase diagram. Another open question is the nature, i. e. the universality class of this transition. For example, it has been suggested that this transition may be controlled by an SO(5)-symmetric fixed point . SO(5) symmetry is thus restored in the long-wavelength limit , and AF and SC can be described in terms of a unique superspin vector in the vicinity of the critical point. Many efforts have been directed towards studying the AF-SC transition in strongly correlated lattice models by numerical techniques such as Quantum-Monte-Carlo (QMC) simulations. As a relevant model, the Hubbard model, is widely accepted for the description of salient features of high-$`T_c`$ materials. Unfortunately, it is quite difficult to study large enough Hubbard-model systems by QMC, due to the occurrence of the minus-sign problem at finite doping. The numerical problem can, in principle, be cured, if one can drive the AF-SC transition by means of a parameter, alternative to the doping, which conserves particle-hole symmetry and, therefore, avoids the tedious minus-sign problem. This idea was followed through by Assaad, Imada, and Scalapino (AIS) in terms of their so-called $`tUW`$ model. It rests on adding an interaction term $`W`$, which depends on the square of the nearest-neighbor hopping. This $`W`$ term can be obtained from a Su-Schrieffer-Heeger type of electron-phonon interaction in the antiadiabatic limit . In QMC simulations this $`tUW`$ model exhibits a transition from an antiferromagnet to a d-wave superconductor at half-filling and at a critical value of the interaction $`W_c0.3t`$ ($`U=4t`$, $`T=0K`$). The QMC data of AIS supports the picture of a continuous quantum phase transition in the sense that the magnetization vanishes continuously at the critical point. The disadvantage of the $`tUW`$ model is that the bandwidth grows substantially with $`W`$. Therefore, one of us suggested to introduce a phase factor in the $`W`$-term which has a d-wave like symmetry. Although this latter model solves the problem of the bandwidth, the existence of a phase transition to a d-wave superconductor remains open. While QMC calculations provide an essentially exact description of the properties of the model, semi-analytical, i. e. diagrammatic, calculations allow for a more direct understanding of the processes which are responsible for a given phenomenon. For this reason in this paper, we carry out a systematic diagrammatic study of the $`tUW`$ model. We first consider in Sec. III the simple Hartree-Fock level, which, due to the complexity of the interaction terms, allows for different broken symmetry phases. However, a careful comparison of the energies of these phases shows that the antiferromagnetic phase is always the stablest one, even for very large values of $`W`$. This holds for both versions of the $`tUW`$ model which are considered, i. e. with and without phase factors. Moreover, the only allowed superconducting solution in the simple $`tUW`$ model has an $`s`$-wave symmetry, while $`d`$-wave symmetry is not allowed. These mean-field results are in strong contrast with the QMC calculations, which predict a transition to a $`d`$-wave SC state at some finite $`W`$ . On the other hand, in the AF region our mean-field results are in very good accord with QMC, in particular concerning single-particle dispersions, as shown in Sec. III. The fact that the transition to the superconducting state does not come out correctly is of no surprise within an Hartree-Fock approximation. Indeed the relevant transition to the d-wave SC state is dominantly driven by an effective attractive interaction mediated by spin fluctuation, as is the case for the Hubbard model . For this reason, one has to consider the effect of spin fluctuations beyond the Hartree-Fock level, in order to reach the SC state. This is done in Sec. IV, where we carry out a complete RPA summation of all particle-hole diagrams (both bubbles and ladders), with Hartree-Fock Green’s functions, in order to obtain the frequency- and momentum-dependent spin and charge susceptibility. The solution of the corresponding Bethe-Salpeter equation is technically quite difficult to achieve and significantly more demanding than the standard case of the simple Hubbard model . This is due to the finite extension of the interaction, as well as its dependence on all (three) momenta, and not on the momentum transfer only. By changing to a mixed real-space momentum-space representation, we demonstrate that it can be reduced, for generic momenta, to the inversion of a $`52\times 52`$ matrix. Next, in some analogy to the RPA-analysis of the $`tU`$ Hubbard model by Schrieffer, Wen, and Zhang , we derive the effective two-particle interaction vertex in the static limit, and solve the associated BCS equation. As for the simple Hubbard model studied in Ref. , the $`d`$-wave solution turns out to be the only stable one, in accordance with QMC results. On the other hand, as expected, the $`s`$-wave solution is unstable, due to the strong on-site Hubbard repulsion $`U`$. In the d-wave phase, we obtain a decreasing superconducting gap as a function of $`W`$, in spite of the fact that the attraction between the quasiparticles should be increased by $`W`$. This is due, on the one hand, to the approximation of taking an energy cutoff for the effective interaction, which has been chosen to be of the order of the AF gap, which, in turn, decreases with increasing $`W`$. On the other hand, the reduction of the density of states at the Fermi level, which is related to the broadening of the bands produced by $`W`$, contributes in reducing the superconducting gap. Our paper is organized as follows. In Sec. II the $`tUW`$ model with and without phase factors is introduced, and briefly summarized In Sec. III, we carry out the Hartree-Fock mean-field study of the antiferromagnetic phase. We discuss the HF results and compare them with QMC calculations. In Sec. IV we derive and solve the Bethe-Salpeter equation, i. e. we account for fluctuation effects in the time-dependent HF or generalized RPA scheme. We obtain spin and charge susceptibilities, as well as the effective interaction vertex. In Sec. V, we write down and solve the BCS gap equation, obtained from this effective interaction within a static approximation. Finally, we present our conclusions in Sec. VI, partly based on detailed comparisons with QMC data. ## II Model The Hamiltonian of the $`tUW`$ model is given by $$=\frac{t}{2}\underset{\stackrel{}{i}}{}K_\stackrel{}{i}+U\underset{\stackrel{}{i}}{}(n_{\stackrel{}{i},}\frac{1}{2})(n_{\stackrel{}{i},}\frac{1}{2})W\underset{\stackrel{}{i}}{}\stackrel{~}{K}_\stackrel{}{i}^2$$ (1) with the hopping kinetic energy $$K_\stackrel{}{i}=\underset{\sigma ,\stackrel{}{\delta }}{}(c_{\stackrel{}{i},\sigma }^{}c_{\stackrel{}{i}+\stackrel{}{\delta },\sigma }^{\text{}}+c_{\stackrel{}{i}+\stackrel{}{\delta },\sigma }^{}c_{\stackrel{}{i},\sigma }^{\text{}})$$ (2) and $$\stackrel{~}{K}_\stackrel{}{i}=\underset{\sigma ,\stackrel{}{\delta }}{}f(\stackrel{}{\delta })(c_{\stackrel{}{i},\sigma }^{}c_{\stackrel{}{i}+\stackrel{}{\delta },\sigma }^{\text{}}+c_{\stackrel{}{i}+\stackrel{}{\delta },\sigma }^{}c_{\stackrel{}{i},\sigma }^{\text{}}),$$ (3) where in its original formulation $`f(\stackrel{}{\delta })=1`$. As mentioned in the Introduction, this model was introduced by Assaad, Imada, and Scalapino , in order to study the antiferromagnetic-superconducting transition at half-filling. As stressed by these authors , the particular choice of $`W`$ was mainly motivated formally as a means of introducing the desired quantum transition from the insulating (AF) to the SC state. The choice of the interaction also guaranteed that no fermion-sign problem was encountered in the QMC simulations at half-filling. While there are various approximate ways to physically justify the form of the microscopic Hamiltonian, it runs into problems when directly compared to the high-$`T_c`$ cuprates. One of the problems of this model is the presence of unphysically broad bands, which also cause problems in the numerical QMC evaluation. For this reason, one of us suggested to introduce in the $`\stackrel{~}{K}_\stackrel{}{i}`$–term of Eq. (3) a d-wave like phase factor of the form $$f(\stackrel{}{\delta })=\{\begin{array}{cc}+1\hfill & \text{for }\stackrel{}{\delta }=\pm \stackrel{}{a}_x\hfill \\ 1\hfill & \text{for }\stackrel{}{\delta }=\pm \stackrel{}{a}_y\hfill \end{array}.$$ (4) The phase factor restores the correct width of the single–particle bands at half-filling. Unfortunately, however no transition to a superconductor in QMC simulations has been observed so far. The $`W`$ term contains four different processes (see e.g. Ref. ), among them single–particle terms that renormalize the chemical potential and permit single–particle hopping between second and third nearest-neighbor sites as well as singlet and triplet scattering terms. These terms are not expected to be relevant for the low-energy physics . However, the most interesting term for the quantum phase transition is the fourth term which generates singlet pair-hopping and produces an antiferromagnetic exchange interaction, i. e. $$_W^{(4)}=2W\underset{\stackrel{}{i},\stackrel{}{\delta },\stackrel{}{\delta }^{}}{}f(\stackrel{}{\delta })f(\stackrel{}{\delta }^{})\mathrm{\Delta }_{\stackrel{}{i},\stackrel{}{\delta }^{}}^{}\mathrm{\Delta }_{\stackrel{}{i},\stackrel{}{\delta }}^{\text{}},$$ (5) where $`\mathrm{\Delta }_{\stackrel{}{i},\stackrel{}{\delta }}^{}=(c_{\stackrel{}{i},}^{}c_{\stackrel{}{i}+\stackrel{}{\delta },}^{}c_{\stackrel{}{i},}^{}c_{\stackrel{}{i}+\stackrel{}{\delta },}^{})/\sqrt{2}`$. For $`\stackrel{}{\delta }=\stackrel{}{\delta }^{}`$ the terms in $`_W^{(4)}`$ contribute to the exchange giving: $$2W\underset{\stackrel{}{i},\stackrel{}{\delta }}{}(\stackrel{}{S}_\stackrel{}{i}\stackrel{}{S}_{\stackrel{}{i}+\stackrel{}{\delta }}\frac{1}{4}n_\stackrel{}{i}n_{\stackrel{}{i}+\stackrel{}{\delta }}).$$ (6) ## III Hartree-Fock Calculations The details of our HF calculation are given in Appendix A. After solving the self-consistent equations for the mean-field parameters in Eqs. (A5) to (A7) and (A21) to (A26) (Appendix A), we arrive at the following results: Figure 1 displays the free energy of the different phases (antiferromagnetic, superconducting and paramagnetic) for the simple $`tUW`$ model (top) and the $`tUW`$ model with phase factors (bottom) as a function of $`W`$ for fixed $`U=4t`$ and $`T=0K`$. For the sake of comparison, we only plot the difference to the paramagnetic energy. In the simple $`tUW`$ model (Fig. 1, top) the antiferromagnetic solution (AF) is always the most favorable. However, with increasing $`W`$ the energy of this solution approaches the paramagnetic solution (PM). The superconducting solution (SC) has a much higher energy than the other solutions. The only possible superconducting solutions have $`s_1=s_3=0`$ (see Appendix A: no on-site pairing, due to $`U`$) and $`s_20`$ (nearest-neighbor singlet pairing), while for $`W0.3t`$ there exists no superconducting solution. The superconducting order parameter $`s_2`$ corresponds to an s-wave like symmetry. The transition from an antiferromagnet to a $`d_{x^2y^2}`$-superconductor observed in QMC simulations is not reproduced at the mean-field level. In the $`tUW`$ model with phase factors (Fig. 1, bottom) the mean-field ground state is also antiferromagnetic (AF). Here, however, in contrast to the simple $`tUW`$ model, the difference in energy with the paramagnetic (PM) solution is increasing with increasing $`W`$. As discussed in Appendix A, there also exist two different superconducting solutions, that lie energetically between the antiferromagnetic and the paramagnetic solution and evolve continuously from the paramagnetic solution at $`W=0t`$. In the first solution, $`s_1`$ and $`s_3`$ are nonvanishing, while $`s_2=0`$ (s-wave). The order parameter has a s-wave symmetry with a superimposed weak modulation of the gap. In the second, energetically more favorable, solution one has $`s_1=s_3=0`$, and $`s_20`$, yielding an order parameter with d-wave symmetry. Notice that also in QMC simulations at half-filling, no transition to a superconductor was found in the $`tUW`$ model with phase factors, in agreement with our results. However, the mean-field result in Fig. 1 is promising in direction of doping away from half-filling, where the AF phase is suppressed. The band structure of the antiferromagnetic solutions is evaluated along the usual paths through the Brillouin zone, as shown in Fig. 2. Fig. 3 (top) gives the bands of the simple $`tUW`$ model for $`W=0.15t`$. One can recognize easily that the bands are much wider than in the Hubbard model but their shape is nearly unaltered. This means that if one would scale the bands by a factor $`\frac{1}{3}`$, they would be almost identical. The effect of $`W`$ seems thus to be a mere “dilatation” of the bands. In the $`tUW`$ model with phase factors, things look quite different. The bands are plotted in Fig. 3 (middle) along the path shown in Fig. 2(a) and in Fig. 3 (bottom) along the path shown in Fig. 2(b) with $`W=0.05t`$. The width of the bands is nearly the same as in the Hubbard model, except for the lifting of the degeneracy along the boundaries of the magnetic Brillouin zone (MBZ). At $`\stackrel{}{k}=(\pi ,0)`$ a kind of double-hump structure can be seen like it appears in $`t`$-$`t^{}`$-$`t^{\prime \prime }`$ models to describe high-$`T_c`$ superconductors . This can be explained as follows: The $`W`$ term contains also hopping processes to second and third nearest neighbor sites which, due to the phase factors, have the same sign as in the standard fit parameters $`t,t^{},t^{\prime \prime }`$ which are often used to adjust the bands to the experimental data of Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> and YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>. If one compares the antiferromagnetic bands to the QMC data of the $`tUW`$ model as in Figs. 4 and 5, one gets a very good agreement. The width of the bands as well as the antiferromagnetic gap were reproduced excellently. In addition, the energy bands of the $`tUW`$ model with phase factors (Fig. 3) show the same double-hump structure at $`\stackrel{}{k}=(\pi ,0)`$ like it is seen in the QMC spectral weight $`A(\stackrel{}{k},\omega )`$ (Fig. 5). They also reproduce well the lift of the degeneracy along the boundaries of the MBZ. Finally, we want to look at two characteristic features of the antiferromagnetic solution: the sublattice magnetization $`m`$ and the Mott-Hubbard gap. The sublattice magnetization defined as $$m=\|c_{\stackrel{}{i},}^{}c_{\stackrel{}{i},}^{\text{}}c_{\stackrel{}{i},}^{}c_{\stackrel{}{i},}^{\text{}}\|,$$ (7) is plotted in Fig. 6 (top) as a function of $`W`$. As expected from the behavior of the free energy (Fig. 1), the sublattice magnetization decreases with increasing $`W`$ in the simple $`tUW`$ model. On the other hand, the sublattice magnetization of the $`tUW`$ model with phase factors is getting stronger with increasing $`W`$. This is also confirmed by QMC data , which show an amplification of the antiferromagnetic correlations with increasing $`W`$. A similar picture occurs for the antiferromagnetic gap (Fig. 6, bottom). Like the sublattice magnetization, the Mott-Hubbard gap is decreasing with increasing $`W`$ in the simple $`tUW`$ model, while it increases (nearly linear) with $`W`$ in the model with phase factors. In summary, the Hartree-Fock calculation gives the antiferromagnetic solution as ground state for any values of $`W`$ in both models. However, qualitatively there are remarkable differences between the two models with and without phase factors. A comparison of these results with the QMC data shows that the antiferromagnetic phase is described quite well by the mean-field approximation. On the other hand, the mean-field level is not able to reproduce the transition to a $`d_{x^2y^2}`$-superconductor at $`W_c0.3t`$ observed in QMC simulations in the simple $`tUW`$ model. This is of no surprise since results of an Hartree-Fock approximation at finite values of the interaction should be taken with due care and cannot give decisive conclusions about the correct phase diagram of a model without a comparison with more reliable calculations, such as, e. g. QMC. Nevertheless, for large $`W`$, for which the AF gap becomes small, one would expect the antiferromagnetic solution to become instable with respect to fluctuations beyond the mean-field level. This is what we analyze in the next section. ## IV Time-dependent Hartree-Fock (generalized RPA) As demonstrated in the previous section, the HF mean-field approximation is not sufficient to describe the transition to a $`d_{x^2y^2}`$-superconductor occurring in the simple $`tUW`$ model according to QMC simulations. For that reason, we carried out an improved calculation, including charge- and spin-density fluctuations. This has been done by means of a time-dependent HF or generalized random phase approximation (RPA), in which we summed both ”bubble” and ”ladder” particle-hole diagrams. In contrast to the fluctuation exchange approximation (FLEX), the Green’s functions are not calculated selfconsistently, but taken over from the Hartree-Fock results, as it has been done in Ref. . ### A Hartree-Fock Correlation Function $`L^0`$ and Interaction Vertex $`\mathrm{\Gamma }^0`$ The $`2\times 2`$ antiferromagnetic Hartree-Fock Green’s function can be written as (see Appendix A, Eqs. (A4) to (A12)): $$𝒢^{HF}(\stackrel{}{k},\omega ,\sigma )=\left(\begin{array}{cc}i\omega +\epsilon (\stackrel{}{k})& \sigma \mathrm{\Delta }(\stackrel{}{k})\\ \sigma \mathrm{\Delta }(\stackrel{}{k})& i\omega \epsilon (\stackrel{}{k})\end{array}\right)\frac{1}{\omega ^2E^2(\stackrel{}{k})}.$$ (8) With this Green’s function we can construct the Hartree-Fock two-particle propagator $`L^0`$: $$L^0{}_{\begin{array}{c}\sigma _1\sigma _2\\ \sigma _1^{}\sigma _2^{}\end{array}}{}^{\begin{array}{c}m_1m_2\\ m_1^{}m_2^{}\end{array}}(\stackrel{}{k}_1,\stackrel{}{k}_2,\stackrel{}{q},\omega _1,\omega _2,\nu )=\delta _{\sigma _1\sigma _2}\delta _{\sigma _1^{}\sigma _2^{}}\delta _{\stackrel{}{k}_1\stackrel{}{k}_2}\delta _{\omega _1\omega _2}𝒢_{}^{HF}{}_{m_2m_1}{}^{}(\stackrel{}{k}_1,\omega _1,\sigma _1)𝒢_{}^{HF}{}_{m_1^{}m_2^{}}{}^{}(\stackrel{}{k}_1\stackrel{}{q},\omega _1\nu ,\sigma _1^{}),$$ (9) where the set of $`m_i`$ stand for the indices of the $`2\times 2`$ matrix in Eq. (8). After a unitary transformation with help of the Pauli matrices, i. e. $$\stackrel{~}{L}^0=U^{}L^0U,$$ (10) where $$U_{\begin{array}{c}\sigma _1\\ \sigma _2\end{array}\alpha }=\frac{1}{\sqrt{2}}\sigma _{\sigma _1\sigma _2}^\alpha ,\alpha =0,x,y,z,$$ (11) we can write the correlation function in the charge-/ spin-channel representation as: $$\begin{array}{cc}& \stackrel{~}{L}_{ab}^0{}_{}{}^{\begin{array}{c}m_1m_2\\ m_1^{}m_2^{}\end{array}}(k_1,k_2,q)=\hfill \\ & \left(\begin{array}{cccc}\stackrel{~}{L}_{\mathrm{0\hspace{0.17em}0}}^0{}_{}{}^{\begin{array}{c}m_1m_2\\ m_1^{}m_2^{}\end{array}}(k_1,k_2,q)& 0& 0& \stackrel{~}{L}_{0z}^0{}_{}{}^{\begin{array}{c}m_1m_2\\ m_1^{}m_2^{}\end{array}}(k_1,k_2,q)\\ 0& \stackrel{~}{L}_+^0{}_{}{}^{\begin{array}{c}m_1m_2\\ m_1^{}m_2^{}\end{array}}(k_1,k_2,q)& 0& 0\\ 0& 0& \stackrel{~}{L}_+^0{}_{}{}^{\begin{array}{c}m_1m_2\\ m_1^{}m_2^{}\end{array}}(k_1,k_2,q)& 0\\ \stackrel{~}{L}_{z\mathrm{\hspace{0.17em}0}}^0{}_{}{}^{\begin{array}{c}m_1m_2\\ m_1^{}m_2^{}\end{array}}(k_1,k_2,q)& 0& 0& \stackrel{~}{L}_{zz}^0{}_{}{}^{\begin{array}{c}m_1m_2\\ m_1^{}m_2^{}\end{array}}(k_1,k_2,q)\end{array}\right).\hfill \end{array}$$ (12) Here and in the following: $`k=(\stackrel{}{k},\omega )`$, $`q=(\stackrel{}{q},\nu )`$ and so on. In contrast to the Hubbard model, also the non-diagonal elements which couple the charge channel to the longitudinal spin channel have to be taken into account. For the following calculations it is also advantageous to transform from the representation $$\stackrel{~}{L}_{ab}^0{}_{}{}^{\begin{array}{c}m_1m_2\\ m_1^{}m_2^{}\end{array}}(\stackrel{}{k}_1,\stackrel{}{k}_2,\stackrel{}{q},\omega _1,\omega _2,\nu )\text{ with }\stackrel{}{k}_1,\stackrel{}{k}_2,\stackrel{}{q}MBZ$$ (13) to the representation $$\overline{L}_{ab}^0(\stackrel{}{k}_1,\stackrel{}{k}_2,\omega _1,\omega _2;\stackrel{}{q}+n\stackrel{}{Q},\stackrel{}{q}+n^{}\stackrel{}{Q},\nu )\text{ with }\stackrel{}{k}_1,\stackrel{}{k}_2BZ,\stackrel{}{q}MBZ,$$ (14) where $`n`$, $`n^{}`$ take the values $`\{0,1\}`$ and $`\stackrel{}{Q}=(\pi ,\pi )`$. In this representation, e. g., the longitudinal spin correlation function can be written as a matrix in the indices $`n`$, $`n^{}`$: $$\begin{array}{cc}\hfill \overline{L}_{zz}^0& (k_1,k_2;\stackrel{}{q}+n\stackrel{}{Q},\stackrel{}{q}+n^{}\stackrel{}{Q},\nu )=\hfill \\ & \delta _{\omega _1\omega _2}\left(\begin{array}{c}\delta _{\stackrel{}{k}_1\stackrel{}{k}_2}\stackrel{~}{𝒢}_{}^{HF}{}_{\mathrm{1\hspace{0.17em}1}}{}^{}(k_1)\stackrel{~}{𝒢}_{}^{HF}{}_{\mathrm{1\hspace{0.17em}1}}{}^{}(k_1q)+\delta _{\stackrel{}{k}_1\stackrel{}{k}_2+\stackrel{}{Q}}\stackrel{~}{𝒢}_{}^{HF}{}_{\mathrm{2\hspace{0.17em}1}}{}^{}(k_1)\stackrel{~}{𝒢}_{}^{HF}{}_{\mathrm{1\hspace{0.17em}2}}{}^{}(k_1q)0\\ 0\delta _{\stackrel{}{k}_1\stackrel{}{k}_2}\stackrel{~}{𝒢}_{}^{HF}{}_{\mathrm{1\hspace{0.17em}1}}{}^{}(k_1)\stackrel{~}{𝒢}_{}^{HF}{}_{\mathrm{2\hspace{0.17em}2}}{}^{}(k_1q)+\delta _{\stackrel{}{k}_1\stackrel{}{k}_2+\stackrel{}{Q}}\stackrel{~}{𝒢}_{}^{HF}{}_{\mathrm{2\hspace{0.17em}1}}{}^{}(k_1)\stackrel{~}{𝒢}_{}^{HF}{}_{\mathrm{2\hspace{0.17em}1}}{}^{}(k_1q)\end{array}\right).\hfill \end{array}$$ (15) Here the $`\stackrel{~}{𝒢}^{HF}`$ are spin independent Green’s functions (i. e. $`𝒢^{HF}`$ (Eq. (8)) with spin index $`\sigma `$ set equal to $`+1`$). The interaction vertex in this representation is given by $$\begin{array}{cc}\hfill \overline{\mathrm{\Gamma }}_{ab}^0(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q}+n\stackrel{}{Q},\stackrel{}{q}+n^{}\stackrel{}{Q})=\frac{1}{\beta N}\{& \left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)U\hfill \\ & +\left(\begin{array}{cccc}2& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right)[\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)V^W(\stackrel{}{k}_2\stackrel{}{q},\stackrel{}{k}_1,\stackrel{}{q})\hfill \\ & +\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)V^W(\stackrel{}{k}_2\stackrel{}{q},\stackrel{}{k}_1,\stackrel{}{q}+\stackrel{}{Q})]\hfill \\ & \left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)V^W(\stackrel{}{k}_2\stackrel{}{q},\stackrel{}{k}_1,\stackrel{}{k}_1\stackrel{}{k}_2)\}.\hfill \end{array}$$ (16) It is written as a direct product of spin- and $`(n,n^{})`$-matrices with the $`W`$-dependent interaction $`V^W`$ given by $$\begin{array}{cc}\hfill V^W(\stackrel{}{k},\stackrel{}{k}^{},\stackrel{}{q})=8W[& (\mathrm{cos}k_x^{}\pm \mathrm{cos}k_y^{})(\mathrm{cos}k_x\pm \mathrm{cos}k_y)\hfill \\ \hfill +& (\mathrm{cos}(k_x^{}q_x)\pm \mathrm{cos}(k_y^{}q_y))(\mathrm{cos}k_x\pm \mathrm{cos}k_y)\hfill \\ \hfill +& (\mathrm{cos}(k_x+q_x)\pm \mathrm{cos}(k_y+q_y))(\mathrm{cos}k_x^{}\pm \mathrm{cos}k_y^{})\hfill \\ \hfill +& (\mathrm{cos}(k_x+q_x)\pm \mathrm{cos}(k_y+q_y))(\mathrm{cos}(k_x^{}q_x)\pm \mathrm{cos}(k_y^{}q_y))].\hfill \end{array}$$ (17) ### B Bethe-Salpeter Equation With help of the HF two-particle propagator $`\overline{L}`$ and the interaction vertex $`\overline{\mathrm{\Gamma }}`$, we can now write down the Bethe-Salpeter equation in the form: $$\begin{array}{cc}\hfill \overline{L}_{ab}(k_1,k_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )=& \overline{L}_{ab}^0(k_1,k_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )\hfill \\ & +\underset{\begin{array}{c}k_3,k_4\\ \stackrel{}{q}^{\prime \prime },\stackrel{}{q}^{\prime \prime \prime }\\ c,d\end{array}}{}\overline{L}_{ac}^0(k_1,k_3;\stackrel{}{q},\stackrel{}{q}^{\prime \prime },\nu )\overline{\mathrm{\Gamma }}_{cd}^0(\stackrel{}{k}_3,\stackrel{}{k}_4;\stackrel{}{q}^{\prime \prime },\stackrel{}{q}^{\prime \prime \prime })\overline{L}_{db}(k_4,k_2;\stackrel{}{q}^{\prime \prime \prime },\stackrel{}{q}^{},\nu ).\hfill \end{array}$$ (18) Unlike for the standard ($`W=0`$) Hubbard model, this equation cannot be easily inverted due to the complicated space and spin structure of the $`W`$ term. For this reason, we apply a method due to Hanke and Sham , which is based on the partial transformation of the interaction vertex back into real space. For short–range interactions, this yields finite (small-sized) matrices in real space. With the Fourier-transformed correlation function and interaction vertex $`\widehat{L}_{ab}(\stackrel{}{R}_1,\stackrel{}{R}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )`$ $`={\displaystyle \frac{1}{\beta N}}{\displaystyle \underset{\begin{array}{c}\stackrel{}{k}_1,\stackrel{}{k}_2\\ \omega _1,\omega _2\end{array}}{}}e^{i\stackrel{}{k}_1\stackrel{}{R}_1}e^{i\stackrel{}{k}_2\stackrel{}{R}_2}\overline{L}_{ab}(k_1,k_2;\stackrel{}{q},\stackrel{}{q}^{},\nu ),`$ (19) $`\widehat{\mathrm{\Gamma }}_{ab}^0(\stackrel{}{R}_1,\stackrel{}{R}_2;\stackrel{}{q},\stackrel{}{q}^{})`$ $`={\displaystyle \frac{\beta }{N}}{\displaystyle \underset{\stackrel{}{k}_1,\stackrel{}{k}_2}{}}e^{i\stackrel{}{k}_1\stackrel{}{R}_1}e^{i\stackrel{}{k}_2\stackrel{}{R}_2}\overline{\mathrm{\Gamma }}_{ab}^0(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{}),`$ (20) we obtain the Bethe-Salpeter equation in matrix form: $$\begin{array}{cc}\hfill \widehat{L}_{ab}(\stackrel{}{R}_1,\stackrel{}{R}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )=& \widehat{L}_{ab}^0(\stackrel{}{R}_1,\stackrel{}{R}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )\hfill \\ & +\widehat{L}_{ac}^0(\stackrel{}{R}_1,\stackrel{}{R}_3;\stackrel{}{q},\stackrel{}{q}^{\prime \prime },\nu )\widehat{\mathrm{\Gamma }}_{cd}^0(\stackrel{}{R}_3,\stackrel{}{R}_4;\stackrel{}{q}^{\prime \prime },\stackrel{}{q}^{\prime \prime \prime })\widehat{L}_{db}(\stackrel{}{R}_4,\stackrel{}{R}_2;\stackrel{}{q}^{\prime \prime \prime },\stackrel{}{q}^{},\nu ).\hfill \end{array}$$ (21) This is now a simple matrix equation which can easily be inverted for the interacting two-particle propagator $`\widehat{L}`$: $$\widehat{L}=\left(1\widehat{L}^0\widehat{\mathrm{\Gamma }}^0\right)^1\widehat{L}^0.$$ (22) In contrast to the Hubbard model we have to deal with complex $`26\times 26`$ matrices for the transverse spin-channel and with complex $`52\times 52`$ matrices for the coupled charge-/ longitudinal spin-channel. The RPA susceptibilities can be constructed by taking the $`(\stackrel{}{0},\stackrel{}{0})`$-matrix element in real space, i. e.: $$\chi _{ab}(\stackrel{}{q},\stackrel{}{q}^{};\nu )=\widehat{L}_{ab}(\stackrel{}{0},\stackrel{}{0};\stackrel{}{q},\stackrel{}{q}^{},\nu ).$$ (23) From this one obtains the retarded susceptibilities by the analytic continuation $`i\nu \omega +i\eta `$. Starting from the idea that spin fluctuations are responsible for the pairing of the quasi-particles, it is reasonable to first concentrate on the dynamic spin susceptibilities for the antiferromagnetic nesting vector $`\stackrel{}{Q}=(\pi ,\pi )`$ as a function of $`\omega `$, as we expect the strongest response there. In Fig. 7 (top), $`\chi _{zz}`$ is plotted for the simple $`tUW`$ model with $`W=0.1t`$, while in Fig. 7 (middle) $`\chi _{zz}`$ is displayed for the model with phase factors and $`W=0.05t`$. Both calculations can be compared with the result for the Hubbard model, reported in Fig. 7 (bottom). One can clearly see that the spectral weight is mainly concentrated at low frequencies and that it is abruptly decreasing at a frequency $`\omega 2\mathrm{\Delta }_{AF}`$, which corresponds to the antiferromagnetic gap. This behavior is most evident in the simple $`tUW`$ model. Moreover, one can recognize that the overall magnitude of the longitudinal spin susceptibility is biggest in the simple $`tUW`$ model and smallest in the $`tUW`$ model with phase factors. ### C Effective Interaction In this section, we calculate the effective 2-particle interaction mediated by the collective charge and spin fluctuations evaluated in the preceding section. Here, we restrict to the model without phase factor, since this is the only one which, according to QMC calculations, displays $`d`$-wave superconductivity. As a first step, we evaluate the fluctuation vertex, from which we can determine the modifications of the bare 2-particle interaction given by the $`tUW`$ Hamiltonian. The calculation of the fluctuation vertex is performed with the same techniques that were used to calculate the Bethe-Salpeter equation. This gives the expression: $$\begin{array}{cc}\hfill \overline{\mathrm{\Gamma }}_{ab}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )& =\underset{\begin{array}{c}k_3,k_4\\ \stackrel{}{q}^{\prime \prime },\stackrel{}{q}^{\prime \prime \prime }\\ c,d\end{array}}{}\overline{\mathrm{\Gamma }}_{ac}^0(\stackrel{}{k}_1,\stackrel{}{k}_3;\stackrel{}{q},\stackrel{}{q}^{\prime \prime })\overline{L}_{cd}(k_3,k_4;\stackrel{}{q}^{\prime \prime },\stackrel{}{q}^{\prime \prime \prime },\nu )\overline{\mathrm{\Gamma }}_{db}^0(\stackrel{}{k}_4,\stackrel{}{k}_2;\stackrel{}{q}^{\prime \prime \prime },\stackrel{}{q}^{})\hfill \\ & =\underset{\begin{array}{c}\stackrel{}{R}_1,\stackrel{}{R}_2\\ \stackrel{}{R}_3,\stackrel{}{R}_4\end{array}}{}\underset{\begin{array}{c}\stackrel{}{q}^{\prime \prime },\stackrel{}{q}^{\prime \prime \prime }\\ c,d\end{array}}{}\left(\frac{1}{\beta N}\right)e^{i\stackrel{}{k}_1\stackrel{}{R}_1}\widehat{\mathrm{\Gamma }}_{ac}^0(\stackrel{}{R}_1,\stackrel{}{R}_3;\stackrel{}{q},\stackrel{}{q}^{\prime \prime })\hfill \\ & \widehat{L}_{cd}(\stackrel{}{R}_3,\stackrel{}{R}_4;\stackrel{}{q}^{\prime \prime },\stackrel{}{q}^{\prime \prime \prime },\nu )\widehat{\mathrm{\Gamma }}_{db}^0(\stackrel{}{R}_4,\stackrel{}{R}_2;\stackrel{}{q}^{\prime \prime \prime },\stackrel{}{q}^{})e^{i\stackrel{}{k}_2\stackrel{}{R}_2},\hfill \end{array}$$ (24) which is diagrammatically represented in Fig. 8. Next, we have to change from the charge/spin channel representation back to the simple spin representation by inverting the transformation given by Eq. (10): $$\begin{array}{cc}\hfill \mathrm{\Gamma }^{\begin{array}{c}\sigma _1\sigma _2\\ \sigma _1^{}\sigma _2^{}\end{array}}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )& =U_{\begin{array}{c}\sigma _1\\ \sigma _1^{}\end{array}a}\overline{\mathrm{\Gamma }}_{ab}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )(U_{\begin{array}{c}\sigma _2\\ \sigma _2^{}\end{array}b})^{}\hfill \\ & =\frac{1}{2}(\overline{\mathrm{\Gamma }}_{\mathrm{0\hspace{0.17em}0}}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )\sigma _{\sigma _1\sigma _1^{}}^0\sigma _{\sigma _2^{}\sigma _2}^0+\overline{\mathrm{\Gamma }}_{0z}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )\sigma _{\sigma _1\sigma _1^{}}^0\sigma _{\sigma _2^{}\sigma _2}^z\hfill \\ & +\overline{\mathrm{\Gamma }}_{z\mathrm{\hspace{0.17em}0}}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )\sigma _{\sigma _1\sigma _1^{}}^z\sigma _{\sigma _2^{}\sigma _2}^0+\overline{\mathrm{\Gamma }}_{zz}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )\sigma _{\sigma _1\sigma _1^{}}^z\sigma _{\sigma _2^{}\sigma _2}^z\hfill \\ & +\overline{\mathrm{\Gamma }}_+(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )\stackrel{~}{\sigma }_{\sigma _1\sigma _1^{}}^{}\stackrel{~}{\sigma }_{\sigma _2^{}\sigma _2}^++\overline{\mathrm{\Gamma }}_+(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{},\nu )\stackrel{~}{\sigma }_{\sigma _1\sigma _1^{}}^+\stackrel{~}{\sigma }_{\sigma _2^{}\sigma _2}^{}).\hfill \end{array}$$ (25) Since we want to use the effective interaction in order to write an effective Hamiltonian, only the static limit of the fluctuation vertex has to be considered. Thus, simple diagrammatic rules yield for the correction to the bare interaction: $$\stackrel{~}{V}^{\begin{array}{c}\sigma _1\sigma _2\\ \sigma _1^{}\sigma _2^{}\end{array}}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{})=\mathrm{\Gamma }^{\begin{array}{c}\sigma _1\sigma _2\\ \sigma _1^{}\sigma _2^{}\end{array}}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{},0)(\beta N)(1).$$ (26) The effective interaction can then be written as (see Fig. 9 for diagrammatic representation) $$V_{eff}^{\begin{array}{c}\sigma _1\sigma _2\\ \sigma _1^{}\sigma _2^{}\end{array}}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{})=V^{\begin{array}{c}\sigma _1\sigma _2\\ \sigma _1^{}\sigma _2^{}\end{array}}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q})+\stackrel{~}{V}^{\begin{array}{c}\sigma _1\sigma _2\\ \sigma _1^{}\sigma _2^{}\end{array}}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q},\stackrel{}{q}^{})$$ (27) with the bare 2-particle interaction $$V^{\begin{array}{c}\sigma _1\sigma _2\\ \sigma _1^{}\sigma _2^{}\end{array}}(\stackrel{}{k}_1,\stackrel{}{k}_2;\stackrel{}{q})=U\delta _{\sigma _1\sigma _1^{}}\delta _{\sigma _2\sigma _2^{}}\delta _{\sigma _1\overline{\sigma }_2}+V^W(\stackrel{}{k}_2\stackrel{}{q},\stackrel{}{k}_1;\stackrel{}{q})\delta _{\sigma _1\sigma _1^{}}\delta _{\sigma _2\sigma _2^{}}$$ (28) and $`V^W`$ given by Eq. (17). Since one has to consider the pairing of the Hartree-Fock quasiparticles, the effective interaction has to be transformed into the $`\gamma `$-base, which produces additional coherence factors . For physical reasons, only the pairing of particles with opposite spin (singlet-pairing) was considered. In order to take into account the effect of the dynamics on top of our static approximation, we follow Ref. , and introduce a cutoff frequency $`\omega _c`$, analogous to the Debye frequency $`\omega _D`$ in the standard BCS theory. This ensures that only particles within an interval of width $`\mathrm{}\omega _c`$ above and below the Fermi energy $`E_F`$ are paired. The motivation for this cutoff frequency $`\omega _c`$ becomes clear if one looks at the spin susceptibilities of the simple $`tUW`$ model in Fig. 7. We have already shown in the preceding section that the spectral weight is concentrated at low frequencies. Under the condition that the spin fluctuations are responsible for the pairing of the quasi-particles, the longitudinal spin susceptibility gives quite naturally a cutoff frequency of the size of the antiferromagnetic gap ($`\omega _c2\mathrm{\Delta }_{AF}`$). This implies that for hole dopings away from half-filling, only the intra-valence band matrix elements have to be considered. The pairing part of the effective Hamiltonian can thus formally be written as $$^{pair}=\frac{1}{2N}\underset{\begin{array}{c}\stackrel{}{k},\stackrel{}{k}^{}\\ \sigma ,\sigma ^{}\end{array}}{^{}}V_{\sigma \sigma ^{}}^{pair}(\stackrel{}{k},\stackrel{}{k}^{})\mathrm{\Theta }(\omega _c\|E^v(\stackrel{}{k})E_F\|)\mathrm{\Theta }(\omega _c\|E^v(\stackrel{}{k}^{})E_F\|)\gamma _{\stackrel{}{k}^{},\sigma ^{}}^v\gamma _{\stackrel{}{k}^{},\sigma ^{}}^v\gamma _{\stackrel{}{k},\sigma }^v\text{}\gamma _{\stackrel{}{k},\sigma }^v\text{},$$ (29) where $`E^v(\stackrel{}{k})=E(\stackrel{}{k})`$ is the valance band energy. The direct interaction $`V_{\sigma \sigma }^{pair}(\stackrel{}{k},\stackrel{}{k}^{})`$, which is given by $`\sigma =\sigma ^{}`$ spin indices contains, besides the longitudinal spin fluctuations, also the bare interactions and the charge fluctuations. The exchange interaction with $`\sigma =\sigma ^{}`$ consists of the transverse spin fluctuations only. In Figs. 10 to 12, the direct interaction and the exchange interaction were plotted for different paths of $`\stackrel{}{k}`$ and $`\stackrel{}{k}^{}`$ through the magnetic Brillouin zone (MBZ) for the simple $`tUW`$ model and the Hubbard model, respectively. The exchange interaction was plotted there as $`V_{\sigma \sigma }^{pair}(\stackrel{}{k},\stackrel{}{k}^{})`$, since it has exactly this form, with negative sign, in the BCS gap equation (see Eq. (31)). Comparing the graphs for the simple $`tUW`$ model ($`W=0.1t`$) with the simple Hubbard model, one can easily see that the $`W`$ term amplifies the attractive parts of the direct interaction, whereas the attractive parts of the exchange interaction remain constant (see e. g. Fig. 11). However, the repulsive parts of the direct interaction and the exchange interaction were both attenuated considerably by increasing $`W`$ (see e. g. Fig. 12). Therefore, the pairing of the quasi-particles is favored in the simple $`tUW`$ model altogether. ## V BCS Gap Equation Finally, we want to solve the BCS gap equation for the effective pairing interaction. Starting point is the effective Hamiltonian, as obtained in the previous Section $$\begin{array}{cc}\hfill _{eff}=& \underset{\stackrel{}{k},\sigma }{^{}}(E^v(\stackrel{}{k})\mu )\gamma _{\stackrel{}{k},\sigma }^v\gamma _{\stackrel{}{k},\sigma }^v\text{}\hfill \\ & +\frac{1}{2N}\underset{\begin{array}{c}\stackrel{}{k},\stackrel{}{k}^{}\\ \sigma ,\sigma ^{}\end{array}}{^{}}V_{\sigma \sigma ^{}}^{pair}(\stackrel{}{k},\stackrel{}{k}^{})\mathrm{\Theta }(\omega _c\|E^v(\stackrel{}{k})E_F\|)\mathrm{\Theta }(\omega _c\|E^v(\stackrel{}{k}^{})E_F\|)\gamma _{\stackrel{}{k}^{},\sigma ^{}}^v\gamma _{\stackrel{}{k}^{},\sigma ^{}}^v\gamma _{\stackrel{}{k},\sigma }^v\text{}\gamma _{\stackrel{}{k},\sigma }^v\text{}.\hfill \end{array}$$ (30) With this Hamiltonian we want to study the superconducting properties of the simple $`tUW`$ model for different hole doping and different values of the model parameter $`W`$. The BCS gap equation becomes $$\mathrm{\Delta }(\stackrel{}{k})=\frac{1}{N}\underset{\stackrel{}{k}^{}}{^{}}(V_{}(\stackrel{}{k},\stackrel{}{k}^{})V_{}(\stackrel{}{k},\stackrel{}{k}^{}))\frac{\mathrm{\Delta }(\stackrel{}{k}^{})}{2E(\stackrel{}{k}^{})}.$$ (31) Here we used the abbreviations $`E(\stackrel{}{k})`$ $`=\sqrt{\xi ^2(\stackrel{}{k})+\mathrm{\Delta }^2(\stackrel{}{k})},`$ (32) $`\xi (\stackrel{}{k})`$ $`=E^v(\stackrel{}{k})\mu ,`$ (33) $`V_{\sigma \sigma ^{}}(\stackrel{}{k},\stackrel{}{k}^{})`$ $`=V_{\sigma \sigma ^{}}^{pair}(\stackrel{}{k},\stackrel{}{k}^{})\mathrm{\Theta }(\omega _A\|E^v(\stackrel{}{k})E_F\|)\mathrm{\Theta }(\omega _A\|E^v(\stackrel{}{k}^{})E_F\|).`$ (34) The gap equation (31) was iterated by assuming different symmetries of the superconducting order parameter, however, only $`d`$-wave solutions turn out to converge. The results for the simple $`tUW`$ model are shown in Fig. 13 as diamonds, triangles and squares. For all values of $`W`$ considered, the superconducting gap becomes zero at half filling, which is consistent with QMC results. It can be shown easily with the aid of equation (31), that even including interband matrix elements doesn’t change this result. This is due to the fact that the band gap is still quite large for these values of $`W`$ (cf. Fig. 3) for interband matrix elements to contribute substantially to the gap equation. On the other hand, all curves seem to indicate that the superconducting phase starts at very small doping, which is in contrast with experiments. While there are no conclusive QMC results about the $`tUW`$ model in this region, due to the minus sign problem, the limitations of our perturbative procedure applied for moderate values of $`U`$ suggest to consider this result with due care. Indeed, strong phase fluctuations, not included in a BCS-type of calculation like Eq.(31) are known to be important and to suppress superconductivity at small doping. Figure 13 also shows, that the optimal hole doping for the simple $`tUW`$ model is moving closer to half-filling with increasing $`W`$. For comparison, we show in Fig. 13 (stars) a result for the $`tUW`$ model with phase factor. One can see, that the superconducting gap is strongly suppressed for hole doping near half-filling, compared to the simple $`tUW`$ (and Hubbard) model. The effect is here even stronger than for the simple $`tUW`$ model, since in the model with phase factors the gap increases with $`W`$. It remains the question as to why the superconducting gap is getting smaller with increasing $`W`$ in the simple $`tUW`$ model, in spite of the fact that the bare attraction between the quasi-particles is enhanced by the $`W`$ term. To understand this, one has to keep in mind that there are other important quantities, like the cutoff frequency or the density of states at the Fermi level, which have an important effect on the magnitude of the superconducting gap. How the superconducting gap depends on these quantities is qualitatively seen already in the weak-coupling solution of the original BCS equation for an attractive $`\delta `$-potential : $$\mathrm{\Delta }=2\mathrm{}\omega _D\mathrm{exp}\left(\frac{1}{N(0)g}\right).$$ (35) A stronger coupling $`g`$ is increasing the superconducting gap, while a smaller cutoff frequency $`\omega _D`$ is decreasing it. Moreover, if the density of states at the Fermi level $`N(0)`$ is reduced, the superconducting gap is getting smaller, as well. This final point is decisive, since the $`W`$ term is broadening the energy bands considerably, thus reducing the density of states dramatically. Therefore, the fact that the superconducting gap is getting smaller with increasing $`W`$ can be understood within our approximation. ## VI Comparison with QMC Results and Conclusions In this paper, we have first shown that the standard Hartree-Fock approximation can describe the antiferromagnetic properties of the $`tUW`$ model, especially the single-particle energy bands, in surprisingly good agreement with the QMC simulations. On the other hand, like for the Hubbard model , it is unable to reproduce the transition to a $`d`$-wave superconductor observed in the simple $`tUW`$ model in numerical simulations (QMC). In order to partly overcome this short-coming we adopt a time-dependent HF or generalized RPA calculations, which we present in the second part of our paper. The standard Hartree-Fock approximation captures only the “high-energy” physics, and is thus capable of reproducing band-widths and the overall features of the single-particle spectral-function. Equivalently, short-range pairing-correlation functions should be well reproduced within this approximation. This is indeed the case. At short length scales (as shown in table 1) the extended s-wave vertex contribution to the pairing correlation functions is dominant in QMC simulations of the simple t-U-W model. It is only at larger distances that the d-wave pairing correlations become dominant. This crossover from short range to long range properties is not reproduced within the standard mean-field approximation. Alternatively, for the t-U-W model with phase factors d-wave pairing dominates in QMC simulations at small length scales (see table 2). On the other hand, by including charge and spin fluctuations within a time-dependent HF, or generalized RPA summation of ladder and bubble diagrams, we were able to obtain an effective attraction between the quasi-particles, which is enhanced by $`W`$. Moreover we obtain a corresponding superconducting order parameter with the correct d-wave symmetry. The QMC results indeed show that as $`W`$ increases at fixed $`U`$ or $`U`$ decreases at fixed $`W`$, an instability towards $`d`$-wave superconductivity occurs. To illustrate this, Fig. 14 plots the vertex contribution to the $`d`$-wave pairing correlations as well as the staggered spin susceptibility. As apparent at low $`U`$ for fixed $`W`$, the superconducting $`d`$-wave becomes the leading instability. Finally, it should be pointed out that our Hartree-Fock-Bethe-Salpeter procedure is of perturbative nature and, therefore, our results should be considered with due care, due to the fact that the interaction is not small. ## Acknowledgments The authors express their deep gratitude to Prof. M. Imada for many insightful discussions and clarifications of physical points. In this connection one of us (W. H.) is grateful to Prof. Imada for the warm hospitality experienced during visits at the ISSP in Tokyo. He acknowledges the support granted by the joint German - Japanese cooperation project (DFG – JSPS: 446 JAP – 113/114/0), which was central for the success of this work. Partial support by the DFG-Project Ha 1537/14-1 is also acknowledged. ## A Details of the HF calculations ### 1 Antiferromagnetic Mean Field We start with the following Ansatz for the mean field parameters: $`c_{\stackrel{}{i},\sigma }^{}c_{\stackrel{}{i},\sigma }^{\text{}}`$ $`=n_1+\sigma e^{i\stackrel{}{Q}\stackrel{}{i}}n_2,`$ (A1) $`c_{\stackrel{}{i},\sigma }^{}f(\stackrel{}{\delta })c_{\stackrel{}{i}+\stackrel{}{\delta },\sigma }^{\text{}}`$ $`=n_3,`$ (A2) $`c_{\stackrel{}{i}+\stackrel{}{\delta },\sigma }^{}f(\stackrel{}{\delta })f(\stackrel{}{\delta }^{})c_{\stackrel{}{i}+\stackrel{}{\delta }^{},\sigma }^{\text{}}`$ $`=n_4+\sigma e^{i\stackrel{}{Q}\stackrel{}{i}}n_5,`$ (A3) where $`\stackrel{}{Q}`$ denotes the antiferromagnetic nesting-vector $`\stackrel{}{Q}=(\pi ,\pi )`$. Here the expectation values $`\mathrm{}`$ contain also an average over all $`\stackrel{}{\delta }`$ and $`\stackrel{}{\delta }^{}`$, whenever explicitly present. At half filling, it is easy to show that $`n_1=\frac{1}{2}`$, and $`n_4=\frac{1}{8}`$. The mean–field Hamiltonian can then be written as $$_{tUW}^{MF}=\underset{\stackrel{}{k},\sigma }{^{}}(\stackrel{~}{c}_{\stackrel{}{k},\sigma }^{},\stackrel{~}{c}_{\stackrel{}{k}+\stackrel{}{Q},\sigma }^{})\left(\begin{array}{cc}\epsilon (\stackrel{}{k})& \sigma \mathrm{\Delta }(\stackrel{}{k})\\ \sigma \mathrm{\Delta }(\stackrel{}{k})& \epsilon (\stackrel{}{k})\end{array}\right)\left(\begin{array}{c}\stackrel{~}{c}_{\stackrel{}{k},\sigma }^{\text{}}\\ \stackrel{~}{c}_{\stackrel{}{k}+\stackrel{}{Q},\sigma }^{\text{}}\end{array}\right)+\stackrel{~}{E}_{tUW}^{MF}$$ (A4) with $`\epsilon (\stackrel{}{k})`$ $`=2t(\mathrm{cos}k_x+\mathrm{cos}k_y)96Wn_3(\mathrm{cos}k_x\pm \mathrm{cos}k_y),`$ (A5) $`\mathrm{\Delta }(\stackrel{}{k})`$ $`=Un_2+32Wn_58Wn_2(\mathrm{cos}k_x\pm \mathrm{cos}k_y)^2`$ (A6) and $$\stackrel{~}{E}_{tUW}^{MF}=+UNn_2^2+192WNn_3^264WN(\frac{1}{16}+n_2n_5).$$ (A7) Here, and in the following equations, the upper sign stands for the simple $`tUW`$ model, while the lower sign is used for the $`tUW`$ model with phase factors. The Hamiltonian can be diagonalized by the usual transformation $$\left(\begin{array}{c}\gamma _{\stackrel{}{k},\sigma }^c\text{}\\ \gamma _{\stackrel{}{k},\sigma }^v\text{}\end{array}\right)=\left(\begin{array}{cc}u(\stackrel{}{k})& \sigma v(\stackrel{}{k})\\ v(\stackrel{}{k})& \sigma u(\stackrel{}{k})\end{array}\right)\left(\begin{array}{c}\stackrel{~}{c}_{\stackrel{}{k},\sigma }^{\text{}}\\ \stackrel{~}{c}_{\stackrel{}{k}+\stackrel{}{Q},\sigma }^{\text{}}\end{array}\right)$$ (A8) with $`u(\stackrel{}{k})`$ $`=\left[{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\epsilon (\stackrel{}{k})}{E(\stackrel{}{k})}}\right)\right]^{\frac{1}{2}},`$ (A9) $`v(\stackrel{}{k})`$ $`=\left[{\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{\epsilon (\stackrel{}{k})}{E(\stackrel{}{k})}}\right)\right]^{\frac{1}{2}}`$ (A10) and $$E(\stackrel{}{k})=\sqrt{\epsilon ^2(\stackrel{}{k})+\mathrm{\Delta }^2(\stackrel{}{k})}.$$ (A11) The resulting Hamiltonian is given by $$_{tUW}^{MF}=\underset{\stackrel{}{k},\sigma }{^{}}E(\stackrel{}{k})(\gamma _{\stackrel{}{k},\sigma }^c\gamma _{\stackrel{}{k},\sigma }^c\text{}\gamma _{\stackrel{}{k},\sigma }^v\gamma _{\stackrel{}{k},\sigma }^v\text{})+\stackrel{~}{E}_{tUW}^{MF}.$$ (A12) Already at this point, one can see from equations (A5) and (A11) that the parameter $`n_3`$ produces extremely wide bands in the simple $`tUW`$ model. On the other hand, it is straightforward to see that for the alternative $`tUW`$ model with phase factors one must have $`n_30`$, in order to preserve the symmetry of the energy bands under interchange of x- and y-directions. This explains why the bandwidth of the $`tUW`$ model with phase factors is drastically smaller, and essentially the same as the one of the simple Hubbard model. ### 2 Superconducting Mean Field Here, the mean field parameters are chosen as: $`c_{\stackrel{}{i},\sigma }^{}c_{\stackrel{}{i},\sigma }^{\text{}}`$ $`=n_1,`$ (A13) $`c_{\stackrel{}{i},\sigma }^{}f(\stackrel{}{\delta })c_{\stackrel{}{i}+\stackrel{}{\delta },\sigma }^{\text{}}`$ $`=n_3,`$ (A14) $`c_{\stackrel{}{i}+\stackrel{}{\delta },\sigma }^{}f(\stackrel{}{\delta })f(\stackrel{}{\delta }^{})c_{\stackrel{}{i}+\stackrel{}{\delta }^{},\sigma }^{\text{}}`$ $`=n_4,`$ (A15) $`c_{\stackrel{}{i},\sigma }^{}c_{\stackrel{}{i},\sigma }^{}`$ $`=\sigma s_1,`$ (A16) $`c_{\stackrel{}{i},\sigma }^{}f(\stackrel{}{\delta })c_{\stackrel{}{i}+\stackrel{}{\delta },\sigma }^{}`$ $`=\sigma s_2,`$ (A17) $`c_{\stackrel{}{i}+\stackrel{}{\delta },\sigma }^{}f(\stackrel{}{\delta })f(\stackrel{}{\delta }^{})c_{\stackrel{}{i}+\stackrel{}{\delta }^{},\sigma }^{}`$ $`=\sigma s_3.`$ (A18) The superconducting parameters can be divided into two groups: $`s_2`$ stands for the nearest–neighbor singlet pairing, which is favored by the $`W`$ term (Eq. (5)). On the other hand, $`s_1`$ and $`s_3`$ for $`\delta =\delta ^{}`$ represent the on-site singlet pairing which is suppressed by the Hubbard $`U`$. The mean–field Hamiltonian thus is $$_{tUW}^{MF}=\underset{\stackrel{}{k}}{}(\stackrel{~}{c}_{\stackrel{}{k},}^{},\stackrel{~}{c}_{\stackrel{}{k},}^{\text{}})\left(\begin{array}{cc}\epsilon (\stackrel{}{k})& \mathrm{\Delta }(\stackrel{}{k})\\ \mathrm{\Delta }(\stackrel{}{k})& \epsilon (\stackrel{}{k})\end{array}\right)\left(\begin{array}{c}\stackrel{~}{c}_{\stackrel{}{k},}^{\text{}}\\ \stackrel{~}{c}_{\stackrel{}{k},}^{}\end{array}\right)+\stackrel{~}{E}_{tUW}^{MF}$$ (A19) with the single–particle energy $`\epsilon (\stackrel{}{k})`$ and the gap parameter $`\mathrm{\Delta }(\stackrel{}{k})`$ given by: $`\epsilon (\stackrel{}{k})`$ $`=2t(\mathrm{cos}k_x+\mathrm{cos}k_y)`$ (A20) $`96Wn_3(\mathrm{cos}k_x\pm \mathrm{cos}k_y),`$ (A21) $`\mathrm{\Delta }(\stackrel{}{k})`$ $`=+Us_1`$ (A22) $`8Ws_1(\mathrm{cos}k_x\pm \mathrm{cos}k_y)^2`$ (A23) $`32Ws_2(\mathrm{cos}k_x\pm \mathrm{cos}k_y)`$ (A24) $`32Ws_3.`$ (A25) The energy constant $`\stackrel{~}{E}_{tUW}^{MF}`$ stands for: $$\stackrel{~}{E}_{tUW}^{MF}=UNs_1^2+64WN(s_1s_3+s_2^2\frac{1}{16})+192WNn_3^2.$$ (A26) Eq. (A19) can be diagonalized with a Bogoliubov – de Gennes transformation, similar to Eq. (A8), i. e. $`\left(\begin{array}{c}\gamma _\stackrel{}{k}^c\text{}\\ \gamma _\stackrel{}{k}^v\text{}\end{array}\right)`$ $`=\left(\begin{array}{cc}u(\stackrel{}{k})& v(\stackrel{}{k})\\ v(\stackrel{}{k})& u(\stackrel{}{k})\end{array}\right)\left(\begin{array}{c}\stackrel{~}{c}_{\stackrel{}{k},}^{\text{}}\\ \stackrel{~}{c}_{\stackrel{}{k},}^{}\end{array}\right),`$ (A27) $`u(\stackrel{}{k})`$ $`=\left[{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\epsilon (\stackrel{}{k})}{E(\stackrel{}{k})}}\right)\right]^{\frac{1}{2}},`$ (A28) $`v(\stackrel{}{k})`$ $`=\left[{\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{\epsilon (\stackrel{}{k})}{E(\stackrel{}{k})}}\right)\right]^{\frac{1}{2}}\mathrm{sign}(\mathrm{\Delta }(\stackrel{}{k})).`$ (A29) The resulting Hamiltonian now becomes $$_{tUW}^{MF}=\underset{\stackrel{}{k}}{}E(\stackrel{}{k})(\gamma _\stackrel{}{k}^c\gamma _\stackrel{}{k}^c\text{}\gamma _\stackrel{}{k}^v\gamma _\stackrel{}{k}^v\text{})+\stackrel{~}{E}_{tUW}^{MF},$$ (A30) with the usual relation (Eq. (A11)) for $`E(\stackrel{}{k})`$. The paramagnetic solutions can be easily obtained by setting the antiferromagnetic parameters in Eq. (A4), or the superconducting parameters in Eq. (A19), equal to zero. ## Figures and Tables
warning/0004/quant-ph0004012.html	ar5iv	text	# Elementary Excitations in Trapped BEC and Zero Mode Problem ## 1 Introduction Since the realization of Bose-Einstein Condensation (BEC) in trapped alkali atomic vapors in 1995, much attentions have been attracted from both viewpoints of theories and experiments. Among all the recent developments, the study of collective excitations occupies an important position in studying the properties of trapped BEC. When the temperature approaches the critical temperature, the theoretical description of collective excitation is still proceeding at present. For the case that the temperature becomes near zero, the linearized GP equations can give correct numerical predictions in agreement with the experiments. However, the linearized GP equations are difficult to be solved analytically, and they include confused zero mode problem. In order to give some new light on this problem, we develop an equivalent formalism to the linearized GP equations based on a natural expansion of the atomic field operator. Our study follows discussions in ref. by Lewenstein and You, and the main difference is that we adopt specific eigenfunctions to expand the atomic field. It is noticed that the advantages of our expansion lie in the following facts: (1) The problem in solving the elementary excitation becomes the problem in diagonalizing the standard quadratic Hamiltonian of Boson operators, which has been completely solved in principle before. (2) Our formalism adopts a simpler form than those formulated by other expansion of the atomic field operator. (3) Since only the several lowest excitations dominate the behavior of trapped BEC system near zero temperature, the limited-level approximation (which is explained in details in Sec. 2) is rational in most cases. This paper is organized as the following: A general formalism for elementary excitations is present in Sec. 2, and then in Sec. 3 it is applied to the homogeneous case, which is interest since it can be used to be a good example to demonstrate the appearance and the physical meaning of the zero mode clearly. In Sec. 4, a conclusion will be made and some possible application of our special expansion will be discussed. ## 2 General Formalism The Hamiltonian of trapped BEC system is $$\widehat{H}=d^3\stackrel{}{r}[\widehat{\psi }^{}(\stackrel{}{r})(\frac{\mathrm{}^2^2}{2M}+V_{tr}(\stackrel{}{r})\mu _0)\widehat{\psi }(\stackrel{}{r})+\frac{1}{2}g\widehat{\psi }^{}(\stackrel{}{r})\widehat{\psi }^{}(\stackrel{}{r})\widehat{\psi }(\stackrel{}{r})\widehat{\psi }(\stackrel{}{r})],$$ (1) where $`\mu _0`$ is the chemical potential, $`g=\frac{4\pi \mathrm{}^2a_s}{M}`$ , $`a_s`$ the s-wave scattering length of the interatomic potential, M the atomic mass, $`V_{tr}(\stackrel{}{r})`$ the external potential, and $`\widehat{\psi }(\stackrel{}{r})`$ ($`\widehat{\psi }^{}(\stackrel{}{r})`$) the atomic annihilation (creation) field operator. Near $`T=0K`$, since almost all the atoms occupy the same one-particle quantum state $`\mathrm{\Phi }_0(\stackrel{}{r})`$, it is convenient to separate out the special mode for the ground state of BEC from the atomic field operator $$\widehat{\psi }(\stackrel{}{r})=\sqrt{N_0}\mathrm{\Phi }_0(\stackrel{}{r})+\delta \widehat{\psi }(\stackrel{}{r}).$$ (2) where $`\delta \widehat{\psi }(\stackrel{}{r})`$ represents the quantum fluctuation of $`\widehat{\psi }(\stackrel{}{r})`$ relative to$`\sqrt{N_0}\mathrm{\Phi }_0(\stackrel{}{r})`$. Notice that, at present $`d^3\stackrel{}{r}\delta \widehat{\psi }^{}(\stackrel{}{r})\delta \widehat{\psi }(\stackrel{}{r})_{ensemble}N_0`$, i.e. the thermal component atoms is negligible. Hence the Hamiltonian can be properly expanded in $`\delta \widehat{\psi }(\stackrel{}{r})`$ series while substituting formula(2) into formula(1). (Only the lower order of $`\delta \widehat{\psi }(\stackrel{}{r})`$ is important, so up to two orders of $`\delta \widehat{\psi }(\stackrel{}{r})`$ is maintained in our consideration.) From the above Hamiltonian, $`\mathrm{\Phi }_0(\stackrel{}{r})`$ can be determined in zero order of $`\delta \widehat{\psi }(\stackrel{}{r})`$ by making use of variation methods, which satisfy the time-independent GP equation $$(\frac{\mathrm{}^2^2}{2M}+V_{tr}(\stackrel{}{r})+gN_0\mathrm{\Phi }_0^{}(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r}))\mathrm{\Phi }_0(\stackrel{}{r})=\mu _0\mathrm{\Phi }_0(\stackrel{}{r}).$$ (3) Combining eqs. (1-3), and ignoring the terms of the third and fourth order of $`\delta \widehat{\psi }(\stackrel{}{r})`$ in eq.(1), the Hamiltonian can be simplified as, $`\widehat{H}`$ $``$ $`{\displaystyle \frac{1}{2}}gN_0^2{\displaystyle d^3\stackrel{}{r}\mathrm{\Phi }_0^{}(\stackrel{}{r})\mathrm{\Phi }_0^{}(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r})}`$ (4) $`+{\displaystyle d^3\stackrel{}{r}\delta \widehat{\psi }^{}(\stackrel{}{r})(\frac{\mathrm{}^2^2}{2M}+V_{tr}(\stackrel{}{r})\mu _0+2gN_0\mathrm{\Phi }_0^{}(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r}))\delta \widehat{\psi }(\stackrel{}{r})}`$ $`+{\displaystyle \frac{1}{2}}gN_0{\displaystyle d^3\stackrel{}{r}\mathrm{\Phi }_0^{}(\stackrel{}{r})\mathrm{\Phi }_0^{}(\stackrel{}{r})\delta \widehat{\psi }(\stackrel{}{r})\delta \widehat{\psi }(\stackrel{}{r})}`$ $`+{\displaystyle \frac{1}{2}}gN_0{\displaystyle d^3\stackrel{}{r}\delta \widehat{\psi }^{}(\stackrel{}{r})\delta \widehat{\psi }^{}(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r})}.`$ In order to further simplify the Hamiltonian, we select a set of special complete wave functions $`\left\{\mathrm{\Phi }_n(\stackrel{}{r})\right\}`$ to expand the atomic field operator, $$\widehat{\psi }(\stackrel{}{r})=\underset{n0}{}\widehat{a}_n\mathrm{\Phi }_n(\stackrel{}{r})+\widehat{A}_0\mathrm{\Phi }_0(\stackrel{}{r}),$$ (5) where$`\left\{\mathrm{\Phi }_n(\stackrel{}{r})\right\}`$ satisfy the following equations $$(\frac{\mathrm{}^2^2}{2M}+V_{tr}(\stackrel{}{r})+gN_0\mathrm{\Phi }_0^{}(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r}))\mathrm{\Phi }_n(\stackrel{}{r})=\mu _n\mathrm{\Phi }_n(\stackrel{}{r}),$$ (6) $`\widehat{a}_n`$ is the annihilation boson operator of the single state $`\mathrm{\Phi }_n(\stackrel{}{r})`$ $`(n0)`$, $`\widehat{A}_0`$ the annihilation boson operator of the single state $`\mathrm{\Phi }_0(\stackrel{}{r})`$. Due to the fact that $`(\frac{\mathrm{}^2^2}{2M}+V_{tr}(\stackrel{}{r})+gN_0\mathrm{\Phi }_0^{}(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r}))`$ is Hermitian, its eigenstates $`\mathrm{\Phi }_n(\stackrel{}{r})`$ $`(n=0,1,\mathrm{})`$ are orthorgonal and complete $$d^3\stackrel{}{r}\mathrm{\Phi }_m^{}(\stackrel{}{r})\mathrm{\Phi }_n(\stackrel{}{r})=\delta _{mn}.$$ (7) In most cases, only the lower part of elementary excitations determines the properties of the system near zero temperature, so it is a good approximation to hold only the lower $`f`$-level of eq.(6), i.e. $`n=\{0,1,\mathrm{},f1\}`$. In addition, when $`f\mathrm{}`$, the present treatment becomes rigorous in principle. Substituting eq. (5) into eq. (2), we obtain $$\delta \widehat{\psi }(\stackrel{}{r})=\underset{n=0}{\overset{f1}{}}\widehat{a}_n\mathrm{\Phi }_n(\stackrel{}{r}),$$ (8) where $$\widehat{a}_0=\widehat{A}_0\sqrt{N_0}.$$ (9) In terms of eqs. (8,6), the Hamiltonian (4) can be reexpressed as $`\widehat{H}`$ $``$ $`{\displaystyle \frac{B_{00}}{2}}N_0+{\displaystyle \underset{m,n}{}}A_{mn}\widehat{a}_m^{}\widehat{a}_n`$ (10) $`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{m,n}{}}(B_{mn}^{}\widehat{a}_m\widehat{a}_n+B_{mn}\widehat{a}_m^{}\widehat{a}_n^{}),`$ where $$A_{mn}=(\mu _m\mu _0)\delta _{mn}+d_{mn},$$ (11) $`B_{mn}`$ $`=`$ $`gN_0{\displaystyle d^3\stackrel{}{r}\mathrm{\Phi }_0(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r})\mathrm{\Phi }_m^{}(\stackrel{}{r})\mathrm{\Phi }_n^{}(\stackrel{}{r})},`$ (12) $`d_{mn}`$ $`=`$ $`gN_0{\displaystyle d^3\stackrel{}{r}\mathrm{\Phi }_m^{}(\stackrel{}{r})\mathrm{\Phi }_0^{}(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r})\mathrm{\Phi }_n(\stackrel{}{r})}.`$ (13) The scheme of diagonalizing the Hamiltonian as the form (10) has been extensively studied by Blaizot and Ripka. Here, their results will be briefly reviewed in the following for the use in our present discussion. First, one write the Hamiltonian in a compact form $$\widehat{H}\frac{B_{00}}{2}N_0+\frac{1}{2}\alpha ^{}M\alpha \frac{1}{2}trA$$ (14) by introducing the vector operator $$\alpha ^{}=\left(\begin{array}{cc}\widehat{a}^{}\hfill & \widehat{a}\hfill \end{array}\right)=\left(\begin{array}{cccccccc}\widehat{a}_0^{}\hfill & \widehat{a}_1^{}\hfill & \mathrm{}\hfill & \widehat{a}_{f1}^{}\hfill & \widehat{a}_0\hfill & \widehat{a}_1\hfill & \mathrm{}\hfill & \widehat{a}_{f1}\hfill \end{array}\right),$$ (15) the $`2f\times 2f`$ coefficients matrix $$M=\left(\begin{array}{cc}A\hfill & B\hfill \\ B^{}\hfill & A^{}\hfill \end{array}\right),A=A^{},B=\stackrel{~}{B}.$$ (16) In ref., it is assumed that the matrix $`M`$ is semi-definite positive. However, in our discussion, this condition on the matrix $`M`$ is not necessary since it is only a condition for the stability of the mode of the ground state of BEC. Second, a unitary canonical transformation is carried out to diagonalize the Hamiltonian, $$\beta =T\alpha $$ (17) where the operator vector $$\beta ^{}=\left(\begin{array}{cc}\widehat{b}^{}\hfill & \widehat{b}\hfill \end{array}\right)=\left(\begin{array}{cccccccc}\widehat{b}_0^{}\hfill & \widehat{b}_1^{}\hfill & \mathrm{}\hfill & \widehat{b}_{f1}^{}\hfill & \widehat{b}_0\hfill & \widehat{b}_1\hfill & \mathrm{}\hfill & \widehat{b}_{f1}\hfill \end{array}\right),$$ (18) the transformation matrix and its inverse matrix $$T=\left(\begin{array}{cc}X^{}\hfill & Y^{}\hfill \\ Y\hfill & X\hfill \end{array}\right),T^1=\left(\begin{array}{cc}\stackrel{~}{X}\hfill & Y^{}\hfill \\ \stackrel{~}{Y}\hfill & X^{}\hfill \end{array}\right),$$ (19) the matrix satisfy the canoical condition $$T\eta T^{}\eta =1,$$ (20) and the matrix $$\eta =\left(\begin{array}{cc}1_f\hfill & 0\hfill \\ 0\hfill & 1_f\hfill \end{array}\right).$$ (21) If the elements of the $`f\times f`$ matrixes $`X`$ and $`Y`$ satisfy the following conditions $$\eta MV^n=\varpi _nV^n,$$ (22) where $$V^n=\left(\begin{array}{c}X^n\hfill \\ Y^n\hfill \end{array}\right),X_i^n=X_{ni},Y_i^n=Y_{ni},(i=1,2,\mathrm{},f)$$ especially, if the above equation have one special solution $`\left\{\varpi _0=0,V^0=P\right\}`$, i.e. $$\eta MP=0$$ (23) the Hamiltonian can be written as $$\widehat{H}\frac{B_{00}}{2}N_0+\underset{n=1}{\overset{f1}{}}\omega _n\widehat{b}_n^{}\widehat{b}_n+\frac{\mathrm{}^2}{2\mu }+\frac{1}{2}\underset{n}{}\omega _n\frac{1}{2}trA$$ (24) where $$\mathrm{}\alpha ^{}\eta P,$$ (25) and $`\mu `$ is a positive constant, which can be determined by the following conditions $`\eta MQ`$ $`=`$ $`i{\displaystyle \frac{P}{\mu }},`$ (26) $`Q^{}\eta P`$ $`=`$ $`i,`$ (27) where the vector $`Q`$ is orthorgonal to all the eigenvectors of the matrix $`\eta M`$. In Hamiltonian (24), the second term is the Hamiltonian of a system of independent oscillators, which represent elementary excitations of the system; However, the third term has the form of a free kinetic energy, which is connected with a collective motion in Fock space arising from a broken $`U(1)`$ symmetry in the procedure of the mean field approximation. Usually, the third term is termed with spurious state or zero mode for the corresponding eigenvalue and the norm of the vector $`P`$ both are zero. The physical meaning of this term will be discussed in details in Sec. 3. In fact, we can determine the vector $`P`$ through observation. Notice that in principle, we can solve the the eigenvalues $`\left\{\omega _n\right\}`$ and the corresponding eigenvectors defined by $`\{X_m^n,Y_m^n\}`$ of eqs. (22), which consist of $`2\times f\times f`$ homogeneous linear equations. Obviously, the eigenvalue $`\omega _0=0`$ and the corresponding vector $`P`$ denoted by $`\left\{X_m^0=\delta _{m0},Y_m^0=\delta _{m0}\right\}`$ is a specific solution of the above equations. Hence $$\mathrm{}=\widehat{a}_0+\widehat{a}_0^{},$$ (28) which is in agreement with that in ref.. From eqs. (24,28), a direct conclusion is that the approximate vacuum state $`\|Vac=_i\|Vac_i`$ of BEC satisfies the following conditions $`\widehat{b}_i\|Vac_i`$ $`=`$ $`0,`$ (29) $`\mathrm{}\|Vac_0`$ $`=`$ $`0.`$ (30) To sum up, in this section, we give a new formalism for elementary excitations in trapped BEC, which is equivalent to the standard linearized GP equation. This equivalence can be easily verified when the complex wave functions $`\left(\begin{array}{c}u(\stackrel{}{r})\hfill \\ v(\stackrel{}{r})\hfill \end{array}\right)`$ in traditional method are expanded with the specific complete wave functions defined by eqs.(6). In this sense, our formalism is a specific representation of the traditional method. ## 3 Homogeneous case $`V_{tr}(\stackrel{}{r})=0`$ In this section, the general formalism obtained in the above section will be demonstrated in the homogeneous case. In this case, the ground wave function satisfy $$(\frac{\mathrm{}^2^2}{2M}+gN_0\mathrm{\Phi }_0^{}(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r}))\mathrm{\Phi }_0(\stackrel{}{r})=\mu _0\mathrm{\Phi }_0(\stackrel{}{r}).$$ (31) Therefore, the ground wave function and the chemical potential are given by $$\mathrm{\Phi }_0(\stackrel{}{r})=\frac{1}{\sqrt{V}},\mu _0=g\frac{N_0}{V}.$$ (32) Eq. (6) which give the complete wave functions now becomes $$(\frac{\mathrm{}^2^2}{2M}+\mu _0)\mathrm{\Phi }_\stackrel{}{k}(\stackrel{}{r})=\mu __\stackrel{}{k}\mathrm{\Phi }_\stackrel{}{k}(\stackrel{}{r})$$ (33) By solving the above equations, the eigen wave functions and the corresponding eigen values are given by $`\mathrm{\Phi }__\stackrel{}{k}(\stackrel{}{r})`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{V}}}e^{i\stackrel{}{k}\stackrel{}{r}},`$ $`\mu __\stackrel{}{k}`$ $`=`$ $`\mu _0+{\displaystyle \frac{\mathrm{}^2k^2}{2M}}.`$ (34) Formulas (12,13) become $`B_{\stackrel{}{k}\stackrel{}{k}^{}}`$ $`=`$ $`gN_0{\displaystyle d^3\stackrel{}{r}\mathrm{\Phi }_0(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r})\mathrm{\Phi }_\stackrel{}{k}^{}(\stackrel{}{r})\mathrm{\Phi }_\stackrel{}{k}^{}^{}(\stackrel{}{r})}=g{\displaystyle \frac{N_0}{V}}\delta _{\stackrel{}{k},\stackrel{}{k}^{}},`$ $`d_{\stackrel{}{k}\stackrel{}{k}^{}}`$ $`=`$ $`gN_0{\displaystyle d^3\stackrel{}{r}\mathrm{\Phi }_\stackrel{}{k}^{}(\stackrel{}{r})\mathrm{\Phi }_0^{}(\stackrel{}{r})\mathrm{\Phi }_0(\stackrel{}{r})\mathrm{\Phi }_\stackrel{}{k}^{}(\stackrel{}{r})}=g{\displaystyle \frac{N_0}{V}}\delta _{\stackrel{}{k},\stackrel{}{k}^{}}.`$ (35) Due to the fact the mode denoted by $`\stackrel{}{k}`$ is only coupled to the mode denoted by $`\stackrel{}{k}`$ implied by the above equation, eqs.(22) can be simplified as $`({\displaystyle \frac{\mathrm{}^2k^2}{2M}}+g{\displaystyle \frac{N_0}{V}}\omega _\stackrel{}{k}^{})X_\stackrel{}{k}^\stackrel{}{k}^{}+g{\displaystyle \frac{N_0}{V}}Y_\stackrel{}{k}^\stackrel{}{k}^{}`$ $`=`$ $`0`$ $`({\displaystyle \frac{\mathrm{}^2k^2}{2M}}+g{\displaystyle \frac{N_0}{V}}+\omega _\stackrel{}{k}^{})Y_\stackrel{}{k}^\stackrel{}{k}^{}+g{\displaystyle \frac{N_0}{V}}X_\stackrel{}{k}^\stackrel{}{k}^{}`$ $`=`$ $`0`$ (36) When $`\stackrel{}{k}^{}0`$, the eigenvalue can be calculated by requiring that the above equations have nontrivial solution, $$\omega _\stackrel{}{k}^{}=\omega _\stackrel{}{k}^{}=\sqrt{(\frac{\mathrm{}^2k^{}_{}{}^{}2}{2M}+g\frac{N_0}{V})^2(g\frac{N_0}{V})^2}$$ (37) The corresponding annihilation operators of the elementary excitation are $$\widehat{b}_\stackrel{}{k}^{}=\sqrt{\frac{1}{2}(\frac{\frac{\mathrm{}^2k^{}_{}{}^{}2}{2M}+g\frac{N_0}{V}}{\omega _\stackrel{}{k}^{}}+1)}\widehat{a}_\stackrel{}{k}^{}\sqrt{\frac{1}{2}(\frac{\frac{\mathrm{}^2k^{}_{}{}^{}2}{2M}+g\frac{N_0}{V}}{\omega _\stackrel{}{k}^{}}1)}\widehat{a}_\stackrel{}{k}^{}^{}$$ or $$\widehat{b}_\stackrel{}{k}^{}=\sqrt{\frac{1}{2}(\frac{\frac{\mathrm{}^2k^{}_{}{}^{}2}{2M}+g\frac{N_0}{V}}{\omega _\stackrel{}{k}^{}}+1)}\widehat{a}_\stackrel{}{k}^{}\sqrt{\frac{1}{2}(\frac{\frac{\mathrm{}^2k^{}_{}{}^{}2}{2M}+g\frac{N_0}{V}}{\omega _\stackrel{}{k}^{}}1)}\widehat{a}_\stackrel{}{k}^{}^{}$$ Clearly, it comes back to the familiar form, which supports that our formalism is equivalent to the traditional one. Since the operators $`\widehat{a}_0`$ and $`\widehat{a}_0^{}`$ are only coupled each other, we can limit ourself in the subspace of the wave vector $`\stackrel{}{k}^{}=0`$. The matrix $`\eta M`$ in this subspace is $$\eta M=\left(\begin{array}{cc}g\frac{N_0}{V}\hfill & g\frac{N_0}{V}\hfill \\ g\frac{N_0}{V}\hfill & g\frac{N_0}{V}\hfill \end{array}\right).$$ Obviously, the eigen vector $`P`$ of the zero mode and the corresponding momentum operator $`\mathrm{}`$ are obtained as $$P=\left(\begin{array}{c}1\hfill \\ 1\hfill \end{array}\right),\mathrm{}=\widehat{a}_0+\widehat{a}_0^{}.=(\widehat{A}_0+\widehat{A}_0^{})2\sqrt{N_0}.$$ According to eqs. (26,27), the other independent vector $`Q`$ and the constant $`\mu `$ are obtained as $$Q=\frac{i}{2}\left(\begin{array}{c}1\hfill \\ 1\hfill \end{array}\right),\mu ^1=g\frac{N_0}{V}.$$ In sum, the Hamiltonian can be written as $$\widehat{H}\frac{gN_0^2}{2V}+\underset{\stackrel{}{k}0}{}\omega _\stackrel{}{k}\widehat{b}_\stackrel{}{k}^{}\widehat{b}_\stackrel{}{k}+\frac{\mathrm{}^2}{2\mu }+\frac{1}{2}\underset{\stackrel{}{k}0}{}\omega _\stackrel{}{k}\frac{1}{2}\underset{\stackrel{}{k}}{}(\frac{\mathrm{}^2k^2}{2M}+g\frac{N_0}{V}).$$ (38) Now, the approximate vacuum state $`\|Vac=_\stackrel{}{k}\|Vac_\stackrel{}{k}`$ of BEC can be obtained analytically by solving eqs.(29,30). When the wave vector $`\stackrel{}{k}0`$, the vacuum state $`\|Vac_\stackrel{}{k}`$ is given by solving the equations $`\widehat{b}_\stackrel{}{k}\|Vac_\stackrel{}{k}=\widehat{b}_\stackrel{}{k}\|Vac_\stackrel{}{k}=0`$, $$\|Vac_\stackrel{}{k}=\underset{n}{}A_\stackrel{}{k}\left(\frac{\frac{\mathrm{}^2k^2}{2M}+g\frac{N_0}{V}\omega _\stackrel{}{k}}{\frac{\mathrm{}^2k^2}{2M}+g\frac{N_0}{V}+\omega _\stackrel{}{k}}\right)^{\frac{n}{2}}\|n_\stackrel{}{k}\|n_\stackrel{}{k},$$ where $`A_\stackrel{}{k}`$ is the constant of normalization, the vector $`\|n_\stackrel{}{k}`$ ($`\|n_\stackrel{}{k}`$) is the eigenstate of the number operator $`\widehat{a}_\stackrel{}{k}^{}\widehat{a}_\stackrel{}{k}`$ ($`\widehat{a}_\stackrel{}{k}^{}\widehat{a}_\stackrel{}{k}`$). When the wave vector $`\stackrel{}{k}=0`$, $`\mathrm{}\|Vac_0=0`$. Therefore, $$\|Vac_0=\frac{1}{\sqrt{2\pi }}\underset{n}{}\|n_0𝑑xe^{i\sqrt{2N_0}x}x\|n^{},$$ (39) where the state vector $`\|n_0`$ is the eigen state of the number operator $`\widehat{A}_0^{}\widehat{A}_0`$, $$x\|n=[\frac{1}{\sqrt{\pi }2^nn!}]^{\frac{1}{2}}e^{\frac{1}{2}x^2}H_n(x),$$ and $`H_n(x)`$ is Hermit polynomial. In this section, we explicitly solve the elementary excitations of BEC in the homogeneous case. In this case, it is easy to see that the term of the zero mode in Hamiltonian (38) originate from the quantum fluctuation of the mode denoted by macroscopic wave function $`\mathrm{\Phi }_0(\stackrel{}{r})`$. In fact, if we adopt the usual Bogoliubov approximation, i.e. $`\widehat{A}_0\widehat{A}_0^{}\sqrt{N_0}`$, thus the momentum operator $`\mathrm{}0`$. However, we have no specific reasons to ignore this quantum fluctuation while maintaining those of the other modes. In our treatment, the conservation of the particle number is destroyed, which is easily seen from the approximate vacuum state $`\|Vac`$. The kinetic term appearing in the Hamiltonian is originated from this symmetry breaking, which represents a collective motion, not an intrinsic elementary excitation of the system. ## 4 Conclusion In this paper, based on a natural choice of the complete wave functions, we expand the atomic field operator and obtain a new formalism for the excitations of trapped BEC system near zero temperature. We argue that our formalism is equivalent to the standard linearized GP equation. In terms of this formalism, we illustrate the relation between the zero mode and the other excited modes. Essentially, the zero mode originates from the quantum fluctuations of the mode denoted by the condensate wave function. When applicating the formalism to the homogeneous case, the formalism comes back to the usual Bogoliubov excitation spectrum, which identifies our theory.. Especially, in this case, the physical meaning of zero mode become obvious and the ground state of BEC can be calculate explicitly up to second order of the quantum fluctuations. ACKNOWLEDGMENT: This work is supported by NSF of China.
warning/0004/astro-ph0004033.html	ar5iv	text	# Status of H i searches for CHVCs beyond the Local Group ## 1 Introduction Many aspects of the high–velocity cloud H emission detected in large–area surveys at radial velocites between about $``$450 and +400 km s<sup>-1</sup> can be understood if this emission is associated with a population of low–mass, dark–matter–dominated structures distributed throughout the Local Group. This scenario was rejuvenated in a modern context by Blitz et al. (blit99 (1999)), based on the HVC tabulation of Wakker & van Woerden (wakk91c (1991)) and data of the Leiden/Dwingeloo H Survey (LDS) of Hartmann & Burton (hart97 (1997)), and by Braun & Burton (brau99 (1999)), based on their extraction of compact, isolated HVCs from the LDS and confirmed with independent data. In the Local Group scenario, extended HVC complexes are nearby objects currently merging (or soon to merge) with the Galaxy, while the compact isolated high–velocity clouds (CHVCs) are the distant counterparts of the complexes which have as yet not strongly interacted with the massive Local Group constituents. At least one of the HVC complexes (Chain A) lies at a distance between about 4 and 10 kpc (van Woerden et al. vanw99 (1999)); another (Complex C) has a metallicity of only 0.09 solar (Wakker et al. wakk99 (1999)), implying an extragalactic origin even if the complex is currently nearby. If located at $`D`$15 kpc, the H mass of Complex C is $`10^7`$ M and its linear size is 15 kpc. Further evidence for the extragalactic nature of the CHVCs follows from the metallicity measurement of 0.05 solar for CHVC 125+41$``$207 by Braun & Burton (brau00 (2000)). The derived distance of this same object, which follows from a comparison of the H column and volume densities, is 650$`\pm `$300 kpc. Comparably large distances and high dark–to–visible mass ratios (10 to 50) are also implied by the rotation signatures displayed by many CHVCs (Braun & Burton brau00 (2000)). In the direct vicinity of massive galaxies we expect that these objects will live but shortly before tidal disruption and merger ensue. Such ongoing mergers are seen regularly in nearby galaxies, e.g. LMC (Luks & Rohlfs luks92 (1992)), M33 (Deul & van der Hulst deul87 (1987)), M101 (Kamphuis kamp93 (1993)), NGC 628 (Kamphuis & Briggs kamp92 (1992)). Gas masses of a few times 10<sup>6</sup> to more than 10<sup>8</sup> M and physical dimensions of 10 to 20 kpc are observed. Counterparts to the CHVCs beyond the Local Group are more difficult to constrain. We argue here that current surveys are only beginning to achieve the appropriate combination of sensitivity and spatial coverage. ## 2 CHVC populations, distance, mass and size Where and in what numbers might CHVC populations be expected to occur? Low–mass objects which have so far escaped assimilation by high–mass galaxies might be found in the vicinity of substantial, dynamically–young mass concentrations at distances less than the turn–around radius of the local over–density. In the case of the Local Group, this radius is at about 1.2 Mpc (Courteau & van den Bergh cour99 (1999)) for a total mass of some 2.3$`\times 10^{12}`$ M. The turn–around radius scales as the cube root of mass (Sandage sand86 (1986)). For field galaxies with a typical luminosity of a few $`10^9`$ L and a mass of a few 10<sup>10</sup> M the turn–around radius would lie at about 250 kpc. The total number of associated low–mass objects in a dynamically young system is expected to scale linearly with total mass (Klypin et al. klyp99 (1999)). Relative to the Local Group, we would then expect some 100 times fewer objects associated with a single L galaxy in the field. What is the typical H mass, linewidth and linear size expected for a population of CHVCs? The average integrated flux of the 65 objects in the CHVC sample of Braun & Burton (brau99 (1999)) is 100 with a dispersion of 75 Jy km s<sup>-1</sup>; the average linewidth 30 km s<sup>-1</sup> FWHM; and the average FWHM angular size is 50 with a dispersion of 25 arcmin when imaged with a 36 arcmin beam. Only a small number of distance estimates for CHVCs are currently in hand and even these are indirect, relying on an assumption about the typical thermal pressure of each source. ### 2.1 CHVC distances utilizing compact opaque cores For objects that have compact opaque cores (in which the H brightness temperature in emission is comparable to the gas kinetic temperature) with a column density enhancement, $`dN_{\mathrm{HI}}`$, of angular size, $`\theta `$, it is possible to estimate the distance from $`D_{\mathrm{Core}}=dN_{\mathrm{HI}}^{\mathrm{Core}}/(fn_{\mathrm{HI}}\theta )`$, assuming crude spherical symmetry, if we have an estimate of the H volume density, $`n_{\mathrm{HI}}`$ and its volume filling factor along the line-of-sight, $`f`$. Since the thermal pressure is given by $`P=k_Bn_{\mathrm{HI}}T_k`$, we can rewrite the expression for distance as, $`D_{\mathrm{Core}}={\displaystyle \frac{670}{f}}\left({\displaystyle \frac{dN_{\mathrm{HI}}^{\mathrm{Core}}}{10^{21}\mathrm{cm}^2}}\right)\left({\displaystyle \frac{T_k}{100\mathrm{K}}}\right)`$ $`\left({\displaystyle \frac{P/k}{100\mathrm{cm}^3\mathrm{K}}}\right)^1\left({\displaystyle \frac{\theta }{100^{\prime \prime }}}\right)^1`$ $`\mathrm{kpc}.`$ (1) In general, only a crude lower limit to the distance follows from this approach since the filling factor is difficult to determine and might be substantially less than unity. However, in those cases where the core H is essentially optically thick, a filling factor near unity can be plausibly assumed. This method can perhaps best be demonstrated by using M31 as an example. The excess column density, $`dN_{\mathrm{HI}}^{\mathrm{Core}}=5\pm 1\times 10^{21}`$cm<sup>-2</sup>, and angular size, $`\theta =60\pm 20`$ arcsec, of opaque H clumps in the North-East half of M31 can be estimated from Fig. 5 of Braun & Walterbos (brau92 (1992)). The H kinetic temperature, $`T_k=175\pm 25`$ K, and estimated thermal pressure, $`P/k=1500\pm 500`$ cm<sup>-3</sup>K, of the M31 mid-disk are tabulated in Table 4 of the same reference. This yields the distance estimate: $`D_{M31}=650\pm 220`$ kpc, assuming a volume filling factor for these clumps of unity. While rather crude, this approach gives a plausible value for the distance to M31. Opaque H cores have only been detected in one CHVC to date, namely CHVC125+41$``$207 (Braun & Burton brau00 (2000)), where an excess column density, $`dN_{\mathrm{HI}}^{\mathrm{Core}}=1.0\pm 0.2\times 10^{21}`$cm<sup>-2</sup>, and angular size, $`\theta =90\pm 15`$ arcsec, is measured in the H clumps. A good estimate of the kinetic temperature follows from the measured linewidth in these clumps of only 2.0 km s<sup>-1</sup> FWHM, for which $`T_k=85\pm 10`$ K, while a peak brightness temperature of 75 K is observed, suggesting that a filling factor near unity might be considered. What must still be estimated to employ eqn.1 is the appropriate thermal pressure. Although the pressure, $`P/k`$, in the solar neighborhood of the Galaxy is well-determined at about 2000 cm<sup>-3</sup>K, it is expected to decline rapidly with distance from the Galactic plane (eg. Wolfire et al. wolf95 (1995)) such that beyond about 10 kpc we should encounter values, $`P/k=100`$ cm<sup>-3</sup>K or less due to the Galactic halo. Within the inter-galactic medium (IGM) of the Local Group, it is difficult to determine what values of thermal pressure might be encountered. Current instrumentation has allowed diffuse X-ray emission to be detected from the IGM in several “poor” galaxy groups (eg. Mulchaey et al. mulc96 (1996)). It must be noted that even the poorest of the currently detected groups is substantially richer than the Local Group and all have proven to be lacking in spirals. The X-ray detected systems have derived gas masses of about 10<sup>13</sup>M (comparable to or less than the mass in galaxies), temperatures of about 1 keV (corresponding to 10<sup>7</sup> K) and radii of about 500 kpc. Mulchaey et al. suggest that spiral-dominated galaxy groups are likely to have lower X-ray temperatures, of perhaps 0.2 keV. Assuming a correspondingly lower total gas mass of about 10<sup>12</sup>M and radii between 500 and 1000 kpc such as might apply to the Local Group, would imply a Local Group IGM thermal pressures in the range 160 to 20 cm<sup>-3</sup>K. Braun & Burton (brau00 (2000)) have considered the total pressure implied by assuming hydrostatic equilibrium of a self-gravitating gaseous disk. The self-gravity of the WNM in a diffuse disk system accounts for a contribution, $`P/k=40(N_{\mathrm{HI}}^{\mathrm{Halo}}(0)/10^{20}`$cm<sup>-2</sup>)<sup>2</sup> cm<sup>-3</sup>K in the absence of a significant dark matter contribution to the disk mass. The peak column densities of WNM found in CHVCs are in the range 10<sup>19.5</sup>–10<sup>20</sup>cm<sup>-2</sup> (Burton et al burt00 (2001)). The resulting thermal pressure, the sum of self-gravitating and external components, at the interface of the warm neutral halo with the cool condensed cores, of about $`P/k=100\pm 50`$ cm<sup>-3</sup>K is in very good agreement with that expected theoretically (Wolfire et al. wolf95 (1995). This transition pressure is found to apply over a wide range of physical conditions (metal abundance, radiation field and dust content). Taking the value, $`P/k=100\pm 50`$ cm<sup>-3</sup>K, allows calculation of the distance to CHVC125+41$``$207, yielding $`D_{\mathrm{C125}}=600\pm 300`$ kpc. ### 2.2 CHVC distances utilizing edge profiles A second method of distance estimation (Burton et al. burt00 (2001)) makes use of the column density profiles of the diffuse WNM halos of these objects. These halos can be well-described by the sky-plane projection of a spherical exponential distribution of the H volume density, $$n_{\mathrm{HI}}(r)=n_\mathrm{o}e^{r/h}$$ (2) in terms of the radial distance, $`r`$, and exponential scale length, $`h`$. The projected column density distribution is then, $$N_{\mathrm{HI}}^{\mathrm{Halo}}(r)=2hn_\mathrm{o}\left[\frac{r}{h}K_1\left(\frac{r}{h}\right)\right]$$ (3) where $`K_1`$ is the modified Bessel function of order 1. In this formulation, the peak halo column density is given simply by, $$N_{\mathrm{HI}}^{Halo}(0)=2hn_\mathrm{o}$$ (4) The distance of the source can then be calculated from, $$D_{\mathrm{Halo}}=\frac{N_{\mathrm{HI}}^{\mathrm{Halo}}(0)}{2hn_o}=\frac{N_{\mathrm{HI}}^{\mathrm{Halo}}(0)k_BT_k^{\mathrm{Halo}}}{2hP}$$ (5) where the implicit assumption has been made of a unit filling factor for the halo gas. If we assign the measured temperature of $`T_k^{\mathrm{Halo}}=10^4`$ K to the halo gas this becomes, $`D_{\mathrm{Halo}}=335\left({\displaystyle \frac{N_{\mathrm{HI}}^{\mathrm{Halo}}(0)}{10^{19}\mathrm{cm}^2}}\right)`$ $`\left({\displaystyle \frac{P/k}{100\mathrm{cm}^3\mathrm{K}}}\right)^1\left({\displaystyle \frac{h}{100^{\prime \prime }}}\right)^1`$ $`\mathrm{kpc}.`$ (6) The peak halo column densities, $`N_{\mathrm{HI}}^{\mathrm{Halo}}(0)`$, and scale-lengths, $`h`$, can be determined from sensitive total power observations, like the deep Arecibo cross-cuts obtained for a sample of ten objects by Burton et al. (burt00 (2001)). Assuming, as before, a nominal thermal pressure of $`P/k=100\pm 50`$ cm<sup>-3</sup>K at the interface of the cool and warm neutral H , yields distance estimates which range from 150 to 850 kpc for the various objects studied. An independent method of distance estimation (Burton et al. burt00 (2001)) follows from simply equating the angular H scale-lengths observed in the CHVCs with the mean physcial scale-length, $`h=1.1\pm 0.4`$ kpc, measured in a sample of low mass dwarf galaxies. This yields distances in the range 320 to 730 kpc. Apparently, a wide range of distances, spanning almost an order of magnitude, is indicated for the Local Group CHVC population. If a distance range of 150$``$850 kpc is representative for the CHVCs, then the corresponding H mass range is $`M_{\mathrm{HI}}=10^{5.7}10^{7.2}`$ M. The distribution of masses over this range is not yet known, but is likely to be more heavily populated at the low end than the high. The critical H mass detection threshold is therefore about $`10^6`$ M, if a large fraction of the population is to be detected. If the typical radial profile of the CHVCs were Gaussian, then the corresponding range of FWHM linear sizes would be 2$``$12 kpc. Since the edge profiles are instead observed to be exponential, these objects will appear slightly resolved when observed with any telescope beam broader than that of the high column density cores of about 5–10 arcmin (Braun & Burton brau00 (2000)), corresponding to 1–2 kpc. Unfortunately, with an exponential edge profile, the observed object “size” and “total mass” will depend not only on the resolution used, but also on the sensitivity achieved. For a 2-D circular exponential distribution, the fractional flux contained within a radius of one, two and three scale-lengths is 26, 59 and 80%. The implication is that only for beam sizes greater than about 20 kpc, is it likely that the entire flux of an object will be detected in cases of modest signal-to-noise ratio. ## 3 Synthesis Surveys One method of achieving “blind” survey coverage for low mass objects is to utilize the relatively large field-of-view of an interferometeric array when observing the H distribution in nearby galaxies to obtain serendipitous detection of nearby uncataloged systems. Some aspects of this problem have been considered previously by Blitz et al. (blit99 (1999)). With a typical interferometric field of view of 30 arcmin, an area with linear size of 90 $`D_{10}`$ kpc (normalized to a source at 10 Mpc) is probed with a single pointing. Since correlator capacity is often a limiting factor for such observations, the typical velocity coverage obtained is often only about 1.5 times the velocity extent of the target galaxy. Depending on the goals of a particular observing program, this entire field and velocity coverage may not even be analyzed; but if it were, the probability of intercepting an associated object would be proportional to the total number of objects, $`N_{\mathrm{Grp}}`$ for a galaxy group or $`N_{\mathrm{Gal}}`$ for an isolated galaxy, as well as the fractional area and velocity extent observed. Assuming that the velocity coverage is sufficient to sample the entire associated population, this corresponds to $`N=0.002N_{\mathrm{Grp}}D_{10}^2`$ for a distribution extending out to 1 Mpc (a distance comparable to the Local Group turn-around radius) or $`N=0.03N_{\mathrm{Gal}}D_{10}^2`$ for an isolated galaxy. Only for values of $`N_{\mathrm{Grp}}`$ as large as 500 or $`N_{\mathrm{Gal}}`$ = 30 would one expect serendipitous detections in individual pointings that achieve a detected mass sensitivity of better than about $`10^6`$ M over 30 km s<sup>-1</sup> with a beam size larger than about 20 kpc. The required mass sensitivity for CHVC detection in external galaxy groups has not yet been generally attained. Taylor et al. (tayl95 (1995),tayl96 (1996)) have searched the fields of 21 apparently isolated H ii galaxies and of 17 low surface brightness dwarf galaxies. The targeted galaxies were at a typical distance of about 25 Mpc, where the half-power field-of-view is 225 kpc and the spatial resolution about 8 kpc. The observations reached an average five sigma mass sensitivity over 30 km s<sup>-1</sup> of about $`2.5\times 10^7`$ M at the field center in the case of the H ii galaxies and $`1.7\times 10^7`$ M in the case of the LSB dwarf galaxies. While a total of about 17 companion galaxies were discovered in this survey, all but one (UM422C) were found to have optical counterparts detectable at surface brightnesses of 23 mag arcsec<sup>-2</sup>. Since the H detection thresholds have very little overlap with the mass range of the Local Group CHVC population this result can be readily understood. Another sensitive search for H emission was carried out by Van Gorkom et al. (vang93 (1993)) in the vicinity of Ly$`\alpha `$ clouds in a single $`30^{}`$ FWHM field centered on 3C 273 and covering $`+840`$ to $`+1840`$ km s<sup>-1</sup>. The field is located some $`10^{}`$ (or about 2.6 Mpc) from the center of the Virgo Cluster (at an assumed distance of 15 Mpc). The rms noise in the center of their field corresponded to a $`5\sigma `$ mass sensitivity of $`2.5`$$`4.5\times 10^6`$M, in a single 41.6km s<sup>-1</sup> channel at linear resolutions of 1.5–4.2 kpc. Van Gorkom et al. estimate a $`5\sigma `$ mass sensitivity at 10 kpc resolution (to better match the expected CHVC size) of $`1.1\times 10^7`$M at the field center. Spatially averaged over the $`30^{}`$ FWHM Gaussian field-of-view yields a $`5\sigma `$ average mass limit of $`1.6\times 10^7`$M over a region 0.013 Mpc<sup>2</sup> in extent. The fractional area sampled relative to the impact parameter of 2.6 Mpc is 0.06%. No H emission was detected, implying an upper limit of some 2000 uniformly distributed objects with $`M_{\mathrm{HI}}>1.6\times 10^7`$ M, if the imaged field is representative. We note that the Virgo Cluster is a dynamically evolved system, for which the long term survival of a low mass population is uncertain. But, in any case, the average mass sensitivity of this search is insufficient to detect large numbers CHVC counterparts. More recently Pisano & Wilcots (pisa99 (1999),pisa00 (2000)) have surveyed the fields of 34 isolated galaxies with distances in the range 21 to 45 Mpc. With a linear resolution of 6–13 kpc, they achieve a 5$`\sigma `$ mass sensitivity over 30 km s<sup>-1</sup> of $`1.4\pm 0.6\times 10^7`$ M at each field center. This sensitivity begins to overlap with the range of expected CHVC masses near the field center. However, averaged over the primary beam out to the half power point, corresponding to radii of 90–190 kpc, the mass sensitivity is decreased to $`2\pm 1\times 10^7`$ M, where effectively no detections are expected. Although several gas-rich companions are detected in this survey, all but possibly one are accompanied by optical counterparts. ## 4 Total Power Surveys One of the earliest systematic searches for H in galaxy groups was that of Lo & Sargent (lo79 (1979)). Large regions of the sky in the vicinity of three nearby galaxy groups were searched using the OVRO 40 m telescope for H emission features without obvious optical counterparts. The 5$`\sigma `$ mass sensitivity for a 35 km s<sup>-1</sup> FWHM linewidth was $`4\times 10^7`$ M for a region of 132 square degrees (or 0.4 Mpc<sup>2</sup>) adjacent to M81, while mass limits of $`8\times 10^7`$ over 0.7 Mpc<sup>2</sup> and $`46\times 10^7`$ M over 5.7 Mpc<sup>2</sup> were obtained in 90 square degree regions near the CVnI group and NGC 1023 respectively. A total of six gas-rich objects were subsequently detected in follow-up work using the 100 m Effelsberg telescope, although all of them do have (faint) optical counterparts. Comparison with our mass estimates above indicates that the non-detection of an extensive CHVC population is entirely in keeping with the likely mass distribution of Local Group CHVCs. ### 4.1 The Arecibo H Strip Survey A sensitive, large area H survey has been the Arecibo H Strip Survey (AHISS) by Zwaan et al. (zwaa97 (1997)). In fact, Zwaan & Briggs (zwaa00 (2000)) claim that the AHISS would have made numerous detections of CHVCs if they did populate the outer environments of galaxies and galaxy groups. In order to critically assess this claim, we have carefully considered the attributes of the survey which are well-documented by Sorar (sora94 (1994)). The spatial coverage of the AHISS consists of two repeatedly–observed strips, one covering the right ascension interval $`18^\mathrm{h}`$ to $`05^\mathrm{h}`$ at $`\delta =14^{}14^{}`$, and the other covering $`01^\mathrm{h}`$ to $`10.^\mathrm{h}7`$ at $`\delta =23^{}09^{}`$. The velocity coverage extends from $`700`$ to +7500 km s<sup>-1</sup>. The noise level of the AHISS data depends on position as well as velocity for a number of reasons. Firstly, varying integration times were obtained for different positions. Secondly, the receiver gain varied significantly with frequency. For the $`\delta =14^{}14^{}`$ strip this implied a factor of two degraded sensitivity at 0 km s<sup>-1</sup> than at +5000 km s<sup>-1</sup>, while for the $`\delta =23^{}09^{}`$ strip, the degradation at 0 km s<sup>-1</sup> was limited to about 15%. And lastly, confusion and spectral baseline residuals due to bright, extended emission from the Galaxy preclude detection of objects within about $`\|V_{GSR}\|<200`$ km s<sup>-1</sup>. The average rms level of the $`\delta =14^{}`$ strip is 0.95 mJy beam<sup>-1</sup> at 16 km s<sup>-1</sup> resolution at velocities near +5000 km s<sup>-1</sup>, rising to 1.9 mJy beam<sup>-1</sup> at low recession velocites. For the $`\delta =23^{}`$ strip, the average noise level is 0.75 mJy beam<sup>-1</sup> at 16 km s<sup>-1</sup> resolution between about +5000 and +1000 km s<sup>-1</sup>, rising to 0.85 mJy beam<sup>-1</sup> at low recession velocities. The spatial coverage of the AHISS is given by the track of the telescope beam (about 180 arcsec FWHM at 1400 MHz) along the indicated right ascension intervals. The maximum sensitivity, noted above, is only obtained for completely unresolved sources which pass within the central 20 arcsec or so of the beam response. If instead we consider a realistic search area comparable in width to the beam FWHM, then the average detected flux is down by a factor of 0.81 for an approximately Gaussian main lobe. This diminished response must be taken into account when assessing the sensitivity for unresolved sources over the full 180 arcsec width. If instead, a source is significantly resolved by the survey beam, then the fraction of detected flux at a given offset position is given by the convolution of the beam and source shapes. In the case of two circular Gaussians of similar FWHM, with random linear offsets in the range $`\pm `$HWHM, the average detected flux is reduced by a factor of 0.64 from the intrinsic flux. The mass sensitivity of the AHISS averaged over the FWHM of the telescope beam can now be assessed as a function of recession velocity. For nearby systems, say at 5 Mpc distance, the telescope beam corresponds to only 4.4 kpc and any detected system will be strongly resolved. The rms noise averaged over the $`\delta =14^{}`$ and $`23^{}`$ strips at V<sub>GSR</sub> = 300 km s<sup>-1</sup>, is 1.4 mJy beam<sup>-1</sup> over 16 km s<sup>-1</sup>, yielding a 5$`\sigma `$ limit on detected signal strength of 160 mJy beam<sup>-1</sup>-km s<sup>-1</sup> over the typical 32 km s<sup>-1</sup> linewidth. The average intrinsic signal strength that gives this response will be at least a factor of 1/0.64 higher (due to offsets from the beam center) yielding a 5$`\sigma `$ mass limit of $`1.5\times 10^6`$ M. At this distance the AHISS is indeed capable of detecting a large part of the expected CHVC mass distribution, even with some further reduction in detected flux due to the exponential edge profiles. At a distance of 15 Mpc, the telescope beam coresponds to 13 kpc, which is comparable to the expected source size. The rms noise averaged over the $`\delta =14^{}`$ and $`23^{}`$ strips at V<sub>GSR</sub> = 1000 km s<sup>-1</sup>, is 1.1 mJy beam<sup>-1</sup> over 16 km s<sup>-1</sup>, yielding a 5$`\sigma `$ limit on detected signal strength of 125 mJy beam<sup>-1</sup>-km s<sup>-1</sup> over the typical 32 km s<sup>-1</sup> linewidth. The average intrinsic signal strength will again be higher by the beam averaging factor of 1/0.64 (for the comparable beam and source size case), yielding a 5$`\sigma `$ mass limit of $`1.0\times 10^7`$ M. This constitutes a practical distance limit to the AHISS for expected detection of the upper mass end of the CHVC mass distribution. We searched the Garcia (garc93 (1993)) catalog of galaxy groups for all those in the range +200 to +1200 km s<sup>-1</sup> and within 1 Mpc projected distance of the AHISS coverage. Only one cataloged galaxy group is intercepted by the AHISS in this range, namely the NGC 628 group centered at $`(\alpha ,\delta ,V)=(01^\mathrm{h}36^\mathrm{m},14^{}55^{},+570`$ km s<sup>-1</sup>), at a distance of about 10 Mpc. There is no other instance of the AHISS coverage passing near a group of galaxies within 15 Mpc. The AHISS sensitivity in the NGC 628 region is 1.4 mJy beam<sup>-1</sup>over 16 km s<sup>-1</sup>, while the telescope beam corresponds to about 9 kpc. With a beam averaging factor of 1/0.64, this yields a 5$`\sigma `$ mass sensitivity of $`5.7\times 10^6`$ M over the expected linewidth of 32 km s<sup>-1</sup>. This sensitivity should be sufficient to allow detection of the more massive members of a CHVC population, if they happened to lie within the survey strip. The AHISS coverage passes within 120 kpc of the nominal NGC 628 group center, so the group area intercepted by this strip corresponds to a region about 9 kpc wide (for the 180 arcsec beam FWHM at $`D=10`$ Mpc) and some 2 Mpc long. This area corresponds to some 0.6% of the total group area ($`\pi R^2`$) of about 3.1 Mpc<sup>2</sup>. For a uniform distribution of objects, the AHISS non–detection implies that the total number of such objects exceeding $`5.7\times 10^6`$ M should be less than about 300. The next nearest group covered by the AHISS survey is the NGC 3227 group at $`(\alpha ,\delta ,V)=(10^\mathrm{h}19^\mathrm{m},21^{}23^{},+1250`$ km s<sup>-1</sup>). Since this is covered in the $`\delta =23^{}`$ strip at a velocity greater than +1000 km s<sup>-1</sup>, the flux sensitivity is 0.75 mJy beam<sup>-1</sup>at 16 km s<sup>-1</sup> resolution. The 5$`\sigma `$ mass sensitivity at an estimated distance of 19.2 Mpc is then $`M_{\mathrm{HI}}=1.1\times 10^7`$ M over 32 km s<sup>-1</sup>. This sensitivity is sufficient to detect only the most massive objects in a CHVC population. Since the survey coverage passes through the group at an impact parameter of about 590 kpc, the 17 kpc beamwidth samples a fractional group area of about 0.9%. The AHISS non-detection implies that the number of objects exceeding $`1.1\times 10^7`$ M should be less than about 200. We have also searched the LEDA database for all cataloged galaxies with recession velocities (corrected for Virgo–centric infall) in the range +200 to +1300 km s<sup>-1</sup> and projected separations of less than 250 kpc from the AHISS strip coverage. Apart from galaxies belonging to the NGC 628 group, and therefore already considered above, a total of 7 other galaxies lie this close to the AHISS coverage. Six of the seven are low mass dwarf irregular galaxies (PGC 86669, 169957, 169968, 169969, and UGC 1561, 2905) with H line widths in the range 30–60 km s<sup>-1</sup> FWHM. Three of these sources were in fact discovered in the “Arecibo Slice” survey discussed below. The seventh is a low luminosity spiral (UGC 5672, $`M_B=15.4`$, M$`{}_{HI}{}^{}=5\times 10^7`$ M). None of these galaxies appear to be fruitful starting points for a search for associated low-mass satellite systems, since even the most massive is only about 0.01$`L^{}`$. Since turn-around radius scales as the cube root of total mass we expect turn-around radii of 55 kpc at most for these systems, rather than the nominal value of 250 kpc for an $`L^{}`$ galaxy. Exactly one of the dwarf irregulars lies within 55 kpc of the AHISS track, PGC 169968. At a distance of 8 Mpc, the AHISS 5$`\sigma `$ mass limit is $`2.2\times 10^6`$ M over 32 km s<sup>-1</sup> and the FWHM beam is some 7 kpc broad. At an impact parameter of 17 kpc from an assumed 55 kpc radius distribution, the strip samples a region of length 100 kpc. The fractional area sampled by the AHISS coverage is then 18%. The absence of detected companions for this object in the AHISS implies that its CHVC population should number less than about 10 for objects with mass greater than $`2.2\times 10^6`$ M. ### 4.2 The “Arecibo Slice” Survey Another recent H survey is that of Spitzak & Schneider (spit98 (1998)) who utilized the Arecibo telescope to survey the region between right ascension $`22^\mathrm{h}`$ to $`03^\mathrm{h}24^\mathrm{m}`$ and $`\delta =22^{}58^{}`$ to $`\delta =23^{}47^{}`$. The flux sensitivity at the hexagonally spaced sampling points was typically 1.7 mJy beam<sup>-1</sup> over 16 km s<sup>-1</sup> at V$`{}_{\mathrm{Hel}}{}^{}=5300`$ km s<sup>-1</sup>, decreasing to 81% of this value at 8340 km s<sup>-1</sup>and to 56% at 100 km s<sup>-1</sup>. Since the sampling points were spaced quite coarsely relative to the beam size (4.1 arcmin spacing and a 3.3 arcmin beam), the flux sensitivity is estimated to vary by about a factor of three, depending on whether sources were located at sampling points or exactly between them. The corresponding average 5$`\sigma `$ mass sensitivity at a distance of 5 Mpc is about $`4\times 10^6`$ M, increasing to $`1.5\times 10^7`$ M at 10 Mpc. We searched for galaxy groups in the Garcia (garc93 (1993)) catalog in the velocity range +200 to +1200 km s<sup>-1</sup> within a projected seperation of 1 Mpc of the survey coverage, but found none. The only four LEDA cataloged galaxies in this velocity range which lie inside the survey coverage or within a projected turn-around radius are PGC 169957, 169968, 169969 and UGC 1561. These are all very low mass systems as already noted above in connection with the AHISS survey discussion, where the same galaxies were also encountered. The first three of these objects were discovered in the “Arecibo Slice” survey. Faint optical counterparts are seen for the PGC 169968 and 169957, which have H masses of about 1.3 and 3$`\times 10^7`$ M respectively. PGC 169969 is the lowest H mass object found, in this and any other blind H survey to date, of only about $`8\times 10^6`$ M. No optical counterpart to this source has yet been detected, although a relatively bright fore-ground star makes this a difficult undertaking. These objects fit into the more general trend implied by Fig. 9 of Spitzak & Schneider (spit98 (1998)) for H selected objects to be increasingly gas-dominated at low galaxy mass. ### 4.3 HIPASS survey of the Cen A galaxy group A portion of the Cen A galaxy group has recently been subjected to a similar search by Banks et al. (bank99 (1999)) using the HIPASS data obtained with the Parkes 64 m telescope. They searched an area of 600 square degrees, corresponding to about 2.2 Mpc<sup>2</sup>, for H detections in the velocity range V<sub>Hel</sub> = 200 – 1000 km s<sup>-1</sup>. The 5$`\sigma `$ mass sensitivity over 35 km s<sup>-1</sup> was $`7\times 10^6`$ M for a nominal distance of 3.5 Mpc, while the 15 arcmin angular resolution corresponds to 15 kpc. This mass sensitivity has substantial overlap with the mass distribution we derive for the Local Group CHVC population, although a modest loss in total flux might be expected in the 15 kpc beam. The Cen A group has six major members within a total extent of about 1.9 Mpc, making it about twice as rich as the Local Group. Banks et al. detect 10 new group members in their survey area, most of which have M<sub>HI</sub> near the $`10^7`$ M sensitivity limit. Five of these were not previously cataloged, despite deep optical searches in this region, and have extremely low optical luminosities (below M$`{}_{B}{}^{}15`$) and surface brightness ($`<\mu _B>26`$ mag arcsec<sup>2</sup>). These five objects have an average H mass of $`2.0\pm 0.8\times 10^7`$ M and a FWHM linewidth of 34$`\pm `$4 km s<sup>-1</sup>. While this is still at the high mass end of the Local Group CHVC distribution, it is possible that a new population of low mass, extremely star-poor systems is starting to emerge. An even more extreme object in this class seems to the possible Cen A group member HIPASS J1712$``$64 discovered by Kilborn et al. (kilb00 (2000)). With V<sub>Hel</sub> = 460 km s<sup>-1</sup> it is well-removed in velocity from any confusing emission and in agreement with that of the group. At an estimated distance of 3.2 Mpc, it has an H mass of $`1.7\times 10^7`$ M, 15 kpc size (at 10<sup>19</sup>cm<sup>-2</sup> column density) and lies at a projected seperation of 1.1 Mpc from the massive Circinus galaxy. No optical counterpart to this object has been detected down to a limiting surface brightness of $`\mu _B27`$ mag arcsec<sup>2</sup>. ### 4.4 Targeted Survey of Groups A targeted survey of galaxy group environments was recently carried out by Zwaan (zwaa00b (2000)). A sparsely sampled grid of 60 pointings was observed with the Arecibo telescope centered on each of five galaxy groups at recession velocities of 1770, 1870, 2280, 2910 and 3000 km s<sup>-1</sup> (corresponding to distances between about 27 and 46 Mpc). The 60 pointings were distributed over rectangular areas of about 2.5$`\times `$1.5 Mpc. At these distances each FWHM beam has an area of 24–40 kpc. Together, each grid of 60 pointings covered an area varying from 0.027–0.075 Mpc<sup>2</sup>, corresponding to a fractional coverage of between 0.7–2.0% for the various groups sampled. An average flux sensitivity of 0.75 mJy beam<sup>-1</sup> over 10 km s<sup>-1</sup> was obtained, corresponding to 5$`\sigma `$ mass limits 1.1–3.2$`\times 10^7`$ M over 30 km s<sup>-1</sup> at each beam center. Since the survey area is sparsely sampled we must consider the sensitivity averaged over the beam area, rather than simply taking the peak value. The average 5$`\sigma `$ mass sensitivity over each beam FWHM then varies between 1.6–4.6$`\times 10^7`$ M over 30 km s<sup>-1</sup>. Only in the case of the nearest two galaxy groups is there a small degree of overlap of the achieved mass sensitivity with the expected CHVC mass distribution. The constraint that follows from non-detections over 0.75% of the group area, is that there must be less than about 100 objects with H mass exceeding 1.7$`\times 10^7`$ M associated with the NGC 5798 and NGC 5962 galaxy groups. Unfortunately this is not a very strong constraint. ## 5 Comparison with other work Zwaan & Briggs (zwaa00 (2000)) discuss detection statistics of primordial gas clouds near galaxies and galaxy groups based on the AHISS data. They state that the survey “… probed the haloes of $``$300 cataloged galaxies and the environments of $``$14 groups with sensitivity to neutral hydrogen masses $`10^7`$ M.” As we have shown, a limiting H mass of $`10^7`$ M is indeed a minimal requirement for the expected detection of such features, and even this will only allow detection of the most massive objects in the expected distribution. The AHISS survey attributes (Sorar sora94 (1994)) lead to a distance limit of 15 Mpc for the 5$`\sigma `$ detection of an H mass of $`10^7`$ M within the 15 kpc Arecibo beam over the expected linewidth of 32 km s<sup>-1</sup>. Out to this distance the AHISS survey coverage passes within the turn-around radius of exactly one galaxy group (the NGC 628 group) and one cataloged galaxy (PGC 169968). Relaxing this distance limit to 20 Mpc, brings only one additional galaxy group (the NGC 3227 group) and no additional galaxies into range. There are several reasons for this apparent discrepancy. The single most important factor is the assumed distance, and hence mass, of the population being sought. Assuming a baryon-to-dark-mass fraction, $`f`$ = 0.1, and calculating a stability distance for the entire Wakker & van Woerden (wakk91c (1991)) HVC sample following Blitz et al. (blit99 (1999)) leads to a typical distance of 1 Mpc and H masses in the range $`10^7`$$`10^{8.5}`$ M. We have shown above that current distance estimates to individual CHVCs in the Local Group lie in the range 150–850 kpc and that the corresponding H mass range for CHVCs cataloged in BB99 is $`M_{\mathrm{HI}}=10^{5.7}10^{7.2}`$ M. Zwaan & Briggs also assume a nominal 5$`\sigma `$ column density limit of 10<sup>18</sup>cm<sup>-2</sup>, corresponding to a flux density rms of 0.75 mJy beam<sup>-1</sup>over 16 km s<sup>-1</sup> that applies uniformly to a 180 arcsec wide strip, in assessing the survey detectability of an extragalactic CHVC population. The AHISS flux density rms is in fact typically higher than 0.75 mJy beam<sup>-1</sup>over 16 km s<sup>-1</sup> as discussed specifically above, especially at low velocities. In addition, the assumed linewidth of 16 km s<sup>-1</sup> is about a factor of two lower than the measured linewidths of this class of object, making actual detection more difficult by a factor of $`\sqrt{2}`$. Furthermore, the average detected response across the 180 arcsec survey strip is reduced by about 50% for sources which are comparable in size to the beam. Together these factors degrade the effective detection sensitivity by more than a factor of three. And finally, Zwaan & Briggs have assumed that each cataloged field galaxy, independant of mass, might host a comparably large population of associated satellites extending out to a radius of 1 Mpc. In fact, it is reasonable to expect (eg. Klypin et al. klyp99 (1999)) that the number of associated low mass satellites might be proportional to the host mass and that the radius over which they might extend be the turn-around radius of the host mass (which scales as $`M^{1/3}`$). ## 6 Conclusions Recent distance estimates for a sample of ten compact high-velocity clouds (CHVCs) yield values in the range 150 to 850 kpc (Burton et al. burt00 (2001)). The corresponding H masses of the BB99 catalog of CHVCs range between $`10^{5.7}10^{7.2}`$ M over an average linewidth of 30 km s<sup>-1</sup>. Interferometric surveys of isolated fields have not generally achieved the combination of sensitivity and sky coverage to put strong constraints on CHVC populations in galaxy groups. The searches for low mass companions reported earlier by Taylor et al. (tayl95 (1995), tayl96 (1996)) and more recently by Pisano & Wilcots (pisa99 (1999), pisa00 (2000)) fall short of probing the expected CHVC mass range. The most sensitive large–area H surveys are just beginning to sample the upper mass end of the implied CHVC mass distribution. The AHISS survey (Zwaan et al. zwaa97 (1997)) could have detected objects with $`M_{\mathrm{HI}}=10^7`$ M out to a maximum distance of 15 Mpc. Unfortunately, out to this distance, the AHISS sky coverage passes within the turn-around radius of only a single galaxy group and a single isolated galaxy. The AHISS non-detections allow placing an upper limit of about 300 objects exceeding $`M_{\mathrm{HI}}=5.7\times 10^6`$ M within a radius of 1 Mpc of the NGC 628 galaxy group. A targeted search of the environments of several galaxy groups (Zwaan zwaa00b (2000)) also falls short of reaching the required mass sensitivity. The Arecibo Slice survey (Spitzak & Schneider spit98 (1998)) has detected several objects which overlap in H mass with the CHVCs. The least massive of these ($`M_{\mathrm{HI}}8\times 10^6`$ M) has no detected optical counterpart, in keeping with a more general trend for low mass systems to be increasingly gas-dominated. The HIPASS survey has allowed detection of a small number of low mass galaxies in the Cen A galaxy group with extremely low surface brightness optical counterparts ($`<\mu _B>26`$ mag arcsec<sup>2</sup>). This was accomplished with a 5$`\sigma `$ H mass sensitivity of about $`7\times 10^6`$ M (Banks et al. bank99 (1999)). In addition, at least one object has been found with no detected optical counterpart down to $`\mu _B27`$ mag arcsec<sup>2</sup> (Kilborn et al. (kilb00 (2000)). These objects appear to be the highest mass counterparts of the Local Group CHVC population. If the estimated distances of the Local Group CHVC population discussed in §2 are correct then pushing the detection limits down by another order of magnitude, to $`M_{\mathrm{HI}}=10^6`$ M (over 15 kpc and 30 km s<sup>-1</sup>) will enable detection of populations comparable to those determined for the Local Group, implying increased numbers by some two orders of magnitude. If these low mass populations are not found, then the Local Group hypothesis for the CHVCs must be seriously reconsidered. ###### Acknowledgements. We are grateful to Martin Zwaan for providing details of the statistical analysis presented by Zwaan and Briggs (zwaa00 (2000)). We acknowledge extensive use of the LEDA database, www-obs.univ-lyon1.fr. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration.
warning/0004/nlin0004031.html	ar5iv	text	# 1 Introduction ## 1 Introduction A number of nonlinear optical phenomena in Kerr-like media are described by the perturbed nonlinear Schrödinger equation (NLSE) which in dimensionless variables reads $$i\frac{u}{t}+\frac{1}{2}\frac{^2u}{x^2}+\|u\|^2u(x,t)=iR[u],$$ (1) and by its multicomponent generalizations. Below we shall discuss several perturbations $`R[u]`$ due to possible linear losses, bandwidth limited and nonlinear amplifications. For $`R[u]=0`$ the NLSE (1) can be solved by means of the inverse scattering method applied to the Zakharov-Shabat system $`L[u]`$. The analytical methods developed allow one: i) to prove that any solution of the NLSE in the limit $`t\mathrm{}`$ tends to a purely solitonic solution; ii) describe the soliton interaction in the generic case when all solitons have pair-wise different velocities. For practical applications one needs to describe the behavior of the so-called $`N`$-soliton trains which are solutions of (1) satisfying the initial condition $$u(x,0)=\underset{k=1}{\overset{N}{}}u_k^{1\mathrm{s}}(x,t=0).$$ (2) Here $`u_k^{1\mathrm{s}}(x,t)`$ is the one-soliton solution of (1): $`u_k^{1\mathrm{s}}(x,t)=2\nu _k\mathrm{e}^{i\varphi _k}\text{sech }(z_k(x,t)),z_k(x,t)=2\nu _k(x\xi _k(t)),`$ $`\varphi _k(x,t)=2\mu _k(x\xi _k(t))+\delta _k(t),\xi _k(t)=2\mu _kt+\xi _{k,0},`$ $`\delta _k(t)=2(\mu _k^2+\nu _k^2)t+\delta _{k,0},`$ (3) where $`\nu _k`$, $`\mu _k`$, $`\xi _k`$ and $`\delta _k`$ are the amplitude, velocity, position and phase of the $`k`$-th soliton-like pulse. Let us first remark that for $`N2`$ the parameters $`\mu _k`$ and $`\nu _k`$ are not directly related to the discrete spectrum of $`L[u]`$. In fact the spectral data of $`L[u]`$ with $`u`$ provided by (2) contains not only $`2N`$ discrete eigenvalues $`\lambda _k^\pm =\kappa _k\pm i\eta _k`$, $`\eta _k>0`$, $`k=1,\mathrm{},N`$, but also nonvanishing ‘radiation’ related to the continuous spectrum of $`L[u]`$. However if we take well separated pulses $`\|\xi _{k+1,0}\xi _{k,0}\|r_01`$ then the energy of the ‘radiation’ is of the order of 1% of the total energy and may well be neglected. As a result the corresponding $`N`$-soliton train may be approximated by an $`N`$-soliton solution whose interactions in the generic case (pair-wise different velocities) are well known . Even if we approximate $`u(x,t)`$ by an exact $`N`$-soliton solution it is not an easy matter to evaluate the discrete eigenvalues $`\lambda _k^\pm `$ and the corresponding ‘normalization’ constants $`C_k^\pm `$ of the Jost solutions; this is easy only in the limit $`r_0\mathrm{}`$, when $`\lambda _k^\pm =\mu _k\pm i\nu _k`$. Other difficulties come from the fact that in many of the applications we need to analyze trains in which: a) the solitons move with nearly the same velocities; b) various perturbations should be taken into account. In such cases the exact approach based on the inverse scattering method can not be directly applied and one should look for other methods . Our aim is to extend our previous results on the $`N`$-soliton interaction in the adiabatic approximation . Firstly we show that the CTC model is an universal one in the sense that it describes the $`N`$-soliton interactions for all NLEE from the NLS hierarchy. We derive the perturbed CTC (PCTC) system and show that small perturbations affect only the center of mass motion and the global phase of the $`N`$-soliton train. Special attention is paid to the (anti-) soliton interaction of the sine-Gordon equation. We show that in the adiabatic approximation their interaction is described by the Toda chain with indefinite metric, which is a special reduction of the CTC. We also outline the peculiarities of the interactions in the non-adiabatic cases. This paper is an extended version of . ## 2 The $`N`$-soliton interactions and the Complex Toda Chain In the quasi-particle approach of Karpman and Solov’ev has been generalized to any $`N>2`$ soliton train in the adiabatic approximation. This means that the solitons initially must have nearly equal amplitudes and velocities and must be well separated, i.e.: $`\|\nu _{k,0}\nu _{j,0}\|\nu _0,\|\mu _{k,0}\mu _{j,0}\|\mu _0,`$ $`\nu _0(\xi _{k+1,0}\xi _{k,0})r_01;\|\nu _{k,0}\nu _0\|(\xi _{k+1,0}\xi _{k,0})1,`$ (4) where the additional zeroes in the subscripts in (2) refer to the value at $`t=0`$ and $`\nu _0`$ and $`\mu _0`$ are the average amplitude and velocity of the $`N`$-soliton train. The result is a dynamical system of equations for the $`4N`$ soliton parameters called the generalized Karpman-Solov’ev system (GKS). The GKS is adapted to treat also the perturbed NLS equation. An exhaustive list of perturbations, relevant for nonlinear optics, which include linear and nonlinear dispersive and dissipative terms, effects of sliding filters, amplitude and phase modulations, etc. is studied . We prove that the linear perturbations affect each of the solitons separately, while the nonlinear ones lead to additional interactive terms between neighboring solitons. Another important step which allowed us to analyze the $`N`$-soliton interactions analytically consists in the fact, that under some additional approximations the GKS reduces to the complex Toda chain (CTC) with $`N`$ nodes : $$\frac{d^2q_k}{dt^2}=16\nu _0^2\left(\mathrm{e}^{q_{k+1}q_k}\mathrm{e}^{q_kq_{k1}}\right),$$ (5) where $`k=1,\mathrm{},N`$ and we assume that $`\mathrm{e}^{q_0}\mathrm{e}^{q_{N+1}}0`$. The complex dynamical variables $`q_k(t)`$ are expressed in terms of the parameters of the $`k`$-th soliton by: $`q_k(t)=2i(\mu _0+i\nu _0)\xi _k(t)i(\delta _k(t)+\delta (t))+kQ_0,`$ $`Q_0=\mathrm{ln}4\nu _0^2+i\pi ,\delta (t)={\displaystyle \frac{1}{N}}{\displaystyle \underset{s=1}{\overset{N}{}}}\delta _k(t).`$ (6) The same result has been derived also by using the variational approach ; this approach however should be applied with care, see . The CTC with $`N`$ nodes, which may be viewed as a natural generalization of the corresponding real Toda chain (RTC), provides a very convenient tool to study the soliton interactions. Numeric simulations show that CTC provides an adequate description for the soliton interactions for a wide class of initial conditions (2) . It is also possible to describe the soliton interactions of other soliton-type nonlinear equations with different dispersion laws. For any such equation one can derive the corresponding GKS and CTC model which could be useful to study the interactions of their $`N`$-soliton trains, see . The adiabatic approximation imposes certain restrictions not only on the soliton parameters (see Eq. (2)), but also on the spectral data of the Zakharov-Shabat system $`L`$. Firstly the discrete eigenvalues in the upper $`\lambda `$-half plane of $`L`$ must be located in a small neighborhood around $`\lambda _0`$: $$\|\lambda _k^+\lambda _0\|^2\epsilon ,\lambda _0=\underset{k=1}{\overset{N}{}}\lambda _k^+/N,$$ (7) where the small number $`\epsilon `$ determines the overlap between the neighboring solitons. With the same precision the eigenvalues $`\lambda _k^+`$ can be approximated by $`2\zeta _k`$ where $`\zeta _k`$ are the eigenvalues of the Lax matrix for the CTC. Secondly we have a condition on the initial values of the constants $`C_k`$ which determine the initial positions and phases of the pulses. Skipping the details we get: $`\mathrm{ln}\left\|C_{k+1}^+(0)/C_k^+(0)\right\|=2\nu _0(\xi _{k+1,0}\xi _{k,0})+𝒪(1)2\mathrm{ln}\epsilon 1.`$ (8) The conditions on the discrete eigenvalues $`\lambda _k^+`$ allow us also to explain the universality of the CTC as a model describing the $`N`$-soliton interactions. Namely, we claim that the CTC model describes in the adiabatic approximation the $`N`$-soliton interactions of all NLEE in the NLS hierarchy. Indeed, let us consider a higher NLS equation with dispersion law $`F(\lambda )`$, regular in the vicinity of $`\lambda _0`$. Then the time-dependence of $`C_k^+`$ is given by $$C_k^+(t)=\mathrm{exp}(2iF(\lambda _k^+)t)C_k^+(0).$$ (9) As a consequence the one-soliton solution will be given by (1) with $`z_k(x,t)`$ and $`\delta _k(t)`$ replaced by $`z_k(x,t)=2\nu _k\left(x{\displaystyle \frac{f_{1,k}}{\nu _k}}t\xi _{0,k}\right),`$ (10) $`\delta _k(t)=\delta _{k,0}+{\displaystyle \frac{2(\mu _kf_{1,k}\nu _kf_{0,k})}{\nu _k}}t,`$ (11) where $`F(\lambda _k^+)=f_{0,k}+if_{1,k}`$. However, due to (7) in fact it is enough to take into account only the first three terms in the Taylor expansion: $`F(\lambda _k^+)`$ $`=`$ $`F(\lambda _0)+(\lambda _k^+\lambda _0)\dot{F}_0+{\displaystyle \frac{1}{2}}(\lambda _k^+\lambda _0)^2\ddot{F}_0+𝒪(\epsilon ^{3/2}),`$ (12) where $`\dot{F}_0=(dF/d\lambda )\|_{\lambda =\lambda _0}`$, $`\ddot{F}_0=(d^2F/d\lambda ^2)\|_{\lambda =\lambda _0}`$. Comparing (9) and (10) we see that in the adiabatic approximation only the first three terms in (12) are important for the soliton parametrization. This explains why the soliton interactions for all the equations from the hierarchy is described by the same universal model: CTC. As an example of higher NLS equations which also finds important applications in nonlinear optics is the one with dispersion law $`F(\lambda )=2\lambda ^2+\eta \lambda ^3`$ introduced in ; the GKS for this equation is derived in . The corresponding CTC model is obtained from (5) by replacing the coefficient $`16\nu _0^2`$ with a factor depending on $`F(\lambda _0)`$ which can be taken care of by redefining $`q_k`$. Another possibility is to choose $`F(\lambda )=1/(2\lambda )`$ which after additional reduction leads to the sine-Gordon equation, see Section 4 below. These results have been further developed by using the fact, that the CTC is a completely integrable dynamical system with $`2N`$ degrees of freedom. The most important consequence of this fact lies in the possibility to predict the asymptotic behavior of the solitons from the set of their initial parameters . Indeed, knowing the initial soliton parameters we can construct the eigenvalues of the Lax matrix for the CTC system, which in turn determine the asymptotic velocities of the solitons. A more detailed study of the solutions of the CTC allowed us to see that it allows much richer class of asymptotic regimes than the RTC . We are also able to describe the class of initial soliton parameters, that lead to each one of these regimes: i) asymptotically free propagation of the solitons (the only regime allowed by RTC); ii) $`N`$-soliton bound states with the possibility of a quasi-equidistant propagation; iii) mixed asymptotic regimes when part of the solitons form bound state(s) and the rest separate from them; iv) regimes corresponding to the degenerate and singular solutions of the CTC. In , a thorough comparison between the CTC predictions with the numerical solutions of the NLS equation has been performed and an excellent match has been established for a number of choices of the initial soliton parameters in each of the regimes listed above. Special attention has been paid to regime ii) and more specifically to the possibility for a quasi-equidistant (QED) propagation of all $`N`$ solitons. A method for the description of the corresponding initial soliton parameters responsible for this regime has been proposed. ## 3 Perturbed NLS and the perturbed CTC In we showed also that the evolution of the $`N`$-soliton train (2) of the perturbed NLS equation (1) is described by the following dynamical system for the ‘slow’ evolution of the soliton parameters: $`{\displaystyle \frac{d\nu _k}{dt}}=16\nu _0^2(S_kS_{k+1})+N_k,`$ (13) $`{\displaystyle \frac{d\mu _k}{dt}}=16\nu _0^2(C_kC_{k+1})+M_k,`$ (14) $`{\displaystyle \frac{d\xi _k}{dt}}=2\mu _k+\mathrm{\Xi }_k^{(0)}+\mathrm{\Xi }_k,`$ (15) $`{\displaystyle \frac{d\delta _k}{dt}}=2(\mu _k^2+\nu _k^2)+X_k^{(0)}+X_k,`$ (16) where $`\mathrm{\Xi }_k^{(0)}=4(S_k+S_{k+1}),`$ (17) $`X_k^{(0)}=2\mu _k\mathrm{\Xi }_k^{(0)}+24\nu _k(C_k+C_{k+1}),`$ (18) $`C_k(t)iS_k(t)={\displaystyle \frac{1}{4\nu _0}}\mathrm{e}^{q_k(t)q_{k1}(t)}.`$ (19) The terms $`N_k`$, …, $`X_k`$ are determined by $`R[u]`$ below. As it was shown in they contain two types of terms: a) ‘self-interaction’ terms depending only on the parameters of the $`k`$-th soliton and b) ‘nearest-neighbour’ interaction terms containing linear combinations of $`S_k`$, $`C_k`$, $`S_{k+1}`$ and $`C_{k+1}`$. In we also derived the explicit expressions for $`M_k`$, …, $`X_k`$ in terms of the soliton parameters for several classes of physically important perturbations. Here we take into account linear and cubic in $`u`$ perturbations including the linear and nonlinear gain, third order dispersion (TOD), intrapulse Raman scattering (IRS) etc, i.e.: $$R[u]=\underset{k=0}{\overset{3}{}}c_k\frac{^ku}{x^k}+d_0\|u\|^2u+\frac{d_1}{4}u(\|u\|^2)_x+\frac{d_2}{4}(\|u\|^2u_xu_x^{}u^2),$$ (20) where $`c_s`$ and $`d_s`$ are generically complex parameters: $$c_s=c_{s0}+ic_{s1},d_s=d_{s0}+id_{s1}.$$ Some of these coefficients, namely $`c_{01}`$, $`c_{21}`$ and $`d_{01}`$ can be put to zero without restrictions; this can be done by conveniently renormalizing $`u`$, $`t`$ and $`x`$. The next argument which we will use is that the coefficients in (20) must be small. We start by assuming that they are, like the terms $`S_k`$ and $`C_k`$, of the order of $`\epsilon `$; at the same time the deviations $`\stackrel{~}{\nu }_k=\nu _k\nu _0`$, $`\stackrel{~}{\mu }_k=\mu _k\mu _0`$ are of the order of $`\sqrt{\epsilon }`$. Therefore in the right hand sides of the equations (13)-(14) we have only terms of the order of $`\epsilon `$, while in the r.h.sides of (15)-(16) we have also terms of the order of 1 and $`\sqrt{\epsilon }`$. That is why we will simplify the perturbative terms in the r.h.sides of (13)-(16) by taking only the first few terms in their Taylor expansions, i.e. $`Z_k(\nu _k,\mu _k)=Z_{00}+\stackrel{~}{\nu }_kZ_{10}+\stackrel{~}{\mu }_kZ_{01},Z_{00}=Z(\nu _0,\mu _0),`$ (21) $`Z_{10}={\displaystyle \frac{Z}{\nu _k}}\|_{\begin{array}{c}\nu _k=\nu _0\\ \mu _k=\mu _0\end{array}},Z_{01}={\displaystyle \frac{Z}{\mu _k}}\|_{\begin{array}{c}\nu _k=\nu _0\\ \mu _k=\mu _0\end{array}},`$ where $`Z`$ stands for each of the functions $`N_k`$, $`M_k`$, $`\mathrm{\Xi }_k`$ and $`X_k`$. The explicit expressions for the coefficients in (21) for each of the four functions are given in the appendix. ### 3.1 Perturbations of order $`\epsilon `$. Note that due to our assumption about the perturbation constants all coefficients in $`Z_{00}`$ are of the order of $`\epsilon `$; so in fact we have to take into account only $`N_{00}`$ and $`M_{00}`$. As a result we derive the following perturbed version of the CTC model: $$\frac{d^2q_k}{dt^2}=U_{00}+16\nu _0^2\left(\mathrm{e}^{q_{k+1}q_k}\mathrm{e}^{q_kq_{k1}}\right),$$ (22) where $`U_{00}=4\mu _0N_{00}4\nu _0(M_{00}+2iN_{00})`$. In deriving (22) we took into account also the fact that now $`\lambda _0=_{k=1}^N\lambda _k/N`$ and consequently $`\nu _0`$ become time-dependent: $$\frac{d\lambda _0}{dt}=\frac{1}{N}\underset{k=1}{\overset{N}{}}(M_k+iN_k)M_{00}+iN_{00},\frac{d\nu _0}{dt}N_{00}.$$ Eq. (22) can be solved exactly with the result: $$q_k(t)=\frac{1}{2}U_{00}t^2+V_{00}t+q_k^{(0)}(t),$$ (23) where $`q_k^{(0)}(t)`$ is a solution of the unperturbed CTC and $`V_{00}`$ is an arbitrary constant. If in addition we assume that $`V_{00}=0`$ and $`\mu _0=0`$ then from the formulae in the appendix we find that: $$U_{00}=\frac{16\nu _0^2}{3}\left\{\nu _0\left[c_{11}+\frac{4}{5}\nu _0^2(d_{11}7c_{31})\right]i\left[3c_{00}+4\nu _0^2(2d_{00}c_{20})\right]\right\}.$$ (24) We remind that $`q_k^{(0)}`$ is related to the $`k`$-th soliton parameters by (2). Then from (23) we see that $`\text{Re }U_{00}`$ and $`\text{Re }V_{00}`$ influence the center of mass motion of the train while $`\text{Im }U_{00}`$ and $`\text{Im }V_{00}`$ drive the global phase $`\delta (t)`$. In particular for the special case when $`\text{Re }U_{00}=0`$ the effect of such perturbation will be to make the phase of all solitons oscillate simultaneously with a rate proportional to $`t^2`$. However the evolution of the phase and coordinate differences $`\delta _{k+1}\delta _k`$ and $`\xi _{k+1}\xi _k`$ will not be influenced. ### 3.2 Perturbations of order $`\sqrt{\epsilon }`$. More complicated and substantially different is the situation when the perturbative constants become of order $`\sqrt{\epsilon }`$. Then the terms $`N_k`$, $`M_k`$ in (13), (14) can be approximated by linear combinations of $`\stackrel{~}{\nu }_k`$, $`\stackrel{~}{\mu }_k`$; as a result Equation (13) acquires the form: $`{\displaystyle \frac{d\nu _k}{dt}}=N_{00}+N_{10}\stackrel{~}{\nu }_k+N_{01}\stackrel{~}{\mu }_k+16\nu _0^2(S_kS_{k+1}),`$ (25) $`{\displaystyle \frac{d\mu _k}{dt}}=M_{00}+M_{10}\stackrel{~}{\nu }_k+M_{01}\stackrel{~}{\mu }_k16\nu _0^2(C_kC_{k+1}),`$ (26) $`{\displaystyle \frac{d\xi _k}{dt}}=2\mu _k+\mathrm{\Xi }_{00}+\mathrm{\Xi }_{10}\stackrel{~}{\nu }_k+\mathrm{\Xi }_{01}\stackrel{~}{\mu }_k+\mathrm{\Xi }_k^{(0)},`$ (27) $`{\displaystyle \frac{d\delta _k}{dt}}=2(\mu _k^2+\nu _k^2)+X_{00}+X_{10}\stackrel{~}{\nu }_k+X_{01}\stackrel{~}{\mu }_k+X_k^{(0)},`$ (28) Then $`N_{00}`$ and $`M_{00}`$ (of order $`\sqrt{\epsilon }`$) will be the leading order terms in (25), (26) while terms like $`c_0\stackrel{~}{\nu }_k`$, $`c_2\stackrel{~}{\mu }_k`$ are of the same order $`\epsilon `$ as the interaction terms (ones with $`S_k`$ and $`C_k`$). The solutions of such PCTC will be qualitatively different from the ones of CTC and require separate studies. ## 4 Sine-Gordon solitons and the CTC Here we shortly discuss the (anti-) soliton interactions of the sine-Gordon equation: $$v_{xt}+\mathrm{sin}v(x,t)=\epsilon R[v].$$ (29) The problem has been attacked long time ago by Spatschek and Karpman and Solov’ev where the interaction of two (anti-) solitons has been studied. The one-soliton solution is given by $$v^{1\mathrm{s}}=2\sigma \mathrm{arcsin}\mathrm{tanh}z\pm \pi ,z=2\nu (x\xi (t)),$$ where $`\xi (t)=t/(4\nu ^2)+\xi _0`$. For $`N>2`$ and $`R[v]=0`$ in the adiabatic approximation only the nearest neighbor interactions are relevant. Then the results of generalize to: $`{\displaystyle \frac{d\nu _k}{dt}}`$ $`=`$ $`4\left(\mathrm{e}^{Q_{k+1}Q_k}\mathrm{e}^{Q_kQ_{k1}}\right),`$ (30) $`{\displaystyle \frac{d\xi _k}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{4\nu _k^2}}+{\displaystyle \frac{1}{\nu _0^2}}\left(\mathrm{e}^{Q_kQ_{k1}}+\mathrm{e}^{Q_{k+1}Q_k}\right),`$ (31) $`Q_k(t)`$ $`=`$ $`2\nu _0\xi _k(t)+{\displaystyle \frac{i\pi }{2}}\left(1\sigma _k\right),`$ (32) where $`\sigma _k=1`$ (or $`1`$) if at position $`k`$ we have soliton (or anti-soliton). Note that the sine-Gordon (anti-) solitons do not have internal degrees of freedom and are characterized only by their amplitudes $`\nu _k`$ and positions $`\xi _k`$; here we choose $`\xi _1<\xi _2<\mathrm{}<\xi _N`$ and $`\nu _0`$ is the average amplitude. We remark also that the sine-Gordon equation is related to the Zakharov-Shabat system $`L[u]`$ if we assume $`u=v_x/2`$ and request that $`v`$ is real. This last condition has important consequences: i) besides the (anti-) soliton solutions related to the purely imaginary eigenvalues $`\lambda _k^\pm =\pm i\nu _k`$ of $`L`$, the sine-Gordon equation has also breather solutions; ii) one can not have two (anti-) solitons moving with the same speed, i.e. $`\nu _k\nu _j`$ for $`kj`$. The adiabatic approximations (7) mean that $`\|\nu _k\nu _{k+1}\|𝒪(\sqrt{\epsilon })`$ and in addition we have (8). Thus we see that in the right hand side of (31) only the first term is the relevant one; the other two are of the order of $`𝒪(\epsilon )`$ and can be neglected. If we now differentiate (31) with respect to $`t`$, use (30) and keep only terms of order $`\epsilon `$ we get: $$\frac{d^2Q_k}{dt^2}=\frac{4}{\nu _0^2}\left(\mathrm{e}^{Q_{k+1}Q_k}\mathrm{e}^{Q_kQ_{k1}}\right),$$ (33) where we replaced $`\nu _k`$ in the denominator by $`\nu _0`$. If all $`\sigma _k`$ are equal we obtain the real Toda chain. It is also known from the spectral properties of the Zakharov-Shabat system, that in the sine-Gordon case we can not have two (anti-) solitons moving with the same velocities. This means that if we have a sequence of solitons (or anti-solitons) only then their interaction is purely repulsive and their asymptotic regime can contain only asymptotically free ‘particles’. This facts are compatible with the analytical results on the sine-Gordon solitons, see . The model (33) with generic $`\sigma _k`$ has been studied by Kodama and Ye and is known as the Toda chain with indefinite metric. In both cases we can view (33) as special reduction of the CTC. Thus we see that the second involution on the Zakharov-Shabat needed for the sine-Gordon equation, carries over as a reduction on the CTC. We should also note that the equations (32), (33) have solutions with singularities which are periodic in time, see . The comparizon between the $`N`$ (anti-) soliton train dynamics of the sine-Gordon equation and the indefinite metric Toda chain is yet to be done. This and the studies of the PCTC for the perturbed sine-Gordon equation will be published elsewhere. ## 5 Non-adiabatic Interactions If one or more of the ‘adiabatic conditions’ (2) are violated then the picture becomes much more complicated. It is possible that due to strong perturbation some of the soliton pulses come very close to each other and strongly overlap. Usualy this is combined with strong deformations of the pulses and substantial emission of ‘radiation’ which is not accounted for in our model. To our knowledge there are no effective models which would provide analytic description of the soliton interactions in such situations. As main tool giving a physical insight of the soliton dynamics is the comparison between the numerical solutions to the NLS equation (1) and the numerical solution of the corresponding Zakharov-Shabat spectral problem . The method consists in the following: first we solve numerically the corresponding (perturbed) NLS equation using the standard fast-Fourier transform (or beam-propagation) method. Then we use the results for the pulse shape evaluated at a certain distance as an initial potential for the Zakharov-Shabat eigenvalue problem and determine numerically its scattering data . As a result we can determine the time evolution of the scattering data (including the data, characterizing the continuous spectrum). The advantage of this method is the possibility to follow up the variations of the amplitudes and the velocities of an arbitrary number of solitons. Here it is possible in a natural way to estimate the energy of the ‘radiation’, related to the continuous spectrum of $`L`$. The disadvantage is in the necessity to know approximately the locations of the eigenvalues of $`L`$ at $`t=0`$. Below we use the basic fact that the unperturbed NLS equation is integrable. As a consequence the evolution of $`u(x,t)`$ preserves the spectrum of corresponding Zakharov-Shabat system $`L`$, which may be determined from the initial condition $`u(x,t=0)`$. In particular the discrete eigenvalues of $`L`$ will be time-independent since they are integrals of motion of the NLSE. If we next consider perturbed NLS equation then generically the perturbation will violate the integrability. However we assume that the perturbation is ‘small’ in the sence that it does not destroy completely the integrability but rather slightly modifies the spectrum of $`L`$. In particular the eigenvalues $`\lambda _k^\pm `$ of $`L`$ start to move; here and below the upperscript $`+`$ ($``$) means that the corresponding eigenvalue is such that $`\text{Im }\lambda _k^+>0`$ ($`\text{Im }\lambda _k^{}<0`$). We remind that the involution $`\lambda _k^+=(\lambda _k^{})^{}`$ holds, so it is enough to know only the discrete eigenvalues $`\lambda _k^+`$. To our knowledge there are no explicit criteria which would allow one to check whether given perturbation is ‘small’ or not. Some inexplicit criteria have been formulated in ; in particular they require that the eigenvalues $`\lambda _k^+`$ remain in the upper half-plane (i.e. $`\text{Im }\lambda _k^+>0`$ for all $`t`$), that they do not come close to the real axis and that they do not coalesce. In terms of the soliton parameters the second of these condition means that the amplitude of the pulses should not become very small. In an investigation of the influence of the intrapulse Raman scattering and the third order dispersion on the discrete eigenvalues of the Zakharov-Shabat system have been performed by numerical means. Two qualitatively different initial conditions approximating two-soliton bound states have been studied. The first one $`u_1(x,t=0)=2\text{sech }(x)`$ correponds to strongly overlapped soliton pulses; so in this case the adiabatic approximation is not valid. The spectrum of $`L_1`$ (i.e. of $`L`$ with potential given by $`u_1(x,t=0)`$) consists of two eigenvalues in the upper half-plane with $`\lambda _k^+=i(k1/2)`$, $`k=1,2`$ (and two more in the lower half-plane). Note also that the distance between these eigenvalues is not small. For the second one $`u_2(x,t=0)=\text{sech }(x\delta )+\text{sech }(x+\delta )`$ with $`\delta 3÷4`$ the pulses are well separated and the adiabatic approximation holds. The spectrum of $`L_2`$ can not be calculated presicely; besides the two pairs of eigenvalues it contains also some small ‘radiation’. The eigenvalues $`\lambda _k^+`$ can be well approximated by the eigenvalues of the Lax matrix for the corresponding CTC. In our case this give: $$\lambda _1^+\frac{i}{2}\left(1+e^\delta \right)\lambda _2^+\frac{i}{2}\left(1e^\delta \right).$$ (34) Note that already for $`\delta 3÷4`$ the quantity $`e^\delta `$ may be considered as small (of the order of $`\sqrt{\epsilon }`$; i.e., these eigenvalues satisfy the adiabaticity condition. Obviously, if take $`\delta `$ to be smaller then the overlap of the solitons grows and the adiabaticity is violated. As a result (34) does not give correct values for the eigenvalues. In the limit of infinitely separated soliton pulses, i.e. $`\delta \mathrm{}`$ the eigenvalues (34) coalesce. This is related to the fact that the solution to the CTC with these initial conditions is singular, see . The effect of these two perturbation on the eigenvalues of $`L`$ are similar for both types of initial conditions. In what follows we describe it for the non-adiabatic case with $`u_2(x,t=0)`$. First we analyse the effect of the third order dispersion (TOD) which is a Hamiltonian perturbation. The time evolution of the eigenvalues is shown on Fig. 1 for different strengths $`c_{30}`$ of TOD. For $`c_{30}0.01`$ it turnes out that the imaginary parts are almost constant while the real parts are zero. It is known also that there exist a critical value $`c_{30,\mathrm{cr}}=0.022`$ where the two-soliton bound state breaks down. For $`c_{30}=0.02`$, which is just below the critical value two strongly fluctuating real parts show up. Minor fluctuations of the imaginary parts can also be identified. For $`c_{30}=c_{30,\mathrm{cr}}`$ the very splitting of the degenerate real parts $`\kappa _1=\kappa _2=0`$ appears. After some transition time both real parts attain constant but different values. The imaginary parts remain almost unchanged. This stage of deformation of the eigenvalues caused by TOD corresponds to the break up of the two-soliton bound state into two single, progressively separating solitons with amplitudes determined by the initial imaginary parts of the eigenvalues. This behavior of the eigenvalues is consistent with the previous results. If $`c_{30}`$ grows even larger (e.g., $`c_{30}0.03`$), the smaller imaginary part changes significantly. This second stage of deformation of the eigenvalues can be described by the ultimate differences between both real and imaginary parts, respectively, which increase with $`c_{30}`$. The change in the imaginary part leads to the creation of ‘radiation’. A similar investigation was performed in order to analyze the effect of dissipative perturbation such as intrapulse Raman scattering on the eigenvalues. The results are shown in Fig. 2. A remarkable fact to be mentioned is that unlike for TOD a very weak perturbation ($`d_{11}=0.0004`$) lifts the degeneracy of the real part. The two-soliton bound state breaks up and two slowly separating solitons with different but constant amplitudes emerge. The second stage in the deformation (changes in the imaginary parts) start from $`d_{11}=0.02`$ and differs from the TOD case in that both imaginary parts change. In contrast to TOD the effect shows up for considerably weaker perturbations. The big change of the larger amplitude ($`d_{11}<0.2`$) causes a strong variation of the corresponding real part due to the amplitude dependence of the soliton self-frequency shift. These results clearly illustrate the qualitatively different effect of Hamiltonian (TOD) and dissipative (IRS) perturbations on the soliton bound states. The numeric evaluation of the spectral data of $`L`$ for each step of propagation of the soliton train also allows one to control the precision of the numerical procedure used to solve the NLSE . Another possible effect of the strong perturbations is that the pulses taken initially to be one-soliton solutions of the NLS, may deform into the exact travelling-wave solutions of the perturbed NLS equation. Such effect has been reported in due to the nonlinear gain and bandwidth limited amplification. Then the perturbed NLS equation goes int the Ginzburg-Landau (GL) equation whose stationary solutions possess characteristic phase modulation (chirp). It is due to this modulation that the soliton interaction reduces substantially. ## 6 Conclusions Starting from the GKS model proposed in we have derived the perturbed CTC system describing the $`N`$-soliton train interaction of the perturbed NLS equation in the adiabatic approximation. For small perturbations the PCTC system is again completely integrable and provides us with an effective tool for analytic study of the asymptotic regimes of the $`N`$-soliton trains. In the non-adiabatic regime we propose a combined numeric solution of the NLS equation and the Zakharov-Shabat problem. Finally we mention that these methods can be applied also to the adiabatic interaction of the multicomponent NLS equation and its perturbed versions. Such equations describe the birefringence effects and soliton interactions in multi-mode fibers. We expect that their soliton interactions will be described by a generalized CTC-model in which the soliton phases $`\delta _k`$ are replaced by ‘polarization’ vectors $`\stackrel{}{n}_k`$, see . These results can be used in soliton-based fiber-optics communications. It is not difficult to treat also the perturbed sine-Gordon equation and derive the corresponding perturbed versions of (33). This and the study of the interactions of (anti-) solitons with breathers will be published elsewhere. Finally we stress on the universal character of the CTC in the sense that it is independent on the dispersion law of the equation whose soliton interactions it describes. ## Acknowledgements We thank Prof. E. Doktorov for useful discussions. This work was supported in part by contract F-807 with the National Science Fund of Bulgaria. ## Appendix A The coefficients $`Z_{\alpha ,\beta }`$. Here we list the coefficients $`Z_{\alpha ,\beta }`$ where $`Z`$ takes the values $`N_k`$, $`M_k`$, $`\mathrm{\Xi }_k`$ and $`X_k`$ while the pair of indices $`(\alpha ,\beta )`$ is one of the following $`(0,0)`$, $`(1,0)`$ and $`(0,1)`$, see formula (21). With $`c_{01}=c_{21}=d_{01}=0`$ (see the remark after eq. (20)) we have: $`N_{00}=4\nu _0({\displaystyle \frac{c_{00}}{2}}c_{11}\mu _0{\displaystyle \frac{2}{3}}c_{20}(\nu _0^2+3\mu _0^2)+8c_{31}\mu _0(\nu _0^2+\mu _0^2)`$ $`+{\displaystyle \frac{4}{3}}\nu _0^2(d_{00}d_{21}\mu _0)),`$ $`N_{10}=2\left(c_{00}2c_{11}\mu _0+(8c_{31}\nu _0\mu _04c_{20})(\nu _0^2+\mu _0^2)+8\nu _0^2(d_{00}d_{21}\mu _0)\right),`$ $`N_{01}=4\nu _0\left(c_{11}+4c_{20}\mu _04c_{31}(\nu _0^2+3\mu _0^2){\displaystyle \frac{4}{3}}d_{21}\nu _0^2\right),`$ $`M_{00}={\displaystyle \frac{4}{3}}\nu _0^2\left(c_{11}+4c_{20}\mu _012c_{31}\left(\mu _0^2+{\displaystyle \frac{7}{15}}\nu _0^2\right)+{\displaystyle \frac{4}{5}}\nu _0^2d_{11}\right),`$ $`M_{10}={\displaystyle \frac{8}{3}}\nu _0\left(c_{11}+4c_{20}\mu _012c_{31}\mu _0^2{\displaystyle \frac{56}{5}}c_{31}\nu _0^2+{\displaystyle \frac{8}{5}}d_{11}\nu _0^2\right),`$ $`M_{01}={\displaystyle \frac{4}{3}}\nu _0^2\left(c_{20}12c_{31}\mu _0\right),\mathrm{\Xi }_{00}=c_{10}+4c_{30}(\nu _0^2+3\mu _0^2){\displaystyle \frac{2}{3}}\nu _0^2d_{10},`$ $`\mathrm{\Xi }_{10}={\displaystyle \frac{4}{3}}\nu _0\left(6c_{30}d_{10}\right),\mathrm{\Xi }_{01}=24c_{30}\mu _0,`$ $`X_{00}=c_{01}16c_{30}\mu _0(\nu _0^2\mu _0^2)+4\nu _0^2\left(\mu _0d_{20}{\displaystyle \frac{1}{3}}\mu _0d_{10}\right),`$ $`X_{10}=8\nu _0\left(4c_{30}\mu _0d_{20}\mu _0+{\displaystyle \frac{1}{3}}d_{10}\mu _0\right),`$ $`X_{01}=48\mu _0^2c_{30}+4\nu _0^2\left(d_{20}{\displaystyle \frac{1}{3}}d_{10}\right).`$
warning/0004/nlin0004038.html	ar5iv	text	# Escape from noisy intermittent repellers \[ ## Abstract Intermittent or marginally-stable repellers are commonly associated with a power law decay in the survival fraction. We show here that the presence of weak additive noise alters the spectrum of the Perron - Frobenius operator significantly giving rise to exponential decays even in systems that are otherwise regular. Version : February 3, 2000 \] There exist a variety of physical situations where one deals with the escape of trajectories from repellers. For instance, nuclear physicists are interested in the escape of particles along fission channels. Similarly, in the study of transport coefficients in two-dimensional ballistic conductors, one has to worry about the trapping time and its relationship with the geometry of the system . By and large, it is now accepted that hyperbolic (chaotic) dynamics leads to an exponential decay in the number of trapped particles while intermittency or marginal stability results in power law decays. Thus, a rectangular billiard table with a hole in the wall shows a power law decrease in the survival fraction while a (hyperbolic) enclosure created by the intersection of three discs shows an exponential decay in the number of trapped particles. The reason for this difference is intuitively clear from the following argument: Consider that there are N particles distributed uniformly in (Birkhoff) phase space and there is a hole along the wall at $`q_0`$ of extent $`\mathrm{\Delta }q_0`$. The (average) fraction of particles that escape at each bounce is identical for a chaotic system and proportional to $`\mathrm{\Delta }q_0`$ since the particles remain uniformly distributed with time. Thus the survival fraction decays exponentially. In a marginally stable system however, an initial uniform distribution does not remain uniform at each bounce since individual particles tend to stick around stable islands. A heuristic derivation of the power law decay can be found in and we merely remark here that the decay exponent is often difficult to determine analytically and an interesting advancement in this direction has been achieved recently by Dalhqvist . We are interested here in a situation where marginal stability or intermittency is accompanied by weak additive noise. Such a situation can arise for instance in an imperfectly fabricated ballistic conductor in the shape of a triangle or stadium where reflection is no longer specular but has additive noise. Thus : $`q_{n+1}`$ $`=`$ $`f_1(q_n,p_n)`$ (1) $`p_{n+1}`$ $`=`$ $`f_2(q_n,p_n)+\xi _n`$ (2) where $`q,p`$ are the Birkhoff coordinates , $`f_1,f_2`$ are the bounce maps and $`\xi _n`$ is a random variable with $`\xi _n=0`$ having a normalized distribution $`g(\xi )`$ (normally taken to be a Gaussian with zero mean). The question that we pose is : does one expect to find a power law decay in such a situation ? The answer we believe is interesting and can significantly alter the way people look at signatures of low-dimensional chaos in various experimental situations where noise is inevitable and often desirable. First, however, we shall consider a 1-dimensional intermittent map and study its spectrum in the presence of weak noise. A trajectory in the presence of additive noise is generated by the iteration $$x_{n+1}=f(x_n)+\xi _n$$ (3) where $`f(x)`$ is a map, $`\xi _n`$ is a random variable as described above and $`x_0[a,b]`$. An initial density of trajectories $`\varphi (x)`$ evolves according to the Perron - Frobenius equation : $$(_0\varphi )(x)=𝑑y\delta (xf(y))\varphi (y)$$ (4) in the unperturbed case. Thus, the eigenvalues and eigenfunctions of $`_0`$ determine the escape rate in an open systems. More specifically, assuming that the spectrum is discreet, an initial density $`\varphi (x)`$ can be expanded as $$\varphi (x)=\underset{\alpha }{}c_\alpha \phi _\alpha (x)$$ (5) so that the fraction of particles that survive $`n`$ iterates of the map is $`\mathrm{\Gamma }(n)`$ $`=`$ $`{\displaystyle \frac{_a^b𝑑x\left(_0^n\varphi \right)(x)}{_a^b𝑑x\varphi (x)}}={\displaystyle \underset{\alpha }{}}\mathrm{\Lambda }_\alpha ^nc_\alpha {\displaystyle \frac{_a^b𝑑x\phi _\alpha (x)}{_a^b𝑑x\varphi (x)}}`$ (6) $``$ $`\mathrm{\Lambda }_0^n=\mathrm{e}^{n\mathrm{ln}(1/\mathrm{\Lambda }_0)}\mathrm{as}n\mathrm{}.`$ (7) In the above, $`\{\mathrm{\Lambda }_\alpha \}`$ are the eigenvalues corresponding to the eigenfunctions $`\{\phi _\alpha (x)\}`$ and $`\mathrm{\Lambda }_0`$ is the leading eigenvalue with the largest real part. The discreetness assumption however holds only when the dynamics is hyperbolic. In the presence of marginally stable cycles, the spectrum has a continuous part leading to a power law decay of correlations (in closed systems) or survival fraction (in open systems). The presence of additive noise results in a modified kernel whose formal expression is well known : $$(\varphi )(x)=𝑑y(x,y)\varphi (y)$$ (8) where $`(x,y)`$ $`=`$ $`\delta _\sigma (xf(y))`$ (9) $`\delta _\sigma (x)`$ $`=`$ $`{\displaystyle \delta (x\xi )g(\xi )𝑑\xi }=g(x).`$ (10) As before, if the spectrum of $``$ is discreet, the decay is exponential and the leading eigenvalue determines the asymptotic decay rate. Note that we have so far steered clear of spill-over effects due to noise . The most commonly adopted technique is the use of periodic boundary conditions which avoids spill-over altogether. Alternately, one can work in the infinite domain $`(\mathrm{},\mathrm{})`$ so that natural boundary conditions may be employed. Yet another approach is to tailor the noise distribution so that the probability of the dynamical variable escaping from the interval is zero. We shall, in this paper, have occasion to use the second and third approaches depending on the problem and it must be noted that there are other approaches to the spill-over problem that may be more realistic in a given situation. Note that the spectrum and eigenfunctions of $``$ can be sensitive to the choice of boundary conditions. With this background, we now introduce the intermittent map : $$f(x)=\{\begin{array}{cc}x(1+x^2)\hfill & x<0\hfill \\ x(1+p(2x)^s)\hfill & 0x<1/2\hfill \\ 2x1\hfill & x>1/2\hfill \end{array},$$ (11) where $`s>0`$ and $`p>1`$. The particle is considered to escape if $`x_{n+1}<0`$ or $`x_{n+1}>1`$. The map is defined in the infinite domain so that natural boundary conditions apply on the density. The intermittency here is due to the fact that $`f^{}(0)=1`$ so that the fixed point $`x=0`$ is marginally or neutrally stable. For an initial uniform distribution of particles in $`[0,1]`$, the fraction that survives one iterate is clearly the sum of the two intervals $`I_L`$ and $`I_R`$ for which $`0f(x)1`$. Similarly, the fraction that survives two iterates is the sum of the four intervals $`I_{LL},I_{LR},I_{RL},I_{RR}`$ for which $`0f^2(x)1`$. Generalization in this case (of binary symbolic dynamics) is simple : the fraction that survives $`n`$ iterates is the sum of the $`2^n`$ intervals for each of which $`0f^n(x)1`$. Each of these intervals contains a periodic point and the larger its (in)stability, the smaller the size of the interval. Thus : $$I_q^{\{n\}}=\frac{a_q}{\mathrm{\Lambda }_q}$$ (12) where $`q`$ is a symbol sequence of length $`n`$ consisting of $`L`$ and $`R`$ which denotes the order in which the left and right branches (with respect to $`x=1/2`$) of the map are visited, $`a_q`$ is a constant and $`\mathrm{\Lambda }_q=\frac{d}{dx}f^n(x)\|_{xI_q}`$ is the stability of the periodic point. The survival fraction can thus be expressed as : $`\mathrm{\Gamma }(n)`$ $`=`$ $`{\displaystyle \underset{q}{\overset{\{n\}}{}}}I_q={\displaystyle \underset{q}{\overset{\{n\}}{}}}{\displaystyle \frac{a_q}{\mathrm{\Lambda }_q}}{\displaystyle \underset{p}{}}{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n_p\delta _{n,rn_p}}{\left\|\mathrm{\Lambda }_p\right\|^r}}=𝒵_n,`$ (13) $`𝒵_n`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _\gamma }z^n\left({\displaystyle \frac{d}{dz}}\mathrm{log}\zeta ^1(z)\right)𝑑z`$ (14) where $`\zeta ^1(z)=_p\left(1z^{n_p}/\left\|\mathrm{\Lambda }_p\right\|\right)`$ is the dynamical zeta function and $`\gamma `$ is a (small) negatively oriented contour around the origin. Dahlqvist has recently shown that in the noiseless case, $`\zeta ^1(z)`$ has a singularity of the type $`(1z)^{1/s}`$. It then follows from eq. (14) that the survival fraction, $`\mathrm{\Gamma }(n)1/n^{1/s}`$ for an initial uniform distribution of particles. The zeta function is also (approximately) related to the eigenvalues of $``$ through the relation : $$\mathrm{Tr}^n=\underset{\alpha }{}\mathrm{\Lambda }_\alpha ^n=\underset{p}{}\underset{r=1}{\overset{\mathrm{}}{}}\frac{n_p\delta _{n,rn_p}}{\left\|1\mathrm{\Lambda }_p\right\|^r}𝒵_n$$ (15) When the zeta function is analytic, its zeroes, $`\{z_k\}`$ are isolated and related to $`\mathrm{\Lambda }_\alpha `$ as $`\mathrm{\Lambda }_\alpha =1/z_k`$. On the other hand, when the system is intermittent and $`\zeta ^1(z)`$ displays a branch cut, the spectrum of $``$ no longer remains discreet. Thus intermittency leads to an asymptotic power law decay. For $`s=0.7,p=1.2`$, the initial decay is however exponential and this is ascribed to a pair of complex conjugate roots . The power law behaviour emerges only after 600 iterations of the map. We now consider the map (eq.(11)) with weak additive Gaussian noise $$g(\xi )=\frac{1}{\sqrt{2\pi \sigma ^2}}e^{\frac{\xi ^2}{2\sigma ^2}}$$ (16) for $`s=0.9,0.7`$ and $`0.5`$ and with $`\sigma =0.002`$ (see figure 1). Clearly there is a transition from a power law to an exponential decay in the presence of weak noise in each of the three cases for large $`n`$. Note, however, that the initial decay, though exponential, is at a significantly different rate and the slope of $`\mathrm{\Gamma }(n)`$ settles down to the asymptotic value gradually after a large number of iterations. Two inferences can be drawn from this transition from power-law to exponential behaviour. First, the presence of weak noise makes the eigenvalues (of $``$, see eq. (8) ) discreet. Also, there are closely spaced eigenvalues with small differences in their real parts around the leading eigenvalue $`\mathrm{\Lambda }_0`$ which leads to the gradual change in slope of $`\mathrm{\Gamma }(n)`$. The discreetness of the spectrum follows from the fact that the noisy kernel is integrable and bounded. The corresponding Fredholm determinant is thus entire whose zeroes ($`1/\mathrm{\Lambda }_\alpha `$) are isolated. Thus, even for very weak noise, the spectrum is discreet though the transition time may be too large for the final exponential decay to be observed experimentally. The closely spaced eigenvalues around $`\mathrm{\Lambda }_0`$ are possibly remnants of the continuous spectrum which exists for the noiseless case. In order to understand this better, we have evaluated the eigenvalues of $``$ by discreetizing the integral equation and diagonalizing the resulting matrix . Recall that a Fredholm integral equation exists as a limit of discreet sum so that a matrix representation is adequate so long as its order is large and spurious eigenvalues are eliminated. Figure 2 shows a plot of the eigenvalues of $``$ for $`s=0.7`$ and $`\sigma =0.002`$ (see eq. (11)) obtained using a matrix of size $`2500\times 2500`$. We have checked that the relevant eigenvalues on the positive real axis have converged to the fourth significant digit. Clearly, the closely spaced eigenvalues along the real line do not allow the survival-fraction to be dominated by the leading eigenvalue for small $`n`$. An order of magnitude evaluation of the transition time can be made by noting that $`\mathrm{\Gamma }(n)e^{n\mathrm{ln}(1/\mathrm{\Lambda }_0)}\left(d_0+d_1e^{n\mathrm{ln}(\mathrm{\Lambda }_0/\mathrm{\Lambda }_1)}\right)`$ where $`\mathrm{\Lambda }_1`$ is the next to leading eigenvalue and $`d_\alpha =c_\alpha 𝑑x\varphi _\alpha (x)`$. Thus, for $`n>>n_{trans}=1/\mathrm{ln}(\mathrm{\Lambda }_0/\mathrm{\Lambda }_1)`$, the leading eigenvalue dominates. For the three values of $`s`$ considered, $`n_{trans}98,105,119`$ for $`s=0.5,0.7,0.9`$ respectively so that exponential decay sets in first for $`s=0.5`$ as observed in fig. 1. Thus the difference between the leading and the next-to-leading eigenvalue gives a good qualitative picture and fixes a lower bound for the transition time. Note also that for each of the three $`s`$-values considered, the leading eigenvalue accurately reproduces the asymptotic decay of the survival-fraction (see Table I). We now turn our attention to the effect of noise in non-chaotic systems. Specifically, we shall consider triangular billiards which are non-chaotic though generically non-integrable. The only integrable examples are the $`(\pi /3,\pi /3,\pi /3)`$, $`(\pi /2,\pi /3,\pi /6)`$ and the $`(\pi /2,\pi /4,\pi /4)`$ triangles while all other rational triangles have at least one internal angle of the form $`m\pi /n,m>1`$ and are non-integrable. Their invariant surface, though $`2`$-dimensional, is not a torus but topologically equivalent to a sphere with multiple holes. Also these systems are non-chaotic though irrational triangles are possibly ergodic and even display the weak mixing property . A linear stability analysis shows that the Jacobian matrix has unit eigenvalues and hence these billiards are marginally stable . Consider such a triangular billiard of unit perimeter and let $`q,p`$ denote the Birkhoff co-ordinates. Here $`q`$ is measured along the boundary while $`p=\mathrm{sin}(\theta )`$ where $`\theta `$ is the angle between the ray and the inward normal at the boundary point $`q`$. Thus $`q[0,1]`$ and $`p[1,1]`$. In a typical experiment, one considers an initial uniform distribution of particles ($`10^8`$) in this phase space evolving freely in between bounces and reflecting specularly from the walls. The particles are allowed to escape through a small opening at $`q_0`$ of extent $`\mathrm{\Delta }q_0(=0.005)`$. For both the integrable $`(\pi /2,\pi /3)`$ and non-integrable $`(18\pi /31,17\pi /97)`$ triangles considered, the initial decay is exponential while the asymptotic decay is a power law, $`\mathrm{\Gamma }(n)n^\beta `$, with $`\beta =1.035`$ in the integrable case and $`\beta =1.085`$ in the non-integrable (NI) case. Thus pre-exponential decays are not exceptional and can persist for a long time in marginally stable systems. A more realistic situation should however include noise. For instance, imperfections can give rise to maps of the type considered in eq. (2). This leads to interesting results. For Gaussian noise (in Birkhoff momentum) with $`\sigma =0.000001`$, a single exponential decay dominates the survival fraction in the NI case for $`n>5000`$ while in the integrable case, the transition continues beyond $`n=14000`$. Moreover, the closely spaced eigenvalues leads to a quasi-algebraic decay in the interval $`6000<n<14000`$ for the integrable case. Note that most trajectories remain largely unaffected for several hundred bounces for the value of $`\sigma `$ considered so that $`\mathrm{\Gamma }(n)`$ closely follows the noiseless case initially. Thus, even when the asymptotic decay is exponential, transition times can be very large. In the weak noise case, the evolution operator can be approximated as $$(\varphi )(x)\xi (x)\frac{\sigma ^2}{2}\xi ^{\prime \prime }(x)$$ (17) where $`\xi (x)=_i\varphi (x)/\left\|f^{}(f_i^1(x))\right\|`$. Using a polynomial basis , a matrix representation of the operator can be constructed where the elements $`_{mn}=\frac{d^m}{dx^m}\{(\frac{x^n}{n!})\}`$. Writing $`\sigma `$ as $`\sigma _1+\sigma _0`$ where $`\sigma _00^+`$, a perturbation calculation shows that the eigenvalues (and thus the difference between $`\mathrm{\Lambda }_0`$ and $`\mathrm{\Lambda }_1`$) decrease as $`\sigma _1^2`$ when $`\sigma _1^2`$ is small. The transition time therefore decreases with noise. For $`\sigma =0.00005`$ (see fig. 3), the transition time decreases significantly and exponential decay sets in for $`n>3500`$ in the integrable case while the leading eigenvalue dominates from the beginning in the non-integrable case. Thus the gap between $`\mathrm{\Lambda }_0`$ and $`\mathrm{\Lambda }_1`$ increases with noise. In conclusion, the broad picture that emerges from these numerical experiments is as follows : (i) Additive noise makes the spectrum of the evolution operator discreet. (ii) When the dynamics is intermittent or regular and the noise weak, exponential decays may emerge only asymptotically due to the presence of closely spaced eigenvalues around the leading eigenvalue, $`\mathrm{\Lambda }_0`$. These are remnants of the continuous spectrum that exists in the zero noise case. The transition phase in such a situation can mimic an algebraic decay. There are important fallouts of this conclusion. In experimental situations where noise is inevitable, signatures of exponential decays are not necessarily indicative of chaotic dynamics. For instance, semiclassical theory links Lorentzian line shapes observed in experiments on ballistic transport in chaotic microstructures to the exponential decay in the survival fraction . The present analysis however shows that noisy intermittent dynamics can also give rise to Lorentzian line shapes and it is interesting to note that there are instances where observations on regular or marginally stable cavities have been found to be no different from chaotic cavities .
warning/0004/cond-mat0004348.html	ar5iv	text	# Cyclically coupled spreading and pair annihilation ## 1. Introduction The study of phase transitions into absorbing states is a fascinating field of nonequilibrium statistical physics . Such phase transitions can be observed in models for the spreading of some (generally non-conserved) agent as a result of a competition between local reproduction and decay. If the rate for reproduction is sufficiently high, the system is able to maintain a fluctuating active phase where the stationary concentration of the spreading agent is positive. On the other hand, if the decay process dominates, the concentration of the spreading agent decreases and tends to zero. Eventually the system becomes trapped in an absorbing state from where it cannot escape. An interesting situation emerges at the borderline between survival and extinction. Here the system undergoes a nonequilibrium phase transition which is characterized by non-trivial critical behavior. It is believed that transitions into absorbing states can be categorized into a finite number of universality classes. Typically each class is associated with certain symmetry properties of the dynamics. At present two universality classes are firmly established. The first and most prominent one is the universality class of directed percolation (DP) . The DP class is characterized by the absence of symmetries (apart from conventional symmetries such as translation and reflection invariance) and covers a wide range of models . Roughly speaking, DP models follow the reaction-diffusion scheme $`X2X`$, $`X\text{Ø}`$. In addition, there has to be a nonlinear mechanism which limits the particle density. In ‘fermionic’ models with at most one particle per site, the density is limited automatically. In ‘bosonic’ models, allowing for infinitely many particles per site, this mechanism has to be implemented by adding the reaction $`2XX`$. It is important to note that DP is a unary spreading process, i.e., individual particles are able to reproduce and destruct themselves. Some time ago Grassberger et al. discovered a second universality class which can be considered as a generalization of directed percolation from one to two absorbing states related by an exact $`Z_2`$-symmetry (DP2). It comprises various models, including nonequilibrium Ising models , certain monomer-dimer models , as well as generalized versions of the Domany-Kinzel model and the contact process . In 1+1 dimensions it is possible to regard kinks between differently oriented absorbing domains as particles. By definition, these particles evolve according to a parity-conserving dynamics, exemplified by branching-annihilating random walks with two offspring $`X3X`$, $`2X\text{Ø}`$ . For this reason, the universality class is often referred to as the parity-conserving (PC) class. However, it should be noted that the PC class and the DP2 class are different in higher dimensions. Three years ago, Howard and Täuber raised the question whether there might exist a third universality class of models with a phase transition from ‘real’ to ‘imaginary’ noise. As a prototype, they introduced the annihilation-fission (AF) process $`2X3X`$, $`2X\text{Ø}`$ with single-particle diffusion. Obviously, this process is neither parity-conserving nor invariant under any other unconventional symmetry transformation. In contrast to DP and DP2, the AF model is a binary process, i.e., only pairs of particles can decay or reproduce themselves. The critical properties at the transition are still poorly understood. Recently, Carlon et. al. analyzed the transition using density matrix renormalization group techniques . Estimating the critical exponents they arrived at the conclusion that the transition of the AF process should belong to the DP2 universality class, although there is no $`Z_2`$ symmetry or parity conservation law. However, as pointed out in Ref. , various physical arguments suggest that the AF process might belong to an independent universality class which has not been investigated before. Since the universality class could not be identified so far, it is of interest to understand the most salient features of the transition in the AF process from a descriptive point of view. To this end, we introduce an effective model which separates the dynamics of pairs and solitary particles in the AF process by introducing two species of particles $`A`$ and $`B`$. The $`A`$’s perform an ordinary DP process while the $`B`$’s are subjected to an annihilating random walk. Both subsystems are cyclically coupled by particle transmutation. It is shown that this three-state model exhibits a nonequilibrium phase transition which is similar to the one observed in the AF process. The following Section briefly summarizes the phenomenological properties of the AF process. In Sec. 3 we demonstrate that the critical behavior is governed by two competing dynamic modes for spreading and diffusion. Based on this interpretation we introduce an effective model involving two species of particles (see Sec. 4). In Sec. 5 the critical exponents of this model are estimated by Monte Carlo simulations. The article ends with several concluding remarks in Sec. 6. ## 2. Phenomenological properties of the annihilation-fission process Assuming the usual scaling picture for phase transitions into absorbing states, we expect the AF process to be characterized by four critical exponents $`\beta `$, $`\beta ^{}`$, $`\nu _{}`$, and $`\nu _{}`$. The first one is associated with the field-theoretic annihilation operator and describes the behavior of the stationary density $`\rho _{\text{stat}}(pp_c)^\beta `$ close to the transition. The exponent $`\beta ^{}`$ is associated with the creation operator and plays a role whenever initial conditions are specified. For example, the survival probability of a cluster grown from a single seed involves the exponent $`\delta ^{}=\beta ^{}/\nu _{}`$. The other two exponents are related to the spatial and temporal correlation lengths $`\xi _{}\|pp_c\|^\nu _{}`$ and $`\xi _{}\|pp_c\|^\nu _{}`$, respectively. It should be noted that in the case of DP the two order parameter exponents $`\beta `$ and $`\beta ^{}`$ coincide because of a duality symmetry (see e.g. Ref. ). In the present case, however, they turn out to be different. Using a density matrix renormalization group approach, Carlon et. al. estimated the critical exponents of the AF process for various diffusion rates $`0.1d0.2`$. Since their estimates $`z=1.73\mathrm{}1.81`$ and $`\beta /\nu _{}=0.46\mathrm{}0.5`$ were in fair agreement with the numerical values of DP2 exponents (see Table 1), they concluded that the AF process should belong to the DP2 universality class, although there is no $`Z_2`$-symmetry or parity conservation law in the system. This conclusion, however, collides with the general believe that critical phenomena are determined by their symmetry properties or, equivalently, by their associated field theories. Without the required symmetries on the microscopic level, a special mechanism would be necessary in order to restore these symmetries effectively on large scales. As there is no such mechanism, it is near at hand to expect that the AF process does not belong to the DP2 class. To support this point of view, preliminary Monte-Carlo simulations were presented in Ref. and later improved by Grassberger and Ódor . Because of strong corrections to scaling, the estimates for the critical exponents depend on the numerical effort. The estimates for $`\delta =\nu _{}/\nu _{}`$, for example, are scattered over the range $`0.25\mathrm{}0.29`$ and seem to decrease with increasing simulation time. A similar tendency was observed for the exponents $`z`$ and $`\beta `$. Moreover, it is not yet fully clear to what extent the exponents depend on the diffusion rate $`d`$. A tentative list of critical exponents, including recent results obtained by simulations on parallel computers , is given in Table 1. In order to understand the transition in the annihilation-fission process from a phenomenological point of view it is helpful to analyze the spatio-temporal structure of critical clusters. To this end we introduce a novel type of scale-invariant space-time plot which can be used to visualize the scaling properties of critical clusters in systems with absorbing states. Starting with a localized seed (a pair of particles) at the origin and simulating the spreading process up to $`10^6`$ time steps, the rescaled position of the particles $`x/t^{1/z}`$ is plotted against $`\mathrm{log}_{10}t`$, where $`z=\nu _{}/\nu _{}`$ is the dynamic exponent of the process under consideration (see Fig. 1). By rescaling the spatial coordinate $`x`$, the cluster is confined to a strip of finite width. Compared to linear space-time plots the scale-invariant representation of a cluster has several advantages. On the one hand, it is possible to survey more than four decades in time. On the other hand, the plot provides a simple visual check of scaling invariance. Roughly speaking, scaling invariance is fulfilled if the cluster’s appearance is time-independent, i.e., spatio-temporal patterns should look similar in the upper and lower parts of the figure. It is needless to say that this visual check does not replace an accurate quantitative analysis. Nevertheless such a scale-invariant plot may improve the intuitive understanding of the asymptotic long-time behavior and may also help to identify relevant and irrelevant contributions. Let us first consider a critical cluster of a branching-annihilating random walk with two offspring (see left part of Fig. 1). Obviously, this process is characterized by an ongoing competition between particle reproduction and decay. Contrarily, the annihilation-fission process admits undisturbed random walks of solitary particles over long distances. As shown in the middle of Fig. 1, this leads to a very different visual appearance of the cluster. It is important to note that these differences persist as time proceeds. However, if both processes were to converge to the same type of long-range critical behavior as suggested in , we would expect the clusters to become increasingly similar in the lower part of the figure. As there is no indication of such a convergence, Fig. 1 supports the viewpoint of Refs. that the AF process might represent a new type of nonequilibrium critical behavior. The corresponding universality class should be characterized by the following properties: 1. Single particles diffuse but do not react. 2. Reproduction requires two particles to meet at neighboring sites. 3. Particles are removed if at least two particles meet at neighboring sites. 4. There is no unconventional symmetry (such as parity conservation). 5. There is no frozen disorder in the system. 6. There is a mechanism limiting the density of particles. Consequently, many other processes, such as the fission-coagulation model $`2X3X`$, $`2XX`$ and the reaction-diffusion process $`2X3X`$, $`3XX`$ are expected to exhibit the same type of nonequilibrium critical behavior. ## 3. Interpretation as a spreading process with two species of particles To what extent is the transition in the AF model different from ordinary DP and DP2 transitions? As Fig. 1 suggests, there are two separate dynamic modes for spreading and diffusion. The spreading mode is characterized by sudden avalanches with a high density of particles. Here the dynamic processes are dominated by interacting pairs of particles. Once an avalanche has stopped, the system enters the diffusive mode, in which solitary particles perform a simple random walk. When two of them meet at neighboring sites, they may release a new avalanche, as illustrated in Fig. 2. The asymptotic critical behavior at the transition will depend on the relevance of the two dynamic modes. In principle there are three possibilities. If the spreading mode becomes dominant we expect a crossover to DP. Conversely, if the diffusive mode governs the asymptotic regime, we expect a purely diffusive behavior with the dynamical critical exponent $`z=2`$. However, Fig. 1 strongly suggests that both modes are equally important and balance one another as time proceeds. In fact, in the rescaled representation the typical spatio-temporal patterns do not change over four decades in time. In order to investigate this transition in more detail, we suggest a phenomenological explanation of the observed spatio-temporal patterns. The basic idea is to describe the two dynamic modes in terms of two separate reaction-diffusion processes involving two different species of particles $`A`$ and $`B`$. The spreading mode is governed by the dynamics of $`A`$-particles. Roughly speaking, the $`A`$-particles can be thought of as representing pairs of particles in the original model. Obviously, there is no parity conservation law for the number of pairs, hence the $`A`$-particles in the new model evolve effectively in the same way as in an ordinary DP process. During the avalanche, several $`A`$’s transmute into $`B`$’s. Therefore, once the avalanche has stopped, several $`B`$-particles are left behind. These $`B`$-particles in turn perform a simple random walk, representing solitary particles in the original model. When two $`B`$-particles meet, they may trigger a new avalanche of $`A`$-particles. The corresponding reaction-diffusion scheme reads (1) $$A2A,A\text{Ø}/B,2BA.$$ Thus the model consists of two subsystems A and B. Subsystem A is a DP process $`A2A,A\text{Ø}`$ which is coupled via transmutation $`AB`$ to subsystem B. In this subsystem the $`B`$-particles diffuse until they annihilate and release a new avalanche of $`A`$-particles. Thus, the reaction-diffusion model (1) can be interpreted as a cyclically coupled sequence of a DP and a pair annihilation process, as sketched in Fig. 3. It should be pointed out that Figs. 2-3 are over-simplified in the sense that they suggest the existence of strictly separated dynamic modes. In reality, however, the two modes are not completely separated, rather they are entangled and sustain each other in a intricate manner. Nevertheless it is the hope of the present study that the AF-process and the simplified reaction-diffusion scheme (1) exhibit at least qualitatively a similar type of nonequilibrium critical behavior. ## 4. A lattice model for cyclically coupled directed percolation and pair annihilation In order to study the process (1) quantitatively, we introduce a three-state model on a square lattice with random-sequential updates which is defined by the following dynamic rules: \| reproduction: \| $`\text{Ø}AAA`$ \| with rate \| $`p/2`$ \| \| --- \| --- \| --- \| --- \| \| \| $`BAAA`$ \| \| $`p/2`$ \| \| \| $`A\text{Ø}AA`$ \| \| $`p/2`$ \| \| \| $`ABAA`$ \| \| $`p/2`$ \| \| decay: \| $`A\text{Ø}`$ \| \| $`(1p)(1\tau )`$ \| \| transmutation: \| $`AB`$ \| \| $`(1p)\tau `$ \| \| diffusion: \| $`\text{Ø}BB\text{Ø}`$ \| \| $`D/2`$ \| \| annihilation: \| $`BBA\text{Ø}`$ \| \| $`r/2`$ \| \| \| $`BB\text{Ø}A`$ \| \| $`r/2`$ \| In the following the numerical analysis will be restricted to the case $`r=D=1`$. Thus, the model is controlled by two parameters, namely the rate for particle reproduction $`p`$ and the transmutation rate $`\tau `$. Obviously, the model has two absorbing states, i.e., the empty lattice and the state with a single diffusing $`B`$-particle. The phase diagram in Fig. 4 comprises two phases. For low values of $`p`$ and $`\tau `$, the system is in the inactive phase where it approaches one of the two absorbing states. If $`p`$ and $`\tau `$ are sufficiently large, an active steady state with non-vanishing particle densities $`\rho _A`$ and $`\rho _B`$ exists on the infinite lattice. Here we are interested in the critical behavior at the phase transition line. It should be noted that the case $`p=0`$ is special. In this case the spreading process does not generate $`B`$ particles. Hence, starting with $`A`$-particles, the critical behavior belongs to the DP universality class. For $`p>0`$, however, we observe non-DP critical behavior. In fact, an ordinary DP transition seems to be unlikely since the inactive phase is charactrized by an algebraic decay of the particle density. Similarly, we can rule out the possibility of a DP2 transition since there is neither a $`Z_2`$-symmetry nor a parity conservation law in the model. ## 5. Numerical results In order to estimate the critical exponents in 1+1 dimensions, we perform Monte Carlo simulations. Starting with randomly distributed $`A`$’s and $`B`$’s, we first measure the densities $`\rho _A(t)`$ and $`\rho _B(t)`$ up to $`310^5`$ time steps. As can be seen in the left panel of Fig. 5, the ratio $`\rho _B(t)/\rho _A(t)`$ approaches a constant value. Assuming an algebraic decay, both quantities should therefore scale with the same critical exponent (2) $$\rho _A(t)\rho _B(t)t^\delta .$$ The temporal decay of the particle densities is shown in the right panel of Fig. 5. As in the case of the AF model, we observe strong corrections to scaling, leading to a considerable curvature of the data. However, compared to the AF model these corrections are less severe. Moreover, the curvatures for $`\rho _A`$ and $`\rho _B`$ have opposite signs. Thus, the local slopes approach the postulated ‘true’ value of $`\delta `$ from both sides. Seeking for the best compromise, we are able to estimate the critical points by \| $`\tau `$ \| $`0`$ \| $`0.1`$ \| $`0.5`$ \| $`1`$ \| \| --- \| --- \| --- \| --- \| --- \| \| $`p_c`$ \| 0.7674(3) \| $`0.757(2)`$ \| $`0.6920(1)`$ \| $`0.540(1)`$ \| For $`\tau =0.5`$ the corresponding exponent is given by $`\delta =\beta /\nu _{}=0.21(2)`$. Similar measurements for $`\tau =0.1`$ and $`\tau =1`$ (not shown here) suggest that this value is the same for all $`\tau >0`$. In order to obtain the dynamic exponent $`z=\nu _{}/\nu _{}`$, we perform finite size simulations at criticality. Here the particle densities should obey the scaling form (3) $$\rho _A(t)\rho _B(t)t^\delta f(t/L^z),$$ where $`f`$ is a universal scaling function. Thus, plotting $`\rho t^\delta `$ against $`t/L^z`$, all data sets should collapse onto a single curve. Comparing different data collapses it turns out that $`\rho _B`$ shows a much cleaner scaling behavior than $`\rho _A`$. As shown in Fig. 6, the best collapse is obtained for $`\delta =0.215(15)`$ and $`z=1.75(5)`$. To obtain the third exponent, we study the behavior $`\rho _B(t)`$ below and above criticality. According to the usual scaling theory for absorbing-state transition, we expect the scaling form (4) $$\rho _B(t)t^\delta g(tϵ^\nu _{}),$$ where $`ϵ=\|pp_c\|`$ denotes the distance from criticality. Plotting $`\rho _B(t)t^\delta `$ against $`tϵ^\nu _{}`$ (see Fig. 6), the best data collapse is obtained for $`\delta =0.215(20)`$ and $`\nu _{}=1.8(1)`$. In order to cross-check these estimates, we perform dynamic simulations starting with a single pair of particles located in the center . As usual in this type of simulations, we measure the survival probability $`P(t)`$ that the system has not yet reached one of the two absorbing states, the average numbers of particles $`N_A(t)`$ and $`N_B(t)`$, and the mean square spreading of all particles from the origin $`R^2(t)`$ averaged over the surviving runs. Assuming that $`N_A(t)`$ and $`N_B(t)`$ scale asymptotically with the same exponent, these quantities should obey the power laws (5) $$P(t)t^\delta ^{},N_A(t)N_B(t)t^\eta ,R^2(t)t^{2/z}$$ with certain dynamical exponents $`\delta ^{}`$ und $`\eta `$. As shown in Fig. 7, the survival probability shows a clean power law over almost five decades. Moreover, the quotient $`N_B(t)/N_A(t)`$ quickly tends to a constant value. Fitting power laws to the data shown in Fig. 7, we obtain the estimates (6) $$\delta ^{}=0.15(1),\eta =0.21(1),2/z=1.16(4).$$ Together with the previous results these estimates satisfy the generalized hyperscaling relation (7) $$\delta +\delta ^{}+\eta d/z=0$$ within numerical errors. Combining all results, the critical exponents are given by (8) $$\begin{array}{cc}\hfill \beta =0.38(6),& \beta ^{}=0.27(3),\hfill \\ \hfill \nu _{}=1.8(1),& \nu _{}=1.0(1).\hfill \end{array}$$ It should be noted that the error bars were obtained by assuming power-law behavior. Thus, they do not include systematic errors due to possible corrections to scaling emerging after very long time. ## 6. Conclusions This work was motivated by recent studies of the annihilation-fission process which exhibits a novel type of nonequilibrium critical behavior. Using a scale-invariant space-time plot we have demonstrated that critical clusters of the AF process are characterized by an interplay of two different dynamic modes. In the high-density mode we observe spreading avalanches whereas the low-density mode is characterized by random walks of solitary particles. In order to understand the interplay between these modes from a phenomenological point of view, we have introduced an effective model which involves two species of particles. The dynamic rules of this model can be regarded as being composed of cyclically coupled DP and pair annihilation processes, i.e., it follows the reaction-diffusion scheme $`A2A`$, $`AB`$, $`2BA`$. The model exhibits a nonequilibrium phase transition which is in many respects similar to the one observed in the AF process. In fact, the AF process and the three-state model proposed in the present work have several features in common: * They both have two non-symmetric absorbing states, namely the empty lattice and the state with a single diffusing particle. * There is no unconventional symmetry (such as parity conservation). * Both models exhibit a continuous nonequilibrium phase transition with non-DP critical exponents. * The visual appearance of critical clusters is very similar (see Fig. 1). Compared to the AF process, the three-state model has several advantages. On the one hand, it is defined as a two-site nearest neighbor process. On the other hand, corrections to scaling are less severe, allowing us to determine the critical exponents more accurately. There are several open questions. Is the critical behavior of the three-state model universal? If this is indeed the case, does it represent a independent universality class different from DP and DP2? Are the AF process and the present model in the same universality class, i.e., do they have the same critical exponents? From the field-theoretic point of view, there is no reason for them to coincide. Yet the two classes may ‘intersect’ in 1+1 dimensions. Very recently, Ódor carried out a systematic study of the annihilation/fission process performing simulations on a parallel computer combined with generalized mean field approximations and coherent anomaly extrapolations . He reports two different universality classes for low and high values of the diffusion constant $`d`$. The corresponding critical exponents are shown in Table 1. According to Ódor, the exponents for $`d0.5`$ are in fair agreement with the exponents observed in the three-state model. However, this conclusion is still speculative and needs to be substantiated by further investigations. Acknowledgements: I would like to thank E. Carlon, P. Grassberger, M. Henkel, M. Howard, J. F. Mendes, G. Ódor, U. Schollwöck, and U. Täuber for fruitful discussions.
warning/0004/cond-mat0004355.html	ar5iv	text	# LPENSL-TH 05/2000 Spin accumulation in the semi classical and quantum regimes ## 1 Introduction The discovery of the Giant Magneto Resistance (GMR) in magnetic multilayers has generated an important interest. These systems are made of a sandwich of alternating ferromagnetic and non magnetic layers. Valet and Fert proposed a semi classical description of the perpendicular GMR, on the basis of a Boltzmann equation incorporating a spin-dependent transport in the presence of spin accumulation (see also ). Spin accumulation occurs in the GMR because the current arising from a ferromagnet is spin polarized, and therefore cannot penetrate a ferromagnet with an opposite magnetization. Instead, spin accumulates at the interface. This phenomenon occurs also at the interface between a ferromagnet and a superconductor, where a spin polarized current cannot penetrate the superconductor . Here, we would like to reconsider two particular aspects of spin accumulation, namely, (i) in the semi classical regime, the possibility of an inverse Drude scaling of the conductance meaning that, in some parameter range, the conductance increases with the length of the conductor; and (ii) in the quantum coherent regime, the existence of a resonance in the spin flip channels. More precisely, we study a ferromagnet - normal metal - ferromagnet spin valve, in which we assume the presence of magnetic scattering in the normal metal . The inverse Drude scaling resulting from spin accumulation is already implicitly contained in the equations obtained by Valet and Fert , but, to our knowledge, this effect has not been studied previously per se in the literature, which we do here. Spin accumulation corresponds to the presence of a different chemical potential for the spin-up and spin-down electrons, which obey a spin diffusion equation . Our treatment is not based on the spin diffusion equation, but relies on an exact solution of the 1D Boltzmann equation where we can make an exact decoupling between the charge and spin sectors. Next, we ask to what extend a quantum model can show a similar physics. We are lead to study a quantum “spin-flip Fabry Perot” interferometer in which a single magnetic impurity is located at a given distance $`a`$ away from a ferromagnet interface. We find the existence of a Fabry-Perot resonance in the spin-flip channels as the parameter $`a`$ is varied. This resonance disappears as the ferromagnet spin polarization is decreased, and can therefore be viewed as the equivalent of the inverse Drude behavior in the quantum coherent regime. Our treatment is based on a Landauer approach, similar to the one used by Zhu and Wang to study the effect of magnetic scattering close to a superconductor interface . The article is organized as follows. Section 2 is devoted to the solution of the semi classical transport equations. We solve the “spin-flip Fabry Perot” interferometer model in section 3. Final remarks are given in the Conclusion. ## 2 Transport in the semi classical regime ### 2.1 Boltzmann equation and boundary conditions We consider a model in which magnetic impurities are present in a normal metal close to a ferromagnet interface (see Fig. 1). We neglect any Kondo correlation , which is an assumption valid above the Kondo temperature. The presence of the ferromagnets close to the normal metal may lead to magnetic flux lines penetrating inside the normal metal, which can orient the magnetic impurities in a preferential direction. We implicitly assume that the temperature is high enough so that the impurities have no preferential orientation. We consider a one dimensional model because only in this geometry can we decouple the spin and charge sectors of the Boltzmann equation. We assume that the interfaces between the ferromagnets and the normal metal are sharp, and that the exchange field has a step function variation at the interface. We note $`f_{R,L}^\sigma (E,x)`$ the semi classical distribution function of right/left moving spin-$`\sigma `$ electrons with an energy $`E`$ at position $`x`$. The Boltzmann equation in the relaxation time approximation reads $$\frac{}{x}\left(\begin{array}{c}f_R^{}(E,x)\\ f_L^{}(E,x)\\ f_R^{}(E,x)\\ f_L^{}(E,x)\end{array}\right)=\left(\begin{array}{cccc}(r+r_s+r_s^{})& r& r_s& r_s^{}\\ r& r+r_s+r_s^{}& r_s^{}& r_s\\ r_s& r_s^{}& (r+r_s+r_s^{})& r\\ r_s^{}& r_s& r& r+r_s+r_s^{}\end{array}\right)\left(\begin{array}{c}f_R^{}(E,x)\\ f_L^{}(E,x)\\ f_R^{}(E,x)\\ f_L^{}(E,x)\end{array}\right),$$ (1) where we have discarded the term involving the electric field. This is valid if the temperature of the electrodes is larger than the applied voltage, in which case the electronic gas has a temperature identical to the one of the electrodes . Therefore, we should consider a finite temperature and calculate the low voltage conductance in the regime $`eVT`$. In practise, we consider the limit $`T0`$, and calculate the linear conductance. The coefficients $`r`$, $`r_s`$ and $`r_s^{}`$ in Eq. 1 denote respectively the rate of backscattering without spin-flip, the rate of forward scattering with spin-flip, and the rate of backward scattering with spin-flip. The coefficients can be related to the $`q=0`$ and $`q=2k_f`$ components of the microscopic scattering potential (see the Appendix). We now explicit the boundary conditions. For this purpose, let us consider an interface between a ferromagnet in the region $`x<0`$ and a normal metal in the region $`x>0`$, and include interface scattering under the form of repulsive potential $`H\delta (x)`$ . Let us first consider a spin-up electron incoming from the left ferromagnet, and denote by $`b^{}`$ and $`t^{}`$ the backscattering and transmission coefficients. The wave function in the region $`x<0`$ is $`\psi _L(x)=\mathrm{exp}(ik^{}x)+b^{}\mathrm{exp}(ik^{}x)`$, and the wave function in the region $`x>0`$ is $`\psi _R(x)=t^{}\mathrm{exp}(ikx)`$. The matching equations are $`\psi _L(0)=\psi _R(0)=\psi (0)`$ and $`\psi _R(0)/x\psi _L(0)/x=(2mH/\mathrm{}^2)\psi (0)`$, from what we deduce $$t^{}=\frac{2ik^{}}{i(k+k^{})2mH/\mathrm{}^2}\text{, and }b^{}=\frac{i(k^{}k)+2mH/\mathrm{}^2}{i(k+k^{})2mH/\mathrm{}^2}.$$ The probability current conservation can be verified easily: $`k^{}=k^{}\|b^{}\|^2+k\|t^{}\|^2`$. The spin-up conductance is found to be $`G^{}=(e^2/h)\|T^{}\|^2`$, with the transmission coefficient $$T^{}=\frac{k}{k^{}}\|t^{}\|^2=\frac{4kk^{}}{(k+k^{})^2+\left[\frac{2mH}{\mathrm{}^2}\right]^2}.$$ (2) The backscattering coefficient is $`B^{}=1T^{}`$. In the spin-down sector, we obtain $`T^{}`$ and $`B^{}`$ by substituting $`k^{}`$ with $`k^{}`$ in Eq. 2. This provides the boundary conditions for the Boltzmann equation: $`f_R^{}(E,0)`$ $`=`$ $`T^{}f_T(EeV)+(1T^{})f_L^{}(E,0)`$ (3) $`f_R^{}(E,0)`$ $`=`$ $`T^{}f_T(EeV)+(1T^{})f_L^{}(E,0)`$ (4) $`f_L^{}(E,L)`$ $`=`$ $`T^{}f_T(E)+(1T^{})f_R^{}(E,L)`$ (5) $`f_L^{}(E,L)`$ $`=`$ $`T^{}f_T(E)+(1T^{})f_R^{}(E,L),`$ (6) where the left and right ferromagnets are assumed to be in equilibrium and $`f_T(E)`$ denotes the Fermi-Dirac distribution function. We consider $`T^{}`$ and $`T^{}`$ to be independent of energy, which amounts to considering the wave vectors $`k`$ and $`k^{}`$ in Eq. 2 to be on the Fermi surface. Eqs. 3 – 6 provide a simple form for the spin-up and spin-down currents at positions $`x=0,L`$. For instance at $`x=L`$, we have $`I^{}(L)`$ $`=`$ $`T^{}{\displaystyle \frac{e}{h}}{\displaystyle \left[f_R^{}(E,L)f_T(E)\right]𝑑E}`$ (7) $`I^{}(L)`$ $`=`$ $`T^{}{\displaystyle \frac{e}{h}}{\displaystyle \left[f_R^{}(E,L)f_T(E)\right]𝑑E}.`$ (8) The Boltzmann equation Eq. 1 and the boundary conditions Eqs. 3– 6 lead to eight equations for eight variables $`f_{R,L}^,(E,x=0,L)`$. We now solve these equations directly and discuss their physics. ### 2.2 Solution of the Boltzmann equation The 4 $`\times `$ 4 Boltzmann equation can block diagonalized into 2 $`\times `$ 2 blocks by changing variables to the charge and spin combinations $`X_{R,L}=f_{R,L}^{}+f_{R,L}^{}`$ and $`Y_{R,L}=f_{R,L}^{}f_{R,L}^{}`$. This spin-charge decoupling allows to solve exactly the Boltzmann equation. We find $$\frac{}{x}\left(\begin{array}{c}X_R\\ X_L\end{array}\right)=\frac{1}{l}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)\left(\begin{array}{c}X_R\\ X_L\end{array}\right)\text{ , and }\frac{}{x}\left(\begin{array}{c}Y_R\\ Y_L\end{array}\right)=\left(\begin{array}{cc}a& b\\ b& a\end{array}\right)\left(\begin{array}{c}Y_R\\ Y_L\end{array}\right),$$ (9) with $`l=1/(r+r_s^{})`$ the mean free path, and $`a=r+2r_s+r_s^{}`$, $`b=rr_s^{}`$. The 2 $`\times `$ 2 block equations can be easily integrated to obtain $$\left(\begin{array}{c}X_R(L)\\ X_L(L)\end{array}\right)=\left(\begin{array}{cc}1x& x\\ x& 1+x\end{array}\right)\left(\begin{array}{c}X_R(0)\\ X_L(0)\end{array}\right),$$ (10) with $`x=L/l`$. Similarly, $$\left(\begin{array}{c}Y_R(L)\\ Y_L(L)\end{array}\right)=\widehat{T}\left(\begin{array}{c}X_R(0)\\ X_L(0)\end{array}\right)\text{ , with }\widehat{T}=\left(\begin{array}{cc}t& u\\ u& \overline{t}\end{array}\right),$$ (11) where $`t=\mathrm{cosh}(\lambda L)\alpha \mathrm{sinh}(\lambda L)`$, $`\overline{t}=\mathrm{cosh}(\lambda L)+\alpha \mathrm{sinh}(\lambda L)`$, and $`u=\beta \mathrm{sinh}(\lambda L)`$. We used the notation $`\alpha =a/\lambda `$, $`\beta =b/\lambda `$, and $`\lambda =\sqrt{a^2b^2}`$. Next, we combine the boundary conditions Eqs. 3– 6 to Eq. 10 to obtain an expression for $`f_R^{}(E,0)f_R^{}(E,0)`$ and $`f_L^{}(E,0)f_L^{}(E,0)`$ as a function of $`f_R^{}(E,L)`$ and $`f_R^{}(E,L)`$. Once injected into Eq. 11, these relations lead to $`\widehat{M}\left(\begin{array}{c}f_R^{}(E,L)\\ f_R^{}(E,L)\end{array}\right)`$ $`=`$ $`2f_T(EeV)\widehat{T}\left(\begin{array}{c}T^{}+T^{}2T^{}T^{}\\ T^{}+T^{}\end{array}\right)`$ (16) $``$ $`f_T(E)\{(T^{}+T^{})\widehat{T}\left(\begin{array}{c}2(T^{}+T^{}T^{}T^{}1)x(T^{}+T^{}2T^{}T^{})\\ 2+T^{}+T^{}x(T^{}+T^{})\end{array}\right)`$ (19) $``$ $`(T^{}T^{})^2\left(\begin{array}{c}0\\ 1\end{array}\right)\}.`$ (22) The matrix $`\widehat{M}`$ appearing in the left hand side of Eq. 16 is $$\widehat{M}=\widehat{T}\left(\begin{array}{cc}A^{}& A^{}\\ B^{}& B^{}\end{array}\right)+(T^{}T^{})\left(\begin{array}{cc}1& 1\\ 1+T^{}& 1T^{}\end{array}\right),$$ (23) with the coefficients $`A^{}`$ $`=`$ $`3T^{}4T^{}T^{}+T^{}2(T^{})^2+2(T^{})^2T^{}+xT^{}\left[T^{}2T^{}T^{}+T^{}\right]`$ (24) $`B^{}`$ $`=`$ $`3T^{}+T^{}(T^{})^2T^{}T^{}+x\left[(T^{})^2+T^{}T^{}\right].`$ (25) The expression of $`A^{}`$ is obtained by exchanging $`T^{}`$ and $`T^{}`$ in Eq. 24. Similarly, $`B^{}`$ is obtained from $`B^{}`$ by exchanging $`T^{}`$ and $`T^{}`$ in Eq. 25. ### 2.3 Fully polarized limit We first consider the solution Eq. 16 in the case of fully polarized ferromagnets with high transparency contacts: $`H=0`$, $`k^{}=T^{}=0`$, $`k^{}=k`$, leading to $`T^{}=1`$. In this limit, only a spin-up current can enter the ferromagnet at $`x=L`$. This is expected on physical grounds, and it can be verified explicitly on the form Eq. 8 of the spin-down current. The total current is found to be $$I=\frac{e}{h}𝑑E\frac{2\left[f_T(EeV)f_T(E)\right]\mathrm{sinh}\left[2\sqrt{r_s(r_s+r)}L\right]}{(2+rL)\mathrm{sinh}\left[2\sqrt{r_s(r_s+r)}L\right]+\sqrt{\frac{r}{r_s}+1}\left\{\mathrm{cosh}\left[2\sqrt{r_s(r_s+r)}L\right]+1\right\}},$$ (26) where we considered only the forward scattering spin flip processes ($`r_s^{}=0`$), and assumed that $`r_sr`$, in which case the elastic mean free path $`l=1/r`$ is much below the spin-flip length $`l_{\mathrm{sf}}=1/[2\sqrt{r_s(r_s+r)}]`$. If $`L`$ is small compared to $`l_{\mathrm{sf}}`$, the conductance $`G2\frac{e^2}{h}r_sL`$ shows an inverse Drude behavior. ### 2.4 Spin polarization profile The spin polarization profile in the diffusive wire can be calculated in a straightforward fashion from the solution of the Boltzmann equation. Once we know the distribution functions at one extremity of the wire, we can use Eqs. 9 to propagate the solution to an arbitrary point. The resulting spin polarization inside the wire is proportional to the applied voltage, and is shown on Fig. 2 for various values of $`L`$. When $`L>l_{\mathrm{sf}}`$, there is a plateau in the spin polarization in the middle of the wire. In the opposite inverse Drude regime, there is no such plateau. ### 2.5 Effect of a partial spin polarization Now, we consider the effect of a partial spin polarization in the ferromagnet. It is expected on physical grounds that a decreasing spin polarization tends to suppress the inverse Drude scaling because this regime is clearly absent in the spin unpolarized case. This is visible on Fig. 3 where we plotted the conductance as a function of the length of the diffusive wire for decreasing spin polarizations. With an arbitrary polarization, there exists a critical length scale $`L_c`$ such that the conductance increases with $`L`$ below $`L_c`$, and decreases with $`L`$ above $`L_c`$. There exists also a critical value of $`T^{}`$ such that $`L_c=0`$ if $`T^{}>T_c^{}`$. To illustrate this, we have shown on Fig. 4 the variations of the critical length $`L_c`$ as a function of the parameter $`T^{}`$. When $`T^{}`$ increases, $`L_c`$ decreases: the maximum in $`G(L)`$ occurs for a smaller $`L_c`$. When $`T^{}`$ is above a critical value $`T_c^{}`$, the conductance decreases monotonically with $`L`$. We now describe the effect of a partial spin polarization on the basis of a small-$`T^{}`$ expansion. The strategy is to express the current to order $`L`$ and determine whether the conductance increases or decreases with $`L`$. We expand the current to first order in the two parameters $`K=\lambda L`$ and $`x=L/l`$, and retain the coefficients of this expansion to leading order in $`T^{}`$. It is first instructive to carry out the expansion with $`L=0`$, and therefore $`K=x=0`$. It is visible on Eqs. 78 that a prefactor $`T^{}`$ enters the spin-up current, and a prefactor $`T^{}`$ enters the spin-down current. We should then express $`f_R^{}(E,L)`$ to first order in $`T^{}`$ while $`f_R^{}(E,L)`$ should be expressed to order $`(T^{})^0`$. The spin-up and spin-down channels appear to play an asymmetric role. Nevertheless, the final expression of the conductance is identical in the spin-up and spin-down channels. An intermediate step in the calculation of $`f_R^{}(E,L)`$ is the derivation of $`\text{Det}\widehat{M}`$ to order $`T^{}`$ (see Eq. 23): $$\text{Det}\widehat{M}=4(T^{})^3\left\{1T^{}\left(\frac{2T^{}1}{T^{}}\right)\right\},$$ leading to an identical current in both spin channels: $$I^{}=I^{}=T^{}\frac{e}{h}\left[f_T(EeV)f_T(E)\right]𝑑E.$$ Now we consider a diffusive wire with a finite length $`L`$, and expand the current to first order in $`x`$ and $`K`$, and to leading order in $`T^{}`$. The determinant of the matrix $`\widehat{M}`$ in Eq. 16 is found to be $$\text{Det}\widehat{M}=4(T^{})^3\left\{1T^{}\left(\frac{2T^{}1}{T^{}}\right)\right\}4(T^{})^2(\alpha \beta )(2T^{})K4x(T^{})^2T^{}.$$ Next we expand the spin-up current to order $`L`$ to obtain $$I^{}=\frac{4(T^{})^3}{\text{Det}\widehat{M}}\left(T^{}(\alpha \beta )K\right)T^{}\left[1+K\frac{\alpha \beta }{T^{}}x\frac{T^{}}{T^{}}\right].$$ If $`T^{}`$ is small, the current increases with $`L`$ while it decreases with $`L`$ if $`T^{}`$ is large. The transition between these two behaviors is obtained for $`T_c^{}=\sqrt{(2r_s/r)T^{}}`$, compatible with the behavior shown on Fig. 4. ### 2.6 Replacement of one of the ferromagnets by a normal metal We now consider the situation where we replace the left-hand-side ferromagnet on Fig. 1 by a normal metal. In the presence of high transparency contacts, the conductance of this junction is of order $`e^2/h`$ in the absence of diffusion while it is of order $`(e^2/h)T^{}`$ in the spin valve geometry on Fig. 1. Replacing one of the ferromagnets by a normal metal is expected to suppress the inverse Drude scaling. The boundary conditions appropriate to describe this situation are $`f_R^{}(E,0)`$ $`=`$ $`Tf_T(EeV)+(1T)f_L^{}(E,0)`$ (27) $`f_R^{}(E,0)`$ $`=`$ $`Tf_T(EeV)+(1T)f_L^{}(E,0)`$ (28) $`f_L^{}(E,L)`$ $`=`$ $`T^{}f_T(E)+(1T^{})f_R^{}(E,0)`$ (29) $`f_L^{}(E,L)`$ $`=`$ $`T^{}f_T(E)+(1T^{})f_R^{}(E,0),`$ (30) that should be solved together with Eqs. 9. The solution is found to be $$\widehat{N}\left(\begin{array}{c}f_R^{}(E,L)\\ f_R^{}(E,L)\end{array}\right)=\left(\begin{array}{c}(T^{}+T^{})(1T+xT)\\ (T^{}T^{})(u+t(1T))\end{array}\right)f_T(E)+2T\left(\begin{array}{c}1\\ 0\end{array}\right)f_T(EeV),$$ with $$\widehat{N}=\left(\begin{array}{cc}C^{}& C^{}\\ D^{}& D^{}\end{array}\right),$$ and $`C^{}`$ $`=`$ $`T+T^{}TT^{}+xTT^{}`$ (31) $`D^{}`$ $`=`$ $`\overline{t}u(1T)(1T^{})(t(1T)+u).`$ (32) We have plotted on Fig. 5 the conductance of the junction with high transparency contacts, where it is visible that the conductance decreases monotonically with the length of the diffusive wire. Now, reducing the contact transparency restores a regime in which the conductance increases with the size of the diffusive wire. This is visible on Fig. 6 where we used $`T=0.01`$ and $`T^{}=1`$. Again, the inverse Drude scaling is obtained for the smallest values of $`T^{}`$ (with strongly polarized magnets). ## 3 Quantum coherent transport: a single magnetic impurity “spin-flip Fabry Perot interferometer” ### 3.1 Matching equations We now consider a single magnetic impurity at $`x=0`$ in a normal metal, in the presence of a normal metal – ferromagnet interface at $`x=a`$ (see Fig. 7). The purpose of this calculation is to study a model in which the interplay between multiple reflections and phase coherence is treated exactly, and to determine whether there exists a signature of spin accumulation in the quantum coherent regime. We find that the quantum model behaves like a Fabry Perot interferometer, with a resonance in the spin flip channels. This can be viewed as the signature of spin accumulation in the quantum coherent regime. We neglect the Kondo effect because we want to describe a situation in which the temperature is above the Kondo temperature. The conduction electrons are scattered through the Hamiltonian $`=V_0+V_1𝐒_i.𝐬`$, where $`𝐒_i`$ is the impurity spin and $`𝐬`$ the spin of the conduction electron, and we further assume a single channel geometry. This type of model has been used by Zhu and Wang to investigate the effect of a magnetic impurity close to a normal metal – superconductor interface. The spin-up and spin-down wave functions are grouped in a two-component spinor $`\widehat{\psi }(x)`$. Clearly, the impurity couples the spin-up and spin-down wave functions. The matching of the wave function at the impurity site reads $$\widehat{\psi }(0^+)=\widehat{\psi }(0^{})\text{, and }\frac{\widehat{\psi }}{x}(0^+)\frac{\widehat{\psi }}{x}(0^{})=\frac{2m}{\mathrm{}^2}\left[\lambda \widehat{1}+\mu \widehat{\sigma }_x\right]\widehat{\psi }(0),$$ (33) with $`\lambda =V_0V_1/4`$ and $`\mu =V_1/2`$. The matching of the wave function at the ferromagnet boundary reads $$\widehat{\psi }(a^+)=\widehat{\psi }(a^{})\text{, and }\frac{\widehat{\psi }}{x}(a^+)\frac{\widehat{\psi }}{x}(a^{})=\frac{2m}{\mathrm{}^2}H\widehat{\psi }(a),$$ (34) where we included a repulsive interface potential $`H\delta (xa)`$ at the normal metal – ferromagnet interface. Eqs. 3334 generate eight constraints, for a set of eight transmission coefficients. This calculation amounts to a resummation to all orders of a series of diagrams in which a conduction electron scatters onto the impurity, scatters back onto the interface, scatters again onto the impurity, … (see Fig. 8 (a)). Note that the diagram with a hole in the intermediate state shown on Fig. 8 (b) generates another series which is not included in the calculation. If one wanted to describe the Kondo effect close to a ferromagnet interface, it would be crucial to incorporate the diagram on Fig. 8 (b), as well as inserting the interface scattering in this diagram. ### 3.2 Scattering in the total spin $`S^z=0`$ sectors #### 3.2.1 Incoming electron with a spin-up We first consider a spin-up electron incoming on the interface while the impurity is supposed to have initially a spin down. The wave functions are $`\widehat{\psi }_i^e(x)`$ $`=`$ $`\left(\begin{array}{c}1\\ 0\end{array}\right)e^{ikx}+\left(\begin{array}{c}b_i^{ee}\\ b_i^{ee}\end{array}\right)e^{ikx}\text{ if }x<0\text{.}`$ (39) $`\widehat{\psi }_i^e(x)`$ $`=`$ $`\left(\begin{array}{c}\alpha \\ \alpha ^{}\end{array}\right)e^{ikx}+\left(\begin{array}{c}\beta \\ \beta ^{}\end{array}\right)e^{ikx}\text{ if }0<x<a\text{.}`$ (44) $`\widehat{\psi }_i^e(x)`$ $`=`$ $`t_i^{ee}\left(\begin{array}{c}1\\ 0\end{array}\right)e^{ik^{}x}+t_i^{ee}\left(\begin{array}{c}0\\ 1\end{array}\right)e^{ik^{}x}\text{ if }x>a\text{,}`$ (49) where $`k^{}`$ and $`k^{}`$ denote the spin-up and spin-down Fermi wave vectors in the ferromagnet. In the notation of the transmission coefficients, the superscript denotes the initial and final spin orientations of the conduction electron while the subscript denotes the initial orientation of the impurity. The solution of the matching equations is straightforward, and we find the transmission coefficients $`t_i^{ee}`$ $`=`$ $`{\displaystyle \frac{1}{𝒟A^{}}}\left[\overline{A}(1+iz)X^{}+AizY^{}\right]`$ (50) $`t_i^{ee}`$ $`=`$ $`{\displaystyle \frac{1}{𝒟A^{}}}iz^{}\left[\overline{A}X^{}+AY^{}\right].`$ (51) We used the notation $`X^\sigma =\frac{1}{2}+iZ+\frac{Z}{2Z^\sigma }`$, $`Y^\sigma =\frac{1}{2}\left(iZ+\frac{Z}{2Z^\sigma }\right)`$, $`A=\mathrm{exp}(ika)`$, $`A^\sigma =\mathrm{exp}(ik^\sigma a)`$. The dimensionless scattering potentials in Eqs. 5051 are $`z=m\lambda /(\mathrm{}^2k)`$ and $`z^{}=m\mu /(\mathrm{}^2k)`$ at the impurity site, and $`Z=mH/(\mathrm{}^2k)`$, $`Z^\sigma =mH/(\mathrm{}^2k^\sigma )`$ at the normal metal – ferromagnet interface. The denominator $`𝒟`$ in Eqs. 5051 is $$𝒟=X^{}X^{}(\overline{A})^2[1z^2+(z^{})^2+2iz]+(X^{}Y^{}+X^{}Y^{})[z^2+(z^{})^2+iz]+Y^{}Y^{}A^2[z^2+(z^{})^2].$$ (52) #### 3.2.2 Incoming electron with a spin-down We consider now an incoming electron with a spin-down while the impurity has initially a spin-up. The wave functions are $`\widehat{\psi }_i^e(x)`$ $`=`$ $`\left(\begin{array}{c}0\\ 1\end{array}\right)e^{ikx}+\left(\begin{array}{c}b_i^{ee}\\ b_i^{ee}\end{array}\right)e^{ikx}\text{ if }x<0\text{.}`$ (57) $`\widehat{\psi }_i^e(x)`$ $`=`$ $`\left(\begin{array}{c}\alpha ^{}\\ \alpha \end{array}\right)e^{ikx}+\left(\begin{array}{c}\beta ^{}\\ \beta \end{array}\right)e^{ikx}\text{ if }0<x<a\text{.}`$ (62) $`\widehat{\psi }_i^e(x)`$ $`=`$ $`t_i^{ee}\left(\begin{array}{c}0\\ 1\end{array}\right)e^{ik^{}x}+t_i^{ee}\left(\begin{array}{c}1\\ 0\end{array}\right)e^{ik^{}x}\text{ if }x>a\text{.}`$ (67) The equations for $`t_i^{ee}`$ and $`t_i^{ee}`$ are obtained from the ones in section 3.2.1 under the transformation $`A^{}A^{}`$, and $`Z^{}Z^{}`$. The amplitude for transmission in the ferromagnet is $`t_i^{ee}`$ $`=`$ $`{\displaystyle \frac{1}{𝒟A^{}}}iz^{}[\overline{A}X^{}+AY^{}]`$ (68) $`t_i^{ee}`$ $`=`$ $`{\displaystyle \frac{1}{𝒟A^{}}}[\overline{A}(1+iz)X^{}+AizY^{}].`$ (69) ### 3.3 Scattering in the total spin $`S^z=\pm 1`$ sectors The incoming electron does not undergo spin-flip scattering in the sectors with a total spin $`S^z=\pm 1`$. The transmission coefficients in the sector $`S^z=1`$ is found to be $$t_i^{ee}=\frac{1}{A^{}[\overline{A}\left(1+i(z+z^{})\right)X^{}+Ai(z+z^{})Y^{}]}.$$ In the sector $`S^z=1`$, we have $$t_i^{ee}=\frac{1}{A^{}[\overline{A}\left(1+i(z+z^{})\right)X^{}+Ai(z+z^{})Y^{}]}.$$ We can check easily that these forms of the transmission coefficients are identical to Eq. 50, with $`z^{}=0`$, and the replacement $`zz+z^{}`$. This is expected since there is no spin-dependent scattering in the limit $`z^{}=0`$ of Eq. 50. ### 3.4 Landauer formula We now evaluate the total conductance and assume that the incoming electron and impurity do not have any preferential direction. The conductance is the sum of four terms, weighted by the probability $`𝒫=1/2`$ to have a spin-up or spin-down impurity: $`G=\left(G_i^e+G_i^e+G_i^e+G_i^e\right)/2`$, with $`G_i^e`$ $`=`$ $`{\displaystyle \frac{e^2}{h}}\left({\displaystyle \frac{k^{}}{k}}\|t_i^{ee}\|^2+{\displaystyle \frac{\text{Re}k^{}}{k}}\|t_i^{ee}\|^2\right)`$ (70) $`G_i^e`$ $`=`$ $`{\displaystyle \frac{e^2}{h}}\left({\displaystyle \frac{k^{}}{k}}\|t_i^{ee}\|^2+{\displaystyle \frac{\text{Re}k^{}}{k}}\|t_i^{ee}\|^2\right)`$ (71) $`G_i^e`$ $`=`$ $`{\displaystyle \frac{e^2}{h}}{\displaystyle \frac{k^{}}{k}}\|t_i^{ee}\|^2`$ (72) $`G_i^e`$ $`=`$ $`{\displaystyle \frac{e^2}{h}}{\displaystyle \frac{\text{Re}k^{}}{k}}\|t_i^{ee}\|^2.`$ (73) We have incorporated the possibility of having a pure imaginary wave vector $`k^{}`$, corresponding to an empty spin-down band. ### 3.5 Resonances We consider the presence of a strong interface scattering at the metal – ferromagnet interface. The electron are multiply reflected before they enter the ferromagnet, and therefore the resonator has a high quality factor. We first choose the parameter $`a`$ in such a way that spin-flip scattering is resonant. As it is visible on Figs. 9 and 10, the presence of a spin polarization in the ferromagnet generates a resonance in the conductance, not present in the unpolarized situation. We have shown the conductance with $`V_0=0`$ but a similar behavior has been obtained with a finite $`V_0`$. The values of $`ka`$ for which a resonance occurs can be worked out by calculating the transmission coefficients in the limit of a large $`z`$, $`z^{}`$. In this limit, we find $`t_i^{ee}`$ $`=`$ $`{\displaystyle \frac{iz^{}\overline{A}^{}}{((z^{})^2z^2)}}{\displaystyle \frac{1}{\overline{A}X^{}+AY^{}}}`$ (74) $`t_i^{ee}`$ $`=`$ $`{\displaystyle \frac{iz\overline{A}^{}}{(z^{})^2z^2}}{\displaystyle \frac{1}{\overline{A}X^{}+AY^{}}}.`$ (75) The resonances occur when $`\mathrm{tan}(ka)=1/(i2Z)`$. For a large $`Z`$, the resonances are close to the real axis: $`\mathrm{tan}(ka)=1/(2Z)`$, in agreement with Fig. 11. The reason why the resonance appear to be sharp as a function of $`a`$ when $`Z`$ is large is that, even without spin flip scattering, the quality factor of such a resonator is large when $`Z`$ is large. The occurrence of a parameter range in which a peak occurs in the spin flip conductance is intriguing. The presence of a specific physics in the spin flip channels can be already understood from the large-$`z`$, $`z^{}`$ behavior, Eqs. 7475. Typically, one has $`\|t_i^{ee}\|^2(z^{})^2/((z^{})^2z^2)^2`$ and $`\|t_i^{ee}\|^2(z)^2/((z^{})^2z^2)^2`$. In the presence of spin flip scattering, one has $`zz^{}`$ and therefore a different conductance in the spin flip and non spin flip channels. ## 4 Conclusions To conclude, we have determined to what extend spin accumulation can result in an inverse Drude behavior in a semi classical spin valve model. Our treatment was based on an exact decoupling between the charge and spin sectors of the Boltzmann equation. We have addressed a similar question in a single impurity quantum model and found the existence of a resonance in the spin flip conductance. It is an open question to determine the quantum coherent behavior of a spin valve with a finite concentration of impurities. ## Appendix A Derivation of the Boltzmann equation with spin-flip scattering We give a brief derivation of the Boltzmann equation Eq. 1 in the presence of a spin-flip scattering potential. The derivation generalizes Ref. to incorporate a spin-flip scattering self energy. The Dyson equation in the spin tensor Keldysh space reads $`(\widehat{G}_0^1\widehat{\mathrm{\Sigma }})(1,2)\widehat{G}(2)=\delta (12)`$. The convolution includes a sum over coordinates, time, and spin. The kinetic equation is obtained from the difference of the Keldysh components of the Dyson equations and its conjugate: $$[\widehat{G}_0^1\text{Re}\widehat{\mathrm{\Sigma }},\widehat{G}^K]_{}[\widehat{\mathrm{\Sigma }}^K,\text{Re}\widehat{G}]_{}=\frac{i}{2}[\widehat{\mathrm{\Sigma }}^K,\widehat{A}]_+\frac{i}{2}[\widehat{\mathrm{\Gamma }},\widehat{G}^K]_+,$$ (76) with $`\left[\right]_{}`$ and $`\left[\right]_+`$ denoting a commutator and an anticommutator respectively. We refer the reader to Ref. for an explanation of the symbols used in Eq. 76. We use the the self energy shown on Fig. 12: $$\widehat{\mathrm{\Sigma }}_\sigma (𝐩,𝐑,T)=n_i\frac{d𝐩^{}}{(2\pi )^3}\|v(𝐩𝐩^{})\|^2\widehat{G}_{\sigma ,\sigma }(𝐩^{},𝐑,T)+n_i^{}\frac{d𝐩^{}}{(2\pi )^3}\|v^{}(𝐩𝐩^{})\|^2\widehat{G}_{\sigma ,\sigma }(𝐩^{},𝐑,T),$$ (77) with $`n_i`$ and $`n_i^{}`$ the concentration of non magnetic and magnetic impurities. The first term in Eq. 77 describes non spin-flip scattering, and the second term describes spin-flip scattering. Notice that this self energy does not incorporate the Kondo effect since we do not incorporate the possibility of a having hole in the intermediate state. We assume the self energy in Eq. 76 to be constant in space, and use the gradient expansion to first order $$\left(AB\right)_{\sigma ,\sigma ^{}}(X,p)\underset{\sigma _2}{}\left[1+\frac{i}{2}\left(_X^A_p^B_p^A_X^B\right)\right]A_{\sigma ,\sigma _2}B_{\sigma _2,\sigma ^{}}.$$ If $`A`$ and $`B`$ are symmetric in spin $`A_{\sigma ,\sigma ^{}}=A_{\sigma ^{},\sigma }`$, and $`B_{\sigma ,\sigma ^{}}=B_{\sigma ^{},\sigma }`$, the commutator reduces to the usual spinless Poisson bracket: $`\left[AB\right]_{,\sigma ,\sigma ^{}}=i_{\sigma _2}\{A_{\sigma ,\sigma _2},B_{\sigma 2,\sigma }\}`$, with $`\{A,B\}=_X^AA_p^BB_p^AA_X^BB`$. Using these relations, we expand the kinetic equation Eq. 76 and integrate over energy to obtain the Boltzmann equation $`_Tf_{𝐩,\sigma }+_𝐩\xi _𝐩_𝐑f_{𝐩,\sigma }_𝐑U_𝐩f_{𝐩,\sigma }`$ $`=`$ $`2\pi n_i{\displaystyle \frac{d𝐩^{}}{(2\pi )^3}\|v(𝐩𝐩^{})\|^2\delta (\xi _𝐩\xi _𝐩^{})[f_{𝐩^{},\sigma }f_{𝐩,\sigma }]}`$ $`+2\pi n_i^{}{\displaystyle \frac{d𝐩^{}}{(2\pi )^3}\|v^{}(𝐩𝐩^{})\|^2\delta (\xi _𝐩\xi _𝐩^{})[f_{𝐩^{},\sigma }f_{𝐩,\sigma }]}.`$ In the one dimensional limit, Eq. A reduces to the Boltzmann equation Eq. 1, with the scattering coefficients related to the $`q=0`$ and $`q=2k_f`$ components of the scattering potential: $`r=n_i\|v_{2k_f}\|^2`$, $`r_s=n_i^{}\|v_0\|^2`$, and $`r_s^{}=n_i^{}\|v_{2k_f}^{}\|^2`$.
warning/0004/quant-ph0004028.html	ar5iv	text	# Quantum Computer as a Probabilistic Inference Engine ## Abstract We propose a new class of quantum computing algorithms which generalize many standard ones. The goal of our algorithms is to estimate probability distributions. Such estimates are useful in, for example, applications of Decision Theory and Artificial Intelligence, where inferences are made based on uncertain knowledge. The class of algorithms that we propose is based on a construction method that generalizes a Fredkin-Toffoli (F-T) construction method used in the field of classical reversible computing. F-T showed how, given any binary deterministic circuit, one can construct another binary deterministic circuit which does the same calculations in a reversible manner. We show how, given any classical stochastic network (classical Bayesian net), one can construct a quantum network (quantum Bayesian net). By running this quantum Bayesian net on a quantum computer, one can calculate any conditional probability that one would be interested in calculating for the original classical Bayesian net. Thus, we generalize the F-T construction method so that it can be applied to any classical stochastic circuit, not just binary deterministic ones. We also show that, in certain situations, our class of algorithms can be combined with Grover’s algorithm to great advantage. ## 1 Introduction In this paper, we use the language of classical Bayesian (CB) and quantum Bayesian (QB) nets. The reader is expected to possess a rudimentary command of this language. We begin this paper with a review of various standard quantum computing algorithms; namely, those due to Deutsch-Jozsa, Simon, Bernstein-Vazirani, and Grover. We discuss these standard algorithms both in terms of qubit circuits (the conventional approach) and QB nets. Then we propose a class of quantum computing algorithms which generalizes the standard ones. Most standard quantum computing algorithms are designed for calculating deterministic or almost deterministic probability distributions. (By a deterministic probability distribution we mean one whose range is restricted to either zero or unit probabilities.) In contrast, our algorithms can also estimate more general probability distributions. Such estimates are useful in, for example, applications of Decision Theory and Artificial Intelligence, where inferences are made based on uncertain knowledge. Since some of the standard algorithms are contained in the class of algorithms that we propose, some algorithms in our class have a time-complexity advantage over the best classical algorithms for performing the same task. Even those algorithms in our class that have no complexity advantage might still be useful for nanoscale quantum computing because they are reversible and thus dissipate less power. Power dissipation is best minimized in nanoscale devices since it can lead to serious performance degradation. The class of algorithms that we propose in this paper is based on a construction method that generalizes a Fredkin-Toffoli (F-T) construction method used in the field of classical reversible computing. F-T showed in Refs. how, given any binary gate $`f`$ (i.e., a function $`f:\{0,1\}^r\{0,1\}^s`$, for some integers $`r,s`$), one can construct another binary gate $`\overline{f}`$ such that $`\overline{f}`$ is a deterministic reversible extension (DRE) of $`f`$. $`\overline{f}`$ can be used to perform the same calculations as $`f`$, but in a reversible manner. Binary gates $`f`$ and $`\overline{f}`$ can be represented as binary deterministic circuits. In this paper, we show how, given any CB net $`𝒩^C`$, one can construct a QB net $`𝒩^Q`$ which is a “q-embedding” (q=quantum) of $`𝒩^C`$. By running $`𝒩^Q`$ on a quantum computer, one can calculate any conditional probability that one would be interested in calculating for the CB net $`𝒩^C`$. Our method for constructing a q-embedding for a CB net is a generalization of the F-T method for constructing a DRE of a binary deterministic circuit. Thus, we generalize their method so that it applies to any classical stochastic circuit, not just binary deterministic ones. A quantum compiler can “compile” a unitary matrix; i.e., it can express the matrix as a SEO (sequence of elementary operations) that a quantum computer can understand. To run a QB net on a quantum computer, we need to replace the QB net by an equivalent SEO. This can be done with the help of a quantum compiler. Thus, the class of algorithms that we propose promises to be fertile ground for the use of quantum compilers. In certain cases, the probabilities that we wish to find are too small to be measurable by running $`𝒩^Q`$ on a quantum computer. However, we will show that sometimes it is possible to define a new QB net, call it $`𝒩_{}^{Q}{}_{}{}^{}`$, that magnifies and makes measurable the probabilities that were unmeasurable using $`𝒩^Q`$ alone. We will refer to $`𝒩_{}^{Q}{}_{}{}^{}`$ as Grover’s Microscope for $`𝒩^Q`$, because $`𝒩_{}^{Q}{}_{}{}^{}`$ is closely related to Grover’s algorithm, and it magnifies the probabilities found with $`𝒩^Q`$. ## 2 Notation and Other Preliminaries In this section, we will introduce certain notation that is used throughout the paper. We will use the word “ditto” as follows. If we say “A (ditto, X) is smaller than B (ditto, Y)”, we mean “A is smaller than B” and “X is smaller than Y”. Let $`Bool=\{0,1\}`$. For integers $`a`$ and $`b`$ such that $`ab`$, let $`Z_{a,b}=\{a,a+1,a+2,\mathrm{}b\}`$. For any statement $`𝒮`$, we define the truth function $`\theta (𝒮)`$ to equal 1 if $`𝒮`$ is true and 0 if $`𝒮`$ is false. For example, $`\theta (x>0)`$ represents the unit step function and $`\delta (x,y)=\theta (x=y)`$ the Kronecker delta function. $``$ will denote addition mod 2. For any integer $`x`$, $`x\%2`$ will mean the remainder from dividing $`x`$ by 2. For example, $`4\%2=0`$ and $`5\%2=1`$. (This same $`\%`$ notation is used in the C programming language.) When speaking of bits with states 0 and 1, we will often use an overbar to represent the opposite state: $`\overline{0}=1`$, $`\overline{1}=0`$. Note that if $`x,kBool`$ then $$\underset{k}{}(1)^{kx}=1+(1)^x=2\delta (x,0).$$ (1) If $`\stackrel{}{x},\stackrel{}{y}Bool^n`$, we will use $`\stackrel{}{x}\stackrel{}{y}=_{\alpha =0}^{n1}x_\alpha y_\alpha `$, where the addition is normal, not mod 2. Given $`xZ_{0,\mathrm{}}`$, let $`x=_{\alpha =0}^{\mathrm{}}x_\alpha 2^\alpha `$, where $`x_\alpha Bool`$ for all $`\alpha `$. Then we will denote the binary representation of $`x`$ by $`bin(x)=(x_0,x_1,x_2,\mathrm{})`$. Thus, $`bin_\alpha (x)=x_\alpha `$. On the other hand, given $`\stackrel{}{x}=(x_0,x_1,x_2,\mathrm{})Bool^{\mathrm{}}`$, let $`x=_{\alpha =0}^{\mathrm{}}x_\alpha 2^\alpha `$. Then we will denote the decimal representation of $`\stackrel{}{x}`$ by $`dec(\stackrel{}{x})=x`$. We will use the symbol $`_{}`$ to denote a sum of whatever is on the right hand side of this symbol over those indices with a dot underneath them. For example, $`_{}f(\text{},\text{},c)=_{a,b}f(a,b,c)`$. Furthermore, $`_{all}`$ will denote a sum over all indices. If we wish to exclude a particular index from the summation, we will indicate this by a slash followed by the name of the index. For example, in $`_{all/a,b}`$ we wish to exclude summation over $`a`$ and $`b`$. Suppose $`f`$ maps set $`S`$ into the complex numbers. We will often use $`\frac{f(x)}{_xnum}`$ to represent $`\frac{f(x)}{_{xS}f(x)}`$. Thus, $`num`$ is shorthand for the numerator of the fraction. We will underline random variables. $`P(\underset{¯}{a}=a)=P_{\underset{¯}{a}}(a)`$ will denote the probability that the random variable $`\underset{¯}{a}`$ assumes value $`a`$. $`P(\underset{¯}{a}=a)`$ will often be abbreviated by $`P(a)`$ when no confusion will arise. $`S_{\underset{¯}{a}}`$ will denote the set of values which the random variable $`\underset{¯}{a}`$ may assume, and $`N_{\underset{¯}{a}}`$ will denote the number of elements in $`S_{\underset{¯}{a}}`$. $`pd(B\|A)`$ will stand for the set of probability distributions $`P(\|)`$ such that $`P(b\|a)0`$ and $`_{b^{}B}P(b^{}\|a)=1`$ for all $`aA`$ and $`bB`$. This paper will also utilize certain notation and nomenclature associated with classical and quantum Bayesian nets. For example, we will use $`(x_{})_A`$ to denote $`\{x_i:iA\}`$. See Ref. for a review of such notation. $`H_1=\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right)`$ is the one bit Hadamard matrix. $`H_{N_B}=H_1^{N_B}`$ (the n-fold tensor product of $`H_1`$) is the $`N_B`$ bit Hadamard matrix. We will also use $`\widehat{H}_1=\frac{1}{\sqrt{2}}H_1`$ and $`\widehat{H}_{N_B}=\widehat{H}_1^{N_B}=\frac{1}{\sqrt{2^{N_B}}}H_1^{N_B}`$. Note that $`(H_1)_{b,b^{}}=(1)^{bb^{}}`$ for $`b,b^{}Bool`$, and $`(H_{N_B})_{\stackrel{}{b},\stackrel{}{b^{}}}=(1)^{\stackrel{}{b}\stackrel{}{b^{}}}`$ for $`\stackrel{}{b},\stackrel{}{b^{}}Bool^{N_B}`$. Any $`2\times 2`$ matrix $`M`$ which acts on bit $`\alpha `$ will be denoted by $`M(\alpha )`$. (We like to use lower case Greek letters for bit labels.) In this notation, a controlled-not (cnot) gate with control bit $`\kappa `$ and target bit $`\tau `$ can be expressed as $`\sigma _x(\tau )^{n(\kappa )}`$. See Ref. for more details about this notation. Let $`\stackrel{}{\kappa }=(\kappa _0,\kappa _1,\mathrm{},\kappa _{N_B1})`$ label $`N_B`$ bits. Assume all $`\kappa _i`$ are distinct. We will often use $`N_S=2^{N_B}`$, where $`N_B`$ stands for number of bits and $`N_S`$ for number of states. If $`\|\varphi _{\kappa _i}=\|\varphi (\kappa _i)`$ is a ket for qubit $`\kappa _i`$, define $`\|\varphi _\stackrel{}{\kappa }=\|\varphi (\stackrel{}{\kappa })=_{i=0}^{N_B1}\|\varphi (\kappa _i)`$. For example, if $$\|0_{\kappa _i}=\left(\begin{array}{c}1\\ 0\end{array}\right)$$ (2) for all $`i`$, then $$\|0_\stackrel{}{\kappa }=\underset{i=0}{\overset{N_B1}{}}\|0_{\kappa _i}=\left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right)\mathrm{}\left(\begin{array}{c}1\\ 0\end{array}\right)=[1,0,0,\mathrm{},0]^T.$$ (3) Likewise, if $`\mathrm{\Omega }(\kappa _i)`$ is an operator acting on qubit $`\kappa _i`$, define $`\mathrm{\Omega }(\stackrel{}{\kappa })=_{i=0}^{N_B1}\mathrm{\Omega }(\kappa _i)`$. For example, $`H_1(\stackrel{}{\kappa })=_{i=0}^{N_B1}H_1(\kappa _i)`$ is an $`N_B`$ bit Hadamard matrix. Next, we will introduce some notation related to Pauli matrices. The Pauli matrices are given by: $$\sigma _x=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _y=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\sigma _z=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$ (4) If $`\|+_z`$ and $`\|_z`$ represent the eigenvectors of $`\sigma _z`$ with eigenvalues $`+1`$ and $`1`$, respectively, then we define $$\|0=\|+_z=\left(\begin{array}{c}1\\ 0\end{array}\right),$$ (5) and $$\|1=\|_z=\left(\begin{array}{c}0\\ 1\end{array}\right).$$ (6) We denote the “number operator” by $`n`$. Thus $$n=\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)=\|_z_z\|=\frac{1\sigma _z}{2},$$ (7) and $$\overline{n}=1n=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)=\|+_z+_z\|=\frac{1+\sigma _z}{2}.$$ (8) Since $`n`$ and $`\sigma _z`$ are diagonal, it is easy to see that $$(1)^n=\sigma _z.$$ (9) It is also useful to introduce symbols for the projectors with respect to $`\|0`$ and $`\|1`$; $$P_0^z=\|00\|=\overline{n},$$ (10) $$P_1^z=\|11\|=n.$$ (11) Most of the definitions and results stated so far for $`\sigma _z`$ have counterparts for $`\sigma _x`$ and $`\sigma _y`$. The counterpart results can be easily proven by applying a rotation that interchanges the coordinate axes. Let $`w\{x,y,z\}`$. If $`\|+_w`$ and $`\|_w`$ represent the eigenvectors of $`\sigma _w`$ with eigenvalues $`+1`$ and $`1`$, respectively, then we define $$\|0_w=\|+_w,$$ (12) and $$\|1_w=\|_w.$$ (13) Let $$n_w=\|_w_w\|=\frac{1\sigma _w}{2},$$ (14) $$\overline{n}_w=1n_w=\|+_w+_w\|=\frac{1+\sigma _w}{2}.$$ (15) As when $`w=z`$, one has $$(1)^{n_w}=\sigma _w.$$ (16) Let $$P_0^w=\|0_w0_w\|=\overline{n}_w,$$ (17) and $$P_1^w=\|1_w1_w\|=n_w.$$ (18) In understanding Grover’s algorithm, it is helpful to be aware of some simple properties of reflections on a plane. Suppose $`\varphi `$ is a normalized ($`\varphi ^{}\varphi =1`$) complex vector. Define the projection and reflection operators for $`\varphi `$ by $$\mathrm{\Pi }_\varphi =\varphi \varphi ^{},R_\varphi =12\mathrm{\Pi }_\varphi .$$ (19) Note that $`\mathrm{\Pi }_\varphi ^2=\mathrm{\Pi }_\varphi `$. Fig.1 shows that if $`x^{}=R_\varphi x`$, then $`x^{}`$ is the reflection of $`x`$ with respect to the plane perpendicular to $`\varphi `$. For example, $`R_\varphi \varphi =\varphi `$. Some simple properties of $`R_\varphi `$ are as follows. $`R_\varphi =R_\varphi ^{}`$ and $`R_\varphi R_\varphi ^{}=R_\varphi ^2=1`$. Since reflections are unitary matrices, a product of reflections is also a unitary matrix. Note that $`(1)^{\mathrm{\Pi }_\varphi }`$ $`=`$ $`e^{i\pi \mathrm{\Pi }_\varphi }=1+(e^{i\pi \mathrm{\Pi }_\varphi }1)`$ (20a) $`=`$ $`1+\mathrm{\Pi }_\varphi (e^{i\pi }1)=12\mathrm{\Pi }_\varphi `$ (20b) $`=`$ $`R_\varphi .`$ (20c) (Eq.(20b) follows from the Taylor expansion of $`e^{i\pi \mathrm{\Pi }_\varphi }`$.) If $`e_1,e_2,\mathrm{},e_n`$ is an orthonormal basis for a vector space, $`\mathrm{\Pi }_i=e_ie_i^{}`$, and $`R_i=12\mathrm{\Pi }_i`$, then the product of the $`R_i`$ in any order is $`1`$. Indeed, $`R_1R_2\mathrm{}R_n`$ $`=`$ $`(12\mathrm{\Pi }_1)(12\mathrm{\Pi }_2)\mathrm{}(12\mathrm{\Pi }_n)`$ (21a) $`=`$ $`12(\mathrm{\Pi }_1+\mathrm{\Pi }_2+\mathrm{}\mathrm{\Pi }_n)`$ (21b) $`=`$ $`1.`$ (21c) Another property of reflection operators which is useful for understanding Grover’s algorithm is the following. Let $$e_0=\left(\begin{array}{c}1\\ 0\end{array}\right),e_1=\left(\begin{array}{c}0\\ 1\end{array}\right).$$ (22) Now suppose that $`e_1^{}`$ is obtained by rotating $`e_1`$ clockwise by an angle $`\theta /2`$: $$e_1^{}=\left(\begin{array}{cc}\mathrm{cos}(\theta /2)& \mathrm{sin}(\theta /2)\\ \mathrm{sin}(\theta /2)& \mathrm{cos}(\theta /2)\end{array}\right)e_1=\left(\begin{array}{c}\mathrm{sin}(\theta /2)\\ \mathrm{cos}(\theta /2)\end{array}\right).$$ (23) $`e_1^{}e_1`$ for small $`\theta `$. It is easy to check that the double reflection $`R_{e_1^{}}R_{e_0}`$ is equivalent to a rotation (also clockwise) by $`\theta `$: $$R_{e_1^{}}R_{e_0}=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right).$$ (24) (That these two successive reflections equal a rotation was to be expected, since the reflections are orthogonal matrices and a product of orthogonal matrices is itself orthogonal.) Above, we have considered plane reflections $`R_\varphi `$ acting on a complex vector space, but our formulas still hold true when $`R_\varphi `$ acts on a real instead of a complex vector space. In the case of real vector spaces, the Hermitian conjugate symbol $``$ is replaced by the matrix transpose symbol $`T`$, and unitary matrices are replaced by orthogonal matrices. ## 3 Some Standard Algorithms Next we will discuss several standard algorithms that are considered among the best that the quantum computation field has to offer at the present time. Later, we will try to generalize these standard algorithms. ### 3.1 Deutsch-Jozsa Algorithm In this section we will discuss the D-J (Deutsch-Jozsa) algorithm. We will do this first in terms of qubit circuits (the conventional approach), and then in terms of QB nets. Let $`\stackrel{}{\kappa }=(\kappa _0,\kappa _1,\mathrm{},\kappa _{N_B1})`$ label $`N_B`$ “control” bits and let $`\tau `$ label a single “target” bit. Assume that $`\tau `$ and all the $`\kappa _i`$ are distinct. We will denote the state of these bits in the preferred basis (the eigenvectors of $`\sigma _z`$) by $`\|x_\stackrel{}{\kappa }\|y_\tau `$, where $`xBool^{N_B}`$ and $`yBool`$. Given a function $`f:Bool^{N_B}Bool`$, define the unitary operator $`\mathrm{\Omega }`$ by $$\mathrm{\Omega }=\sigma _x(\tau )\widehat{H}_1(\tau )\widehat{H}_1(\stackrel{}{\kappa })\sigma _x^{f(\stackrel{}{n}(\stackrel{}{\kappa }))}(\tau )\widehat{H}_1(\stackrel{}{\kappa })\widehat{H}_1(\tau )\sigma _x(\tau ),$$ (25) where $`\stackrel{}{n}(\stackrel{}{\kappa })=(n(\kappa _0),n(\kappa _1),\mathrm{},n(\kappa _{N_B1}))`$. The operation $`\sigma _x^{f(\stackrel{}{n}(\stackrel{}{\kappa }))}(\tau )`$, because it depends on $`f`$, is often called an “oracle” and each use of it is called a “query”. The right hand side of Eq.(25) may be represented by the circuit diagram shown in Fig.2. The D-J algorithm consists of applying $`\mathrm{\Omega }`$ to an initial state $`\|0_\stackrel{}{\kappa }\|0_\tau `$ of bits $`\stackrel{}{\kappa }`$ and $`\tau `$, and then measuring the final state of these bits in the preferred basis. Fig.2 and the right hand side of Eq.(25) are two equivalent ways of representing a particular SEO. There are infinitely many SEOs that yield $`\mathrm{\Omega }`$. Fig.2 is just one of them. In fact, the original D-J paper gave a different SEO for $`\mathrm{\Omega }`$, one with two queries instead of one. For $`XBool^{N_B},YBool`$, let $$\|\psi _0=\|X_\stackrel{}{\kappa }\|Y_\tau ,$$ (26) and $$\|\psi _i=\mathrm{\Omega }_i\|\psi _{i1}\mathrm{for}i=1,2,\mathrm{},$$ (27) where $$\mathrm{\Omega }_1=\widehat{H}_1(\stackrel{}{\kappa })\widehat{H}_1(\tau )\sigma _x(\tau ),$$ (28) $$\mathrm{\Omega }_2=\sigma _x^{f(\stackrel{}{n}(\stackrel{}{\kappa }))}(\tau ),$$ (29) and $$\mathrm{\Omega }_3=\mathrm{\Omega }_1^{}.$$ (30) Then it is easy to show using simple identities (such as $`(H_1)_{b,b^{}}=(1)^{bb^{}}`$, $`\overline{0}=1`$, $`\overline{1}=0`$, and $`(1)^b=(1)^b`$ for $`b,b^{}Bool`$) that $$\|\psi _1=\frac{1}{\sqrt{2^{N_B+1}}}\underset{x,y}{}(1)^{xX+y\overline{Y}}\|x_\stackrel{}{\kappa }\|y_\tau ,$$ (31) $$\|\psi _2=\frac{1}{\sqrt{2^{N_B+1}}}\underset{x,y}{}(1)^{xX+y\overline{Y}}\|x_\stackrel{}{\kappa }\|yf(x)_\tau ,$$ (32) $$\|\psi _3=\frac{1}{2^{N_B+1}}\underset{x,y,X^{},Y^{}}{}(1)^{x(X^{}X)+y(\overline{Y^{}}\overline{Y})+\overline{Y^{}}f(x)}\|X^{}_\stackrel{}{\kappa }\|Y^{}_\tau .$$ (33) Applying $`X^{},Y^{}\|`$ to the right hand side of Eq.(33) and using the identity Eq.(1) finally yields: $$X^{},Y^{}\|\mathrm{\Omega }\|X,Y=\delta (Y^{},Y)\frac{1}{2^{N_B}}\underset{xBool^{N_B}}{}(1)^{x(X^{}X)+\overline{Y^{}}f(x)}$$ (34) for all $`X^{},XBool^{N_B}`$ and $`Y^{},YBool`$. Thus, if the initial states of $`\stackrel{}{\kappa }`$ and $`\tau `$ are $`\underset{¯}{X}=0`$ and $`\underset{¯}{Y}=0`$, then the probability of obtaining $`\underset{¯}{X}^{}=X^{}`$ for the final state of $`\stackrel{}{\kappa }`$ is $`P(X^{}\|\underset{¯}{X}=\underset{¯}{Y}=0)`$ $`=`$ $`{\displaystyle \underset{Y^{}}{}}\|X^{},Y^{}\|\mathrm{\Omega }\|X=0,Y=0\|^2`$ (35) $`=`$ $`{\displaystyle \frac{1}{4^{N_B}}}\|{\displaystyle \underset{x}{}}(1)^{xX^{}+f(x)}\|^2.`$ Let $`_{bal}`$, the set of “balanced” functions, be the set of all $`f:Bool^{N_B}Bool`$ such that $`f`$ maps exactly half of its domain to zero and half to one. Let $`_{con}`$, the set of “constant” functions, be the set of all $`f:Bool^{N_B}Bool`$ such that $`f`$ maps all its domain to zero or all of it to one. From Eq.(35), if $`X^{}=0`$ and $`f_{bal}_{con}`$, then $$P(\underset{¯}{X}^{}=0\|\underset{¯}{X}=\underset{¯}{Y}=0)=\{\begin{array}{cc}1\hfill & \mathrm{if}f_{con}\hfill \\ 0\hfill & \mathrm{if}f_{bal}\hfill \end{array}.$$ (36) Now consider the QB net defined by Fig.3 and Table 3.1. \| nodes \| states \| amplitudes \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{X}`$ \| $`XBool^{N_B}`$ \| $`\delta (X,0)`$ \| \| \| $`\underset{¯}{Y}`$ \| $`YBool`$ \| $`\delta (Y,0)`$ \| \| \| $`\underset{¯}{x}`$ \| $`xBool^{N_B}`$ \| $`(1)^{xX}/\sqrt{2^{N_B}}`$ \| \| \| $`\underset{¯}{y}`$ \| $`yBool`$ \| $`(1)^{y\overline{Y}}/\sqrt{2}`$ \| \| \| $`\underset{¯}{c}`$ \| $`c=(c_x,c_y),c_xBool^{N_B},c_yBool`$ \| $`\delta (c_x,x)\delta (c_y,yf(x))`$ \| \| \| $`\underset{¯}{x}^{}`$ \| $`x^{}Bool^{N_B}`$ \| $`\delta (x^{},c_x)`$ \| \| \| $`\underset{¯}{y}^{}`$ \| $`y^{}Bool`$ \| $`\delta (y^{},c_y)`$ \| \| \| $`\underset{¯}{X}^{}`$ \| $`X^{}Bool^{N_B}`$ \| $`(1)^{X^{}x^{}}/\sqrt{2^{N_B}}`$ \| \| \| $`\underset{¯}{Y}^{}`$ \| $`Y^{}Bool`$ \| $`(1)^{\overline{Y^{}}y^{}}/\sqrt{2}`$ \| \| Table 1 For this net, the amplitude $`A(x.)`$ of net story $`x.`$ is the product of all the terms in the third column of Table 3.1. If $`\underset{¯}{X}=0`$ and $`\underset{¯}{Y}=0`$, then the probability of obtaining $`\underset{¯}{X}^{}=X^{}`$ is $$P(X^{}\|\underset{¯}{X}=\underset{¯}{Y}=0)=\frac{_Y^{}\|_{all/X^{}Y^{},X,Y}A(x.)\|_{X=Y=0}\|^2}{_X^{}num},$$ (37) where $`A(x.)`$ on the right hand side is evaluated at $`X=Y=0`$. Substituting the value of $`A(x.)`$ into Eq.(37) immediately yields Eq.(35). Note that one can calculate the probability distribution Eq.(35) by means of a CB net instead of a QB net. One can do this with the CB net defined by the graph $`\underset{¯}{X}^{}\underset{¯}{Y}^{}`$, with: \| nodes \| states \| probabilities \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{X}^{}`$ \| $`X^{}Bool^{N_B}`$ \| $`P_{\underset{¯}{X}^{}}(X^{})`$ \| \| \| $`\underset{¯}{Y}^{}`$ \| $`Y^{}Bool`$ \| $`P_{\underset{¯}{Y}^{}\|\underset{¯}{X}^{}}(Y^{}\|X^{})`$ \| \| Table 2 where $`P_{\underset{¯}{X}^{}}`$ and $`P_{\underset{¯}{Y}^{}\|\underset{¯}{X}^{}}`$ are calculated from $$P_{\underset{¯}{X}^{},\underset{¯}{Y}^{}}(X^{},Y^{})=\|X^{},Y^{}\|\mathrm{\Omega }\|X=0,Y=0\|^2.$$ (38) We will say that the CB net defined by the graph $`\underset{¯}{X}^{}\underset{¯}{Y}^{}`$ and Table 3.1 is “q-embedded” in the QB net defined by Fig.3 and Table 3.1. In subsequent sections, we will say much more about q-embedding of CB nets. ### 3.2 Simon’s Algorithm In this section we will discuss Simon’s algorithm. We will do this first in terms of qubit circuits (the conventional approach), and then in terms of QB nets. Simon’s algorithm uses $`N_B`$ “control” bits, just like the D-J algorithm. However, it uses $`N_B`$ target bits whereas the D-J algorithm uses only one. Simon’s algorithm deals with a vector-valued function $`f:Bool^{N_B}Bool^{N_B}`$, whereas D-J’s algorithm deals with a scalar-valued function $`f:Bool^{N_B}Bool`$. Let $`\stackrel{}{\kappa }=(\kappa _0,\kappa _1,\mathrm{},\kappa _{N_B1})`$ label $`N_B`$ “control” bits and let $`\stackrel{}{\tau }=(\tau _0,\tau _1,\mathrm{},\tau _{N_B1})`$ label $`N_B`$ “target” bits. Assume all $`\tau _i`$ and $`\kappa _i`$ are distinct. We will denote the state of these bits in the preferred basis (the eigenvectors of $`\sigma _z`$) by $`\|x_\stackrel{}{\kappa }\|y_\stackrel{}{\tau }`$, where $`xBool^{N_B}`$ and $`yBool^{N_B}`$. Given a function $`f=(f_0,f_1,\mathrm{},f_{N_B1})`$ where $`f_i:Bool^{N_B}Bool`$, define the unitary operator $`\mathrm{\Omega }`$ by $$\mathrm{\Omega }=\widehat{H}_1(\stackrel{}{\kappa })\left[\underset{i=0}{\overset{N_B1}{}}\sigma _x^{f_i(\stackrel{}{n}(\stackrel{}{\kappa }))}(\tau _i)\right]\widehat{H}_1(\stackrel{}{\kappa }).$$ (39) The operator $`\mathrm{\Omega }`$ for Simon’s algorithm is analogous to the $`\mathrm{\Omega }`$ defined by Eq.(25) for the D-J algorithm. The right hand side of Eq.(39) may be represented by the circuit diagram of Fig.4. Simon’s algorithm consists of applying $`\mathrm{\Omega }`$ given by Eq.(39) to an initial state $`\|0_\stackrel{}{\kappa }\|0_\stackrel{}{\tau }`$ of bits $`\stackrel{}{\kappa }`$ and $`\stackrel{}{\tau }`$, and then measuring the final state of these bits in the preferred basis. One performs this routine several times. The measurement outcomes allow one to determine the period of the function $`f`$ if $`f`$ is of a special periodic type that will be specified later. Using the same techniques that we used to evaluate the matrix elements of $`\mathrm{\Omega }`$ for the D-J algorithm, one finds $$X^{},Y^{}\|\mathrm{\Omega }\|X,Y=\frac{1}{2^{N_B}}\underset{xBool^{N_B}}{}(1)^{x(X^{}X)}\delta (Y^{},Yf(x)),$$ (40) for all $`X^{},Y^{},X,YBool^{N_B}`$. If the initial states of $`\stackrel{}{\kappa }`$ and $`\stackrel{}{\tau }`$ are $`\underset{¯}{X}=0`$ and $`\underset{¯}{Y}=0`$, then the probability of obtaining $`\underset{¯}{X}^{}=X^{}`$ for the final state of $`\stackrel{}{\kappa }`$ is $`P(X^{}\|\underset{¯}{X}=\underset{¯}{Y}=0)`$ $`=`$ $`{\displaystyle \underset{Y^{}}{}}\|X^{},Y^{}\|\mathrm{\Omega }\|X=0,Y=0\|^2`$ (41) $`=`$ $`{\displaystyle \frac{1}{4^{N_B}}}{\displaystyle \underset{Y^{}}{}}\|{\displaystyle \underset{x}{}}(1)^{xX^{}}\delta (Y^{},f(x))\|^2.`$ Now suppose $`_S`$ is the set of those functions $`f:Bool^{N_B}Bool^{N_B}`$ such that $`f`$ is 2 to 1 (i.e., $`f`$ maps exactly two domain points into each image point) and has a “period” $`\mathrm{\Delta }`$. By a period $`\mathrm{\Delta }`$, we mean a non-zero element of $`Bool^{N_B}`$ such that $`f(x)=f(x\mathrm{\Delta })`$ for all $`xBool^{N_B}`$. For any $`f_S`$ and any $`yBool^{N_B}`$, there exist exactly two elements of $`Bool^{N_B}`$, call them $`x_1`$ and $`x_2`$, such that $`x_1=x_2\mathrm{\Delta }`$ and $`f(x_1)=f(x_2)=y`$. Call $`f_p^1(y)`$ one of these $`x`$ values, and call $`f_p^1(y)\mathrm{\Delta }`$ the other. (The $`p`$ subscript stands for “principal part”, in analogy with Complex Analysis.) If $`f_S`$, and $`I(f)`$ is the image of $`f`$, then $$\delta (Y^{},f(x))=\{\begin{array}{cc}\delta (f_p^1(Y^{}),x)+\delta (f_p^1(Y^{})\mathrm{\Delta },x),\hfill & \mathrm{if}Y^{}I(f)\hfill \\ 0\hfill & \mathrm{otherwise}\hfill \end{array}.$$ (42) Substituting this expression for $`\delta (Y^{},f(x))`$ into Eq.(41) and using Eq.(1) yields $$P(X^{}\|\underset{¯}{X}=\underset{¯}{Y}=0)=\frac{1}{2^{N_B1}}\delta (X^{}\mathrm{\Delta },0).$$ (43) To calculate the period $`\mathrm{\Delta }`$ of $`f`$, run the experiment $`\nu `$ times, measuring $`X^{}`$ each time. Let $`X^{}(i)`$ represent the $`i`$th measurement outcome. Then, for sufficiently large $`\nu `$, one can find $`\mathrm{\Delta }`$ by solving the equations $`X^{}(1)\mathrm{\Delta }=0`$, $`X^{}(2)\mathrm{\Delta }=0`$, … , $`X^{}(\nu )\mathrm{\Delta }=0`$. Now consider the QB net defined by Fig.5 and Table 3.2. \| nodes \| states \| amplitudes \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{X}`$ \| $`XBool^{N_B}`$ \| $`\delta (X,0)`$ \| \| \| $`\underset{¯}{Y}`$ \| $`YBool^{N_B}`$ \| $`\delta (Y,0)`$ \| \| \| $`\underset{¯}{x}`$ \| $`xBool^{N_B}`$ \| $`(1)^{xX}/\sqrt{2^{N_B}}`$ \| \| \| $`\underset{¯}{c}`$ \| $`c=(c_x,c_y);c_x,c_yBool^{N_B}`$ \| $`\delta (c_x,x)\delta (c_y,Yf(x))`$ \| \| \| $`\underset{¯}{x}^{}`$ \| $`x^{}Bool^{N_B}`$ \| $`\delta (x^{},c_x)`$ \| \| \| $`\underset{¯}{X}^{}`$ \| $`X^{}Bool^{N_B}`$ \| $`(1)^{X^{}x^{}}/\sqrt{2^{N_B}}`$ \| \| \| $`\underset{¯}{Y}^{}`$ \| $`Y^{}Bool^{N_B}`$ \| $`\delta (Y^{},c_y)`$ \| \| Table 3 For this net, the amplitude $`A(x.)`$ of net story $`x.`$ is the product of all the terms in the third column of Table 3.2. If $`\underset{¯}{X}=0`$ and $`\underset{¯}{Y}=0`$, then the probability of obtaining $`\underset{¯}{X}^{}=X^{}`$ is $$P(X^{}\|\underset{¯}{X}=\underset{¯}{Y}=0)=\frac{_Y^{}\|_{all/X^{}Y^{},X,Y}A(x.)\|_{X=Y=0}\|^2}{_X^{}num},$$ (44) where $`A(x.)`$ on the right hand side is evaluated at $`X=Y=0`$. Substituting the value of $`A(x.)`$ into Eq.(44) immediately yields Eq.(41). It is possible to calculate the probability distribution Eq.(41) by means of a CB net instead of a QB net. One can do this with the CB net defined by the graph $`\underset{¯}{X}^{}\underset{¯}{Y}^{}`$, with: \| nodes \| states \| probabilities \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{X}^{}`$ \| $`X^{}Bool^{N_B}`$ \| $`P_{\underset{¯}{X}^{}}(X^{})`$ \| \| \| $`\underset{¯}{Y}^{}`$ \| $`Y^{}Bool^{N_B}`$ \| $`P_{\underset{¯}{Y}^{}\|\underset{¯}{X}^{}}(Y^{}\|X^{})`$ \| \| Table 4 where $`P_{\underset{¯}{X}^{}}`$ and $`P_{\underset{¯}{Y}^{}\|\underset{¯}{X}^{}}`$ are calculated from $$P_{\underset{¯}{X}^{},\underset{¯}{Y}^{}}(X^{},Y^{})=\|X^{},Y^{}\|\mathrm{\Omega }\|X=0,Y=0\|^2.$$ (45) We will say that the CB net defined by the graph $`\underset{¯}{X}^{}\underset{¯}{Y}^{}`$ and Table 3.2 is “q-embedded” in the QB net defined by Fig.5 and Table 3.2. ### 3.3 Bernstein-Vazirani Algorithm In this section we will discuss the B-V (Bernstein-Vazirani) algorithm. To understand the B-V algorithm, it is helpful to first establish the following simple single qubit identities. First note that the single qubit Hadamard matrix rotates the $`Z`$-direction number operator into the $`X`$-direction number operator: $$\widehat{H}_1n_z\widehat{H}_1=\frac{1}{2}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)=n_x.$$ (46) Thus, $$\sigma _x^b=[(1)^{n_x}]^b=(1)^{bn_x}=\widehat{H}_1(1)^{bn_z}\widehat{H}_1.$$ (47) Next note that $`\sigma _x`$ exchanges the components of any vector it acts on: $$\sigma _x\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)=\left(\begin{array}{c}\beta \\ \alpha \end{array}\right),$$ (48) for any complex numbers $`\alpha ,\beta `$. In particular, if $`bBool`$, then $$\sigma _x^b\|0=\|b.$$ (49) Now we are ready to discuss the B-V algorithm. Let $`\stackrel{}{\kappa }=(\kappa _0,\kappa _1,\mathrm{},\kappa _{N_B1})`$ label $`N_B`$ “control” bits and let $`\tau `$ label a single “target” bit. Assume that $`\tau `$ and all the $`\kappa _i`$ are distinct. We will denote the state of these bits in the preferred basis (the eigenvectors of $`\sigma _z`$) by $`\|x_\stackrel{}{\kappa }\|y_\tau `$, where $`xBool^{N_B}`$ and $`yBool`$. For $`\stackrel{}{b}Bool^{N_B}`$, define the unitary operator $$\omega _\stackrel{}{b}=\underset{i=0}{\overset{N_B1}{}}\sigma _x(\kappa _i)^{b_i}.$$ (50) The B-V algorithm is simply the following multi-qubit generalization of Eq.(49) $$\omega _\stackrel{}{b}\|0_\stackrel{}{\kappa }=\|\stackrel{}{b}_\stackrel{}{\kappa }.$$ (51) That’s all there is to B-V! Eq.(51) can be represented by a qubit circuit consisting of a single wire for $`\stackrel{}{\kappa }`$, with a single node representing $`\omega _\stackrel{}{b}`$. Eq.(51) can also be represented by a QB net defined by the graph $`\underset{¯}{X}\underset{¯}{X}^{}`$, with \| nodes \| states \| amplitudes \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{X}`$ \| $`X=(X_0,X_1,X_{N_B1})Bool^{N_B}`$ \| $`\delta (X,0)`$ \| \| \| $`\underset{¯}{X}^{}`$ \| $`X^{}=(X_0^{},X_1^{},X_{N_B1}^{})Bool^{N_B}`$ \| $`_{i=0}^{N_B1}\delta ^{b_i}(X_i^{},\overline{X_i})`$ \| \| Table 5 We should mention that it is common in the literature to dress up and obfuscate Eq.(50) as follows. By virtue of Eq.(47), one can re-express $`\omega _\stackrel{}{b}`$ as $$\omega _\stackrel{}{b}=\underset{i=0}{\overset{N_B1}{}}\sigma _x(\kappa _i)^{b_i}=\widehat{H}_1(\stackrel{}{\kappa })(1)^{_{i=0}^{N_B1}b_in_z(\kappa _i)}\widehat{H}_1(\stackrel{}{\kappa }).$$ (52) Some workers ascend to an even higher peak of obfuscation by adding a totally unnecessary target qubit. They define an operator, call it $`\mathrm{\Omega }_\stackrel{}{b}`$, obtained by replacing the $`(1)`$ in Eq.(52) by the operator $`\sigma _x(\tau )`$ acting on a target qubit $`\tau `$: $$\mathrm{\Omega }_\stackrel{}{b}=\widehat{H}_1(\stackrel{}{\kappa })[\sigma _x(\tau )]^{_{i=0}^{N_B1}b_in_z(\kappa _i)}\widehat{H}_1(\stackrel{}{\kappa }).$$ (53) At the beginning of the experiment, they put the target qubit in a state which is an eigenvector of $`\sigma _x(\tau )`$ with eigenvalue $`1`$. Thus, the obfuscated version of the B-V algorithm with a target qubit can be summarized by $$\mathrm{\Omega }_\stackrel{}{b}\|_x_\tau \|0_\stackrel{}{\kappa }=\omega _\stackrel{}{b}\|_x_\tau \|0_\stackrel{}{\kappa }=\|_x_\tau \|\stackrel{}{b}_\stackrel{}{\kappa }.$$ (54) We emphasize that for the B-V algorithm, the target qubit is a totally unnecessary affectation. So far we have given an unconventional presentation of the B-V algorithm. For completeness, we now give a conventional one. Define $$\|\psi _0=\|0_\stackrel{}{\kappa }\|_x_\tau ,$$ (55) and $$\|\psi _i=\mathrm{\Omega }_i\|\psi _{i1}\mathrm{for}i=1,2,\mathrm{},$$ (56) where $$\mathrm{\Omega }_1=\widehat{H}_1(\stackrel{}{\kappa }),$$ (57) $$\mathrm{\Omega }_2=[\sigma _x(\tau )]^{_{i=0}^{N_B1}b_in_z(\kappa _i)},$$ (58) and $$\mathrm{\Omega }_3=\mathrm{\Omega }_1.$$ (59) It follows that $$\|\psi _1=\frac{1}{\sqrt{2^{N_B}}}\underset{\stackrel{}{x}Bool^{N_B}}{}\|\stackrel{}{x}_\stackrel{}{\kappa }\|_x_\tau ,$$ (60) $$\|\psi _2=\frac{1}{\sqrt{2^{N_B}}}\underset{\stackrel{}{x}}{}(1)^{\stackrel{}{b}\stackrel{}{x}}\|\stackrel{}{x}_\stackrel{}{\kappa }\|_x_\tau ,$$ (61) and $`\|\psi _3`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2^{N_B}}}}{\displaystyle \underset{\stackrel{}{x}}{}}(1)^{\stackrel{}{b}\stackrel{}{x}}{\displaystyle \frac{1}{\sqrt{2^{N_B}}}}{\displaystyle \underset{\stackrel{}{y}}{}}(1)^{\stackrel{}{y}\stackrel{}{x}}\|\stackrel{}{y}_\stackrel{}{\kappa }\|_x_\tau `$ (62a) $`=`$ $`{\displaystyle \frac{1}{2^{N_B}}}{\displaystyle \underset{\stackrel{}{x},\stackrel{}{y}}{}}(1)^{(\stackrel{}{b}\stackrel{}{y})\stackrel{}{x}}\|\stackrel{}{y}_\stackrel{}{\kappa }\|_x_\tau `$ (62b) $`=`$ $`\|\stackrel{}{b}_\stackrel{}{\kappa }\|_x_\tau .`$ (62c) To go from step (b) to step (c) of Eq.(62), we used the orthogonality property given by Eq.(1). ### 3.4 Grover’s Algorithm In this section we will discuss Grover’s algorithm . Let $`\stackrel{}{\kappa }=(\kappa _0,\kappa _1,\mathrm{},\kappa _{N_B1})`$ label $`N_B`$ bits. Assume all $`\kappa _i`$ are distinct. We begin by defining the following $`N_S`$-dimensional column vectors: $$\|\mu _\stackrel{}{\kappa }=\mu =\mu _{N_S}=\frac{1}{\sqrt{N_S}}[1,1,1,1\mathrm{},1]^T=\frac{1}{\sqrt{N_S}}H_{N_B}[1,0,0,0,\mathrm{},0]^T,$$ (63) $$\|\varphi _\stackrel{}{\kappa }=\varphi =[0,\mathrm{},0,0,1,0,0,\mathrm{},0]^T.$$ (64) All components of $`\varphi `$ are zero except for one predetermined component, located at position $`j_{targ}Z_{0,N_S1}`$, which equals one. We will refer to $`j_{targ}`$ as the target state (not to be confused with a target qubit). Note that we chose a special basis (or, equivalently, a special matrix representation) from the start. Note that $`\varphi \|\mu =\frac{1}{\sqrt{N_S}}`$, so $`\mu `$ and $`\varphi `$ are nearly orthogonal for large $`N_S`$. It is also convenient to define the component-wise negation of $`\varphi `$: $$\|\varphi ^{not}_\stackrel{}{\kappa }=\varphi ^{not}=[1,\mathrm{},1,1,0,1,1,\mathrm{},1]^T.$$ (65) Note that $`\varphi ^{not}`$ is not normalized. Define projection and reflection operators for $`\mu `$ and $`\varphi `$: $$\mathrm{\Pi }_\mu =\|\mu \mu \|,R_\mu =12\mathrm{\Pi }_\mu =(1)^{\mathrm{\Pi }_\mu },$$ (66) and $$\mathrm{\Pi }_\varphi =\|\varphi \varphi \|,R_\varphi =12\mathrm{\Pi }_\varphi =(1)^{\mathrm{\Pi }_\varphi }.$$ (67) Grover’s algorithm can be summarized by the following equation: $$(R_\mu R_\varphi )^r\mu \varphi ,$$ (68) for some integer $`r`$ to be determined, where “$``$” means approximation at large $`N_S`$. Thus, starting with an $`N_B`$ qubit system in a state $`\mu `$, one applies the operator $`(R_\mu R_\varphi )`$ consecutively $`r`$ times, so that the $`N_B`$ qubit system ends in a state as close to $`\varphi `$ as possible. Measuring state $`\varphi `$ in the special basis yields the target state $`j_{targ}`$. Eq.(68) can be represented by a qubit circuit consisting of a single wire for $`\stackrel{}{\kappa }`$, with $`r`$ nodes, each representing $`R_\mu R_\varphi `$. Eq.(68) can also be represented by a QB net defined by a Markov chain graph $`\underset{¯}{X}_0\underset{¯}{X}_1\underset{¯}{X}_2\mathrm{}\underset{¯}{X}_{r1}`$, with \| nodes \| states \| amplitudes \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{X}_0`$ \| $`X_0Bool^{N_B}`$ \| $`\delta (X_0,0)`$ \| \| \| $`\underset{¯}{X}_i`$ for $`iZ_{1,r1}`$ \| $`X_iBool^{N_B}`$ \| $`X_i\|(R_\mu R_\varphi )\|X_{i1}`$ \| \| Table 6 To find the optimum number $`r`$ of iterations, one can proceed as follows. First, notice that Eq.(68) describes a process which is entirely confined to the vector subspace spanned by $`\mu `$ and $`\varphi `$. Since $`\mu `$ and $`\varphi `$ are not orthogonal, it is convenient to define an orthonormal basis $`e_0,e_1`$ for the space $`span(\mu ,\varphi )`$. Let $$e_0=\varphi ,e_1=\frac{\varphi ^{not}}{\sqrt{N_S1}}.$$ (69) Then $$\mu =\frac{1}{\sqrt{N_S}}(e_0+\sqrt{N_S1}e_1).$$ (70) Fig.6 portrays various vectors that arise in explaining Grover’s algorithm. Since we plan to stay within the two dimensional vector space with orthonormal basis $`e_0,e_1`$, it is convenient to switch matrix representations. Within $`span(e_0,e_1)`$, $`e_0,e_1`$ can be represented more simply by: $$e_0=\left(\begin{array}{c}1\\ 0\end{array}\right),e_1=\left(\begin{array}{c}0\\ 1\end{array}\right).$$ (71) If $`e_0,e_1`$ are represented in this way, then $$\varphi =\left(\begin{array}{c}1\\ 0\end{array}\right),\mu =\frac{1}{\sqrt{N_S}}\left(\begin{array}{c}1\\ \sqrt{N_S1}\end{array}\right),$$ (72) and $$R_\mu R_\varphi =\frac{1}{N_S}\left(\begin{array}{cc}N_S2& 2\sqrt{N_S1}\\ 2\sqrt{N_S1}& N_S2\end{array}\right).$$ (73) Thus, $$R_\mu R_\varphi =\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right),$$ (74) where $$\mathrm{sin}\theta =\frac{2\sqrt{N_S1}}{N_S}\frac{2}{\sqrt{N_S}}.$$ (75) Eq.(74) is just Eq.(24) with $`e_1^{}=\mu `$ and $`e_0=\varphi `$. It follows that $$(R_\mu R_\varphi )^r=\left(\begin{array}{cc}\mathrm{cos}(r\theta )& \mathrm{sin}(r\theta )\\ \mathrm{sin}(r\theta )& \mathrm{cos}(r\theta )\end{array}\right),$$ (76) and $`(R_\mu R_\varphi )^r\mu `$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}(r\theta )& \mathrm{sin}(r\theta )\\ \mathrm{sin}(r\theta )& \mathrm{cos}(r\theta )\end{array}\right){\displaystyle \frac{1}{\sqrt{N_S}}}\left(\begin{array}{c}1\\ \sqrt{N_S1}\end{array}\right)`$ (81) $``$ $`\left(\begin{array}{c}\mathrm{sin}(r\theta )\\ \mathrm{cos}(r\theta )\end{array}\right).`$ (84) We want the final state of the system to be parallel or anti-parallel to $`e_0=\varphi `$ ; therefore, we want $$\left(\begin{array}{c}\mathrm{sin}(r\theta )\\ \mathrm{cos}(r\theta )\end{array}\right)\left(\begin{array}{c}\pm 1\\ 0\end{array}\right).$$ (85) This will occur if $$r\theta \frac{\pi }{2}(1+2k),r\frac{\pi }{4}(1+2k)\sqrt{N_S}$$ (86) for some integer $`k`$. Note that, in Grover’s algorithm, the number of “queries” (calls to a unitary matrix that depends on $`\varphi `$) is far from unique. To illustrate this, let $`Q`$ be a permutation matrix that satisfies $$Q\varphi =\|0=[1,0,0,\mathrm{},0]^T.$$ (87) Since all the components of $`\mu `$ are equal, $`Q\mu =\mu `$. Thus $$(R_\mu R_\varphi )^r\mu =Q^T(R_\mu R_{\|0})^rQ\mu =Q^T(R_\mu R_{\|0})^r\mu .$$ (88) Hence, it is possible to accomplish the full Grover transformation of $`\mu `$ with only a single query $`Q^T`$. Since $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{c}a\\ b\end{array}\right)=\left(\begin{array}{c}b\\ a\end{array}\right)`$, the matrix $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ is just a clockwise rotation by $`\pi /2`$. Let $`U_{Grov}`$ $`=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ $`=`$ $`e_1e_0^T+e_0e_1^T`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_S1}}}[\varphi ^{not}\varphi ^T+\varphi (\varphi ^{not})^T].`$ (92) Note that $`U_{Grov}\mu `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_S1}}}\left(\varphi ^{not}[\varphi ^T\mu ]+\varphi [(\varphi ^{not})^T\mu ]\right)`$ (93) $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_S(N_S1)}}}[\varphi ^{not}+(N_S1)\varphi ].`$ From the point of view of quantum compiling, what Grover found is that the $`\pi /2`$ rotation $`U_{Grov}`$ is (approximately) equal to the $`r`$-fold product of $`R_\mu R_\varphi `$, where $`R_\mu R_\varphi `$ can be shown to have a SEO of low (polynomial in $`N_B`$) complexity. Grover’s algorithm has been modified in various, minor ways since it was first published. For example, Brassard et al. pointed out in Ref. that the vector $`\mu `$ need not be the vector whose components are all equal. Other vectors $`\mu `$ will do just as well. Another modification of Grover’s algorithm due to Younes-Miller adds an extra qubit to the original $`N_B`$ qubits. Next we will discuss the Younes-Miller modification of Grover’s algorithm, because it resembles a modification of Grover’s algorithm that we will use in a future section. Let $`\stackrel{}{\kappa }=(\kappa _0,\kappa _1,\mathrm{},\kappa _{N_B1})`$ label $`N_B`$ bits. Let $`\tau `$ label a single bit. Assume $`\tau `$ and all the $`\kappa _i`$ are distinct. Let $`\mu `$ and $`\varphi `$ denote the same $`N_S`$ dimensional column vectors that we used in discussing the original Grover algorithm. In addition, define the following $`2N_S`$ dimensional column vectors: $$\|\stackrel{~}{\mu }=\|+_z_\tau \|\mu _\stackrel{}{\kappa }=\left(\begin{array}{c}1\\ 0\end{array}\right)\mu _{N_S}=\left(\begin{array}{c}\mu _{N_S}\\ 0\end{array}\right),$$ (94) $$\|\stackrel{~}{\varphi }=\|_x_\tau \|\varphi _\stackrel{}{\kappa }=\frac{1}{\sqrt{2}}\left(\begin{array}{c}1\\ 1\end{array}\right)\varphi =\frac{1}{\sqrt{2}}\left(\begin{array}{c}\varphi \\ \varphi \end{array}\right).$$ (95) Note that $`\stackrel{~}{\varphi }\|\stackrel{~}{\mu }=\frac{1}{\sqrt{2N_S}}`$, so $`\stackrel{~}{\varphi }`$ and $`\stackrel{~}{\mu }`$ are nearly orthogonal for large $`N_S`$. Define projection and reflection operators for $`\stackrel{~}{\varphi }`$ in the usual way: $$\mathrm{\Pi }_{\stackrel{~}{\varphi }}=\|\stackrel{~}{\varphi }\stackrel{~}{\varphi }\|,R_{\stackrel{~}{\varphi }}=12\mathrm{\Pi }_{\stackrel{~}{\varphi }}=12\mathrm{\Pi }_{\|\varphi _\stackrel{}{\kappa }}\mathrm{\Pi }_{\|_x_\tau }.$$ (96) $`R_{\stackrel{~}{\varphi }}`$ can be re-expressed as $`R_{\stackrel{~}{\varphi }}`$ $`=`$ $`1+\mathrm{\Pi }_{\|\varphi _\stackrel{}{\kappa }}(\sigma _x(\tau )1)=\mathrm{exp}[\mathrm{\Pi }_{\|\varphi _\stackrel{}{\kappa }}\mathrm{ln}\sigma _x(\tau )]=`$ (97) $`=`$ $`[\sigma _x(\tau )]^{\mathrm{\Pi }_{\|\varphi _\stackrel{}{\kappa }}}.`$ Define projection and reflection operators for $`\stackrel{~}{\mu }`$ in the usual way: $$\mathrm{\Pi }_{\stackrel{~}{\mu }}=\|\stackrel{~}{\mu }\stackrel{~}{\mu }\|,R_{\stackrel{~}{\mu }}=12\mathrm{\Pi }_{\stackrel{~}{\mu }}=12\mathrm{\Pi }_{\|\mu _\stackrel{}{\kappa }}\mathrm{\Pi }_{\|+_z_\tau }.$$ (98) $`R_{\stackrel{~}{\mu }}`$ can be re-expressed as $$R_{\stackrel{~}{\mu }}=\widehat{H}_1(\stackrel{}{\kappa })\left(12\mathrm{\Pi }_{\|0_\stackrel{}{\kappa }}\mathrm{\Pi }_{\|0_\tau }\right)\widehat{H}_1(\stackrel{}{\kappa }).$$ (99) In analogy with the original Grover’s algorithm, the Younes-Miller version can be summarized by $$(R_{\stackrel{~}{\mu }}R_{\stackrel{~}{\varphi }})^r\|\stackrel{~}{\mu }\|\stackrel{~}{\varphi },$$ (100) for some integer $`r`$ to be determined, where “$``$” means approximation at large $`N_S`$. Thus, starting with an $`N_B+1`$ qubit system in a state $`\stackrel{~}{\mu }`$, one applies the operator $`(R_\mu R_\varphi )`$ consecutively $`r`$ times, so that the final state of the $`N_B+1`$ qubit system ends in a state as close to $`\stackrel{~}{\varphi }`$ as possible. Measuring state $`\stackrel{~}{\varphi }`$ in the special basis yields the target state $`j_{targ}`$. To find the optimum number $`r`$ of iterations, one can proceed as follows. First, notice that Eq.(100) describes a process which is entirely confined to the vector subspace spanned by $`\stackrel{~}{\mu }`$ and $`\stackrel{~}{\varphi }`$. Since $`\stackrel{~}{\mu }`$ and $`\stackrel{~}{\varphi }`$ are not orthogonal, it is convenient to define an orthonormal basis $`e_0,e_1`$ for the space $`span(\stackrel{~}{\mu },\stackrel{~}{\varphi })`$. Let $$e_0=\stackrel{~}{\varphi }=\frac{1}{\sqrt{2}}\left(\begin{array}{c}\varphi \\ \varphi \end{array}\right),$$ (101) and $$e_1=\frac{1}{K}[\stackrel{~}{\mu }(\stackrel{~}{\mu }e_0)e_0],$$ (102) where $`K`$ is chosen so that $`e_1^2=1`$. It is easy to show that $$K=\|\stackrel{~}{\mu }(\stackrel{~}{\mu }e_0)e_0\|=\sqrt{\frac{N_S\frac{1}{2}}{N_S}}.$$ (103) Thus, $$e_1=\frac{1}{\sqrt{N_S\frac{1}{2}}}\left(\begin{array}{c}\varphi ^{not}+\frac{\varphi }{2}\\ \frac{\varphi }{2}\end{array}\right).$$ (104) Furthermore, $$\stackrel{~}{\mu }=\frac{1}{\sqrt{2N_S}}[e_0+\sqrt{2N_S1}e_1].$$ (105) Fig.7 portrays various vectors that arise in explaining Younes’ version of Grover’s algorithm. Since we plan to stay within the two dimensional vector space with orthonormal basis $`e_0,e_1`$, it is convenient to switch matrix representations. Within $`span(e_0,e_1)`$, $`e_0,e_1`$ can be represented more simply by: $$e_0=\left(\begin{array}{c}1\\ 0\end{array}\right),e_1=\left(\begin{array}{c}0\\ 1\end{array}\right).$$ (106) If $`e_0,e_1`$ are represented this way, then $$\stackrel{~}{\varphi }=\left(\begin{array}{c}1\\ 0\end{array}\right),\stackrel{~}{\mu }=\frac{1}{\sqrt{2N_S}}\left(\begin{array}{c}1\\ \sqrt{2N_S1}\end{array}\right),$$ (107) and $$R_{\stackrel{~}{\mu }}R_{\stackrel{~}{\varphi }}=\frac{1}{N_S}\left(\begin{array}{cc}N_S1& \sqrt{2N_S1}\\ \sqrt{2N_S1}& N_S1\end{array}\right).$$ (108) Thus, $$R_{\stackrel{~}{\mu }}R_{\stackrel{~}{\varphi }}=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right),$$ (109) where $$\mathrm{sin}\theta =\frac{\sqrt{2N_S1}}{N_S}\sqrt{\frac{2}{N_S}}.$$ (110) A comparison of Eq.(72) (for the original Grover’s algorithm) and Eq.(107) (for Younes’s version of Grover’s algorithm) reveals that for the purpose of finding the optimal number $`r`$ of iterations, Younes’ algorithm is the same as Grover’s algorithm if one replaces $`N_S`$ in Grover’s algorithm by $`2N_S`$. This comes from the fact that Younes’ algorithm uses $`N_B+1`$ bits whereas Grover’s uses $`N_B`$. ## 4 Generalization of Standard Algorithms, <br>a list of Desiderata So far we have analyzed several standard quantum computing algorithms, namely those attributed to Deutsch-Jozsa, Bernstein-Vazirani, Simon and Grover. (Two other standard algorithm’s that we didn’t analyze are Shor’s algorithm and the algorithm for Teleportation.) In this section, we will try to point out those features of the standard algorithms that would be, in our opinion, fruitful to generalize. Bear in mind that generalizations are seldom unique, but some are more natural, fruitful and far-reaching than others. (a) Allow more complicated graph topologies The standard algorithms discussed here can all be represented by QB nets with trivial topologies such as 2 body scattering graphs or Markov chains. However, other important quantum algorithms, such as the one for Teleportation, can be represented by QB nets with more complicated graph topologies (e.g., with loops). (b) Estimate more general probability distributions The goal of most standard algorithms is to estimate a deterministic probability distribution. However, estimating non-deterministic ones is also very useful. Such estimates are useful in, for example, applications of Decision Theory and Artificial Intelligence, where inferences are made based on uncertain knowledge. (c) Allow multiple runs and the rejection of some If one is estimating a non-deterministic probability distribution, it will be necessary to do multiple runs. It may also be necessary to allow rejection of runs. Obviously, the number of rejected runs is best kept as small as possible. (d) Allow more general measurements Suppose $`\underset{¯}{x}`$ is a node of a QB net. Let $`S_{\underset{¯}{x}}`$ be the set of its possible states. We will say that node $`\underset{¯}{x}`$ has been measured if during the experiment which the QB net describes, a measurement is performed that restricts the possible states of $`\underset{¯}{x}`$ to a proper subset $`S_{\underset{¯}{x}}^{}`$ of $`S_{\underset{¯}{x}}`$. When $`\underset{¯}{x}`$ is an internal (ditto, external) node of the QB net, we will refer to its measurement as an internal (ditto, external) measurement. The standard algorithms discussed here use external but no internal measurements. However, other important quantum algorithms, such as the one for Teleportation, do use internal ones. ## 5 Q-Embeddings The remainder of this paper will be devoted to discussing a class of algorithms which generalizes some standard algorithms and achieves some of the desiderata given in the previous section. Our algorithms are based on the idea that, given a CB net, one can always embed it in a QB net. Simple examples of such q-embeddings have already been given in the sections dealing with standard algorithms. We start by defining some terminology that will be useful. A probability matrix $`P(y\|x)`$ is a rectangular (not necessarily square) matrix with row index $`yS_{\underset{¯}{y}}`$ and column index $`xS_x`$ such that $`P(y\|x)0`$ for all $`x,y`$, and $`_yP(y\|x)=1`$ for all $`x`$. The set of all probability matrices $`P(y\|x)`$ where $`xS_{\underset{¯}{x}}`$ and $`yS_{\underset{¯}{y}}`$ will be denoted by $`pd(S_{\underset{¯}{y}}\|S_{\underset{¯}{x}})`$ (pd = probability distribution). A probability matrix is assigned to each node of a CB net. A probability matrix $`P(y\|x)`$ is deterministic if for each column $`x`$, there exists a single row $`y`$, call it $`y(x)`$, such that $`P(y\|x)=\delta (y(x),y)`$. Any map $`f:S_{\underset{¯}{x}}S_{\underset{¯}{y}}`$ uniquely specifies (and is uniquely specified) by the deterministic probability matrix $`P`$ with matrix elements $`P(y\|x)=\delta (y,f(x))`$ for all $`xS_{\underset{¯}{x}}`$ and $`yS_{\underset{¯}{y}}`$. We will often talk about a map $`f`$ and its associated probability matrix $`P(y\|x)`$ as if they were the same thing. Given two matrices $`A`$ and $`B`$ of the same dimensions, their Hadamard product $`C=AB`$ is defined by $`C_{i,j}=A_{i,j}B_{i,j}`$ for all $`i,j`$. We will call $`\mathrm{HAS}(A)=AA^{}`$ the Hadamard Absolute Square (HAS) of matrix $`A`$. If $`U`$ is a unitary matrix, then $`\mathrm{HAS}(U)`$ is a probability matrix. For example, for any angle $`\theta `$, $$\mathrm{HAS}(\left[\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right])=\left[\begin{array}{cc}\mathrm{cos}^2\theta & \mathrm{sin}^2\theta \\ \mathrm{sin}^2\theta & \mathrm{cos}^2\theta \end{array}\right].$$ (111) Another example is $$\mathrm{HAS}(\widehat{H}_1)=\frac{1}{2}\left[\begin{array}{cc}1& 1\\ 1& 1\end{array}\right].$$ (112) A CB net $`𝒩^C`$ is the HAS of QB net $`𝒩^Q`$ if $`𝒩^Q`$ and $`𝒩^C`$ have the same graph, and their node matrices are related as follows. For each node $`\underset{¯}{x}_i`$, if $`A[x_i\|(x.)_{\mathrm{\Gamma }_i}]`$ is the amplitude of node $`\underset{¯}{x}_i`$ in $`𝒩^Q`$, and $`P[x_i\|(x.)_{\mathrm{\Gamma }_i}]`$ is the probability of node $`\underset{¯}{x}_i`$ in $`𝒩^C`$, then $`\|A[x_i\|(x.)_{\mathrm{\Gamma }_i}]\|^2=P[x_i\|(x.)_{\mathrm{\Gamma }_i}]`$. In such a case, we will write $`\mathrm{HAS}(𝒩^Q)=𝒩^C`$. A unitary matrix $`A(y,\stackrel{~}{x}\|x,\stackrel{~}{y})`$ (with rows labelled by $`y,\stackrel{~}{x}`$ and columns by $`x,\stackrel{~}{y}`$) is a q-embedding of probability matrix $`P(y\|x)`$ if $$\underset{\stackrel{~}{x}}{}\|A(y,\stackrel{~}{x}\|x,\stackrel{~}{y}=0)\|^2=P(y\|x)$$ (113) for all possible values of $`y`$ and $`x`$. (the “q” in “q-embedding” stands for “quantum”). We say $`\stackrel{~}{y}`$ is a source index and $`\stackrel{~}{x}`$ is a sink index. We also refer to $`\stackrel{~}{x}`$ and $`\stackrel{~}{y}`$ collectively as ancilla indices. Note that any unitary matrix is a q-embedding of its HAS. Indeed, in this case Eq.(113) is satisfied with the indices $`\stackrel{~}{x}`$ and $`\stackrel{~}{y}`$ each ranging over a single value (i.e., $`\stackrel{~}{x}`$ and $`\stackrel{~}{y}`$ are fixed). If a q-embedding satisfies $`A(y,\stackrel{~}{x}\|x,\stackrel{~}{y})Bool`$ for all $`y,\stackrel{~}{x},x,\stackrel{~}{y}`$, we say that it is a deterministic q-embedding or a deterministic reversible extension (DRE) of its probability matrix (note that its probability matrix must also be deterministic). By an extension of a matrix we mean adding extra rows and/or columns to it. General q-embeddings use the square root of the entries of the original probability matrix so they are not simply extensions of the original matrix; they are, however, reversible since they are unitary matrices. Given a QB net $`𝒩^Q`$, let $$P[(x.)_L]=\|\underset{(x.)_{\mathrm{\Gamma }_QL}}{}A(x.)\|^2.$$ (114) On the right hand side of Eq.(114), $`A(x.)`$ is the amplitude of story $`(x.)`$, $`\mathrm{\Gamma }_Q`$ is the set of indices of all the nodes of $`𝒩^Q`$, and $`L`$ is the set of indices of all leaf (aka external) nodes of $`𝒩^Q`$. We say $`𝒩^Q`$ is a q-embedding of CB net $`𝒩^C`$ if $`P[(x.)_L]`$ defined by Eq.(114) satisfies $$P[(x.)_{\mathrm{\Gamma }_C}]=\underset{L_1}{}P[(x.)_L],$$ (115) where $`L_1L`$, and $`\mathrm{\Gamma }_C`$ is the set of indices of all nodes of $`𝒩^C`$. Thus, the probability distribution associated with all nodes of $`𝒩^C`$ can be obtained from the probability distribution associated with the external nodes of $`𝒩^Q`$. Some examples of q-embeddings of CB nets have already been given during our discussion of standard algorithms. More examples will be given in subsequent sections. For some positive integers $`r`$ and $`s`$, we will say a map $`f:Bool^rBool^s`$ is a binary gate from $`r`$ to $`s`$ bits. $`f`$ uniquely specifies (and is uniquely specified) by the deterministic probability matrix with entries $`P(y\|x)=\delta (f(x),y)`$, where $`x=(x_0,x_1,\mathrm{},x_{r1})Bool^r`$ and $`y=(y_0,y_1,\mathrm{},y_{s1})Bool^s`$. If $`f`$ is an invertible map, we will say that the gate is reversible. For example, the AND gate which takes $`(x_1,x_0)y_0`$ with $`y_0=x_0x_1`$ is a binary gate. So are the OR and NOT gates. Out of these 3 gates, only the NOT gate is reversible. Another example of a reversible binary gate is the Toffoli gate. It maps 3 bits into 3 bits as follows: $$\begin{array}{c}y_0=T_0(x)=x_0,\hfill \\ y_1=T_1(x)=x_1,\hfill \\ y_2=T_2(x)=x_2x_0x_1.\hfill \end{array}$$ (116) The Toffoli gate can also be defined as the following deterministic probability matrix $$P(y\|x)=\delta (y,T(x))=\delta (y_2,x_2x_0x_1)\delta (y_1,x_1)\delta (y_0,x_0).$$ (117) Consider 3 bits labelled 0, 1, and 2, and suppose the $`i`$th bit changes value from $`x_i`$ to $`y_i`$. Then bits 0 and 1 do not change whereas bit 2 flips iff the product $`x_0x_1`$ equals one. Thus, the probability matrix with entries given by Eq.(117) is simply a doubly controlled not: $$[P(y\|x)]=\sigma _x(2)^{n(1)n(0)}.$$ (118) It is convenient to use the term Toffoli gate to refer not only to the gate defined by Eq.(117), but also to the 3 other gates that one obtains by replacing $`x_0x_1`$ in Eq.(117) by $`x_0\overline{x_1}`$, or $`\overline{x_0}x_1`$, or $`\overline{x_0}\overline{x_1}`$. This corresponds to replacing $`n(1)n(0)`$ in Eq.(118) by $`n(1)\overline{n}(0)`$, or $`\overline{n}(1)n(0)`$, or $`\overline{n}(1)\overline{n}(0)`$. Fig.8 shows the 4 doubly-controlled nots that we call Toffoli gates as well as the circuit diagrams usually used to represent them. ### 5.1 Q-Embedding of Probability Matrices In this section we will first give some examples of q-embeddings of probability matrices. Then we will show that any probability matrix has a q-embedding. Any unitary matrix is a q-embedding of its HAS, but such q-embeddings are trivial in the sense that they have no ancilla indices. As first shown in Refs., the Toffoli gates can be used to build q-embeddings (in fact, DREs) of the elementary binary gates AND, XOR, NOT, FANOUT. See Fig.9. Let $`x=(x_0,x_1,x_2)Bool^3`$ and $`y=(y_0,y_1,y_2)Bool^3`$. For the AND gate, $$\underset{y_1,y_0}{}\left\|y\|\sigma _x(2)^{n(1)n(0)}\|x_2=0,x_1,x_0\right\|^2=\delta (y_2,x_1x_0).$$ (119a) For the FANOUT gate, $$\underset{y_1}{}\left\|y\|\sigma _x(2)^{\overline{n}(1)n(0)}\|x_2=0,x_1=0,x_0\right\|^2=\delta (y_2,x_0)\delta (y_0,x_0).$$ (119b) For the XOR gate, $$\underset{y_1,y_0}{}\left\|y\|\sigma _x(2)^{n(1)\overline{n}(0)}\|x_2,x_1,x_0=0\right\|^2=\delta (y_2,x_2x_1).$$ (119c) For the NOT gate, $$\underset{y_1,y_0}{}\left\|y\|\sigma _x(2)^{\overline{n}(1)\overline{n}(0)}\|x_2,x_1=0,x_0=0\right\|^2=\delta (y_2,x_21)=\delta (y_2,\overline{x_2}).$$ (119d) Note that the NOT gate is just $`\sigma _x`$, which is a DRE of itself. Eq.(119d) gives a different DRE of $`\sigma _x`$. In the left hand side of Eqs.(119), the $`x_i`$ indices that are set to zero are called source indices, and the $`y_i`$ indices that are summed over are called sink indices. Sink and source indices are collectively called ancilla indices. Next we will prove that any probability matrix has a q-embedding. Suppose that we are given a probability matrix $`P(y\|x)`$ where $`xS_{\underset{¯}{x}}`$ and $`yS_{\underset{¯}{y}}`$. Let $`N_{\underset{¯}{x}}`$ (ditto, $`N_{\underset{¯}{y}}`$) denote the number of elements in $`S_{\underset{¯}{x}}`$ (ditto, $`S_{\underset{¯}{y}}`$). Let $`\xi ^{(x)}`$ for $`xS_{\underset{¯}{x}}`$ be any orthonormal basis of the complex $`N_{\underset{¯}{x}}`$ dimensional vector space. The components of $`\xi ^{(x)}`$ will be denoted by $`\xi _{\stackrel{~}{x}}^{(x)}`$, where $`\stackrel{~}{x}S_{\underset{¯}{x}}`$. If the $`\xi ^{(x)}`$’s are the standard basis, then $`\xi _{\stackrel{~}{x}}^{(x)}=\delta (x,\stackrel{~}{x})`$. Define matrix $`A`$ by $$A(y,\stackrel{~}{x}\|x,\stackrel{~}{y})=\{\begin{array}{cc}\sqrt{P(y\|x)}\xi _{\stackrel{~}{x}}^x\hfill & \text{ if}\stackrel{~}{y}=0\hfill \\ \text{obtained by Gram-Schmidt method}\hfill & \text{ if}\stackrel{~}{y}0\hfill \end{array}.$$ (120) To understand the last equation, consider Fig.10. In that figure we have assumed for definiteness that $`S_{\underset{¯}{x}}=\{0,1,2\}`$ and $`S_{\underset{¯}{y}}=\{0,1,2,3\}`$. The shaded (ditto, unshaded) columns have $`\stackrel{~}{y}0`$ (ditto, $`\stackrel{~}{y}=0`$). It is easy to see that the unshaded columns are orthonormal because the vectors $`\xi ^x`$ are orthonormal and $`_yP(y\|x)=1`$. Since the unshaded columns are orthonormal, one can use the Gram-Schmidt method to fill the shaded columns so that all the columns of $`A`$ are orthonormal and therefore $`A`$ is unitary. Note that by virtue of Eq.(120), $$\underset{\stackrel{~}{x}}{}\|A(y,\stackrel{~}{x}\|x,\stackrel{~}{y}=0)\|^2=\underset{\stackrel{~}{x}}{}P(y\|x)\xi _{\stackrel{~}{x}}^{(x)}\xi _{\stackrel{~}{x}}^{(x)}=P(y\|x)$$ (121) so that the $`A`$ defined by Eq.(120) does indeed satisfy Eq.(113). Note that the matrix $`A`$ defined by Eq.(120) has dimensions $`N_{\underset{¯}{x}}N_{\underset{¯}{y}}\times N_{\underset{¯}{x}}N_{\underset{¯}{y}}`$. It is sometimes possible to find a smaller q-embedding of an $`N_{\underset{¯}{y}}\times N_{\underset{¯}{x}}`$ probability matrix $`P(y\|x)`$. For example, $`\sigma _x`$ is a q-embedding of itself. As a less trivial example, suppose $$P(y\|x_1,x_2)=\delta (y,x_1x_2),$$ (122) for $`y,x_1,x_2Bool`$. Then define $$A(y,e\|x_1,x_2)=\frac{(1)^{x_1e}}{\sqrt{2}}\delta (y,x_1x_2),$$ (123) for $`y,e,x_1,x_2Bool`$. It is easy to check that matrix $`A`$ is unitary. Furthermore, $$\underset{e}{}\|A(y,e\|x_1,x_2)\|^2=\frac{1}{2}\underset{e}{}\delta (y,x_1x_2)=\delta (y,x_1x_2).$$ (124) ### 5.2 Q-Embedding of CB Nets As we’ve said before, F-T showed in Refs. how, given any binary gate $`f`$, one can construct another binary gate $`\overline{f}`$ such that $`\overline{f}`$ is a DRE of $`f`$. Their method for constructing $`\overline{f}`$ is to first represent $`f`$ as a binary deterministic circuit composed of elementary gates (AND, XOR, NOT, FANOUT), and then to modify the circuit by replacing each of its gates by a DRE of it. The desired gate $`\overline{f}`$ is then specified by the modified circuit. In this section we will show how, given any CB net $`𝒩^C`$, one can construct a QB net $`𝒩^Q`$ which is a q-embedding of $`𝒩^C`$. So far we’ve shown how to construct a q-embedding for any probability matrix. Now remember that each node of $`𝒩^C`$ has a probability matrix assigned to it. The main step in constructing a q-embedding of $`𝒩^C`$ is to replace each node matrix of $`𝒩^C`$ with a q-embedding of it. Thus, our method for constructing a q-embedding of a CB net is a generalization of the F-T method for constructing a DRE of a binary deterministic circuit. We generalize their method so that it can be applied to any classical stochastic circuit, not just binary deterministic ones. Before describing our construction method, we need some definitions. We say a node $`\underset{¯}{m}`$ is a marginalizer node if it has a single input arrow and a single output arrow. Furthermore, the parent node of $`\underset{¯}{m}`$, call it $`\underset{¯}{x}`$, has states $`x=(x_1,x_2,\mathrm{},x_n)`$, where $`x_iS_{\underset{¯}{x}_i}`$ for each $`iZ_{1,n}`$. Furthermore, for some particular integer $`i_0Z_{1,n}`$, the set of possible states of $`\underset{¯}{m}`$ is $`S_{\underset{¯}{m}}=S_{\underset{¯}{x}_{i_0}}`$, and the node matrix of $`\underset{¯}{m}`$ is $`P(\underset{¯}{m}=m\|\underset{¯}{x}=x)=\delta (m,x_{i_0})`$. Let $`𝒩^C`$ be a CB net for which we want to obtain a q-embedding. Our construction has two steps: (Step 1) Add marginalizer nodes. More specifically, replace $`𝒩^C`$ by a modified CB net $`𝒩_{}^{C}{}_{mod}{}^{}`$ obtained as follows. For each node $`\underset{¯}{x}`$ of $`𝒩^C`$, add a marginalizer node between $`\underset{¯}{x}`$ and every child of $`\underset{¯}{x}`$. If $`\underset{¯}{x}`$ has no children, add a child to it. As an example of this step, consider the net $`𝒩^C`$ (“two body scattering net”) defined by Fig.11 and Table 5.2. \| nodes \| states \| probabilities \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{a}`$ \| $`aS_{\underset{¯}{a}}`$ \| $`P(a)`$ \| \| \| $`\underset{¯}{b}`$ \| $`bS_{\underset{¯}{b}}`$ \| $`P(b)`$ \| \| \| $`\underset{¯}{c}`$ \| $`cS_{\underset{¯}{c}}`$ \| $`P(c\|x)`$ \| \| \| $`\underset{¯}{d}`$ \| $`dS_{\underset{¯}{d}}`$ \| $`P(d\|x)`$ \| \| \| $`\underset{¯}{x}`$ \| $`xS_{\underset{¯}{x}}`$ \| $`P(x\|a,b)`$ \| \| Table 7 Applying Step 1 to $`𝒩^C`$ for two body scattering yields $`𝒩_{}^{C}{}_{mod}{}^{}`$ defined by Fig.12 and Table 5.2. \| nodes \| states \| probabilities \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{a}_2`$ \| $`a_2S_{\underset{¯}{a}}`$ \| $`P_{\underset{¯}{a}}(a_2)`$ \| \| \| $`\underset{¯}{a}_3`$ \| $`a_3S_{\underset{¯}{a}}`$ \| $`\delta (a_3,a_2)`$ \| \| \| $`\underset{¯}{b}_2`$ \| $`b_2S_{\underset{¯}{b}}`$ \| $`P_{\underset{¯}{b}}(b_2)`$ \| \| \| $`\underset{¯}{b}_3`$ \| $`b_3S_{\underset{¯}{b}}`$ \| $`\delta (b_3,b_2)`$ \| \| \| $`\underset{¯}{c}_2`$ \| $`c_2S_{\underset{¯}{c}}`$ \| $`P_{\underset{¯}{c}\|\underset{¯}{x}}(c_2\|x_{3c})`$ \| \| \| $`\underset{¯}{c}_3`$ \| $`c_3S_{\underset{¯}{c}}`$ \| $`\delta (c_3,c_2)`$ \| \| \| $`\underset{¯}{d}_2`$ \| $`d_2S_{\underset{¯}{d}}`$ \| $`P_{\underset{¯}{d}\|\underset{¯}{x}}(d_2\|x_{3d})`$ \| \| \| $`\underset{¯}{d}_3`$ \| $`d_3S_{\underset{¯}{d}}`$ \| $`\delta (d_3,d_2)`$ \| \| \| $`(\underset{¯}{x}_{2c},\underset{¯}{x}_{2d})`$ \| $`(x_{2c},x_{2d})S_{\underset{¯}{x}}^2`$ \| $`P_{\underset{¯}{x}\|\underset{¯}{a},\underset{¯}{b}}(x_{2c}\|a_3,b_3)\delta (x_{2d},x_{2c})`$ \| \| \| $`\underset{¯}{x}_{3c}`$ \| $`x_{3c}S_{\underset{¯}{x}}`$ \| $`\delta (x_{3c},x_{2c})`$ \| \| \| $`\underset{¯}{x}_{3d}`$ \| $`x_{3d}S_{\underset{¯}{x}}`$ \| $`\delta (x_{3d},x_{2d})`$ \| \| Table 8 (Step 2) Replace node probability matrices by their q-embeddings. Add ancilla nodes. More specifically, replace $`𝒩_{}^{C}{}_{mod}{}^{}`$ by a QB net $`𝒩^Q`$ obtained as follows. For each node of $`𝒩_{}^{C}{}_{mod}{}^{}`$, except for the marginalizer nodes that were added in the previous step, replace its node matrix by a new node matrix which is a q-embedding of the original node matrix. Add a new node for each ancilla index of each new node matrix. These new nodes will be called ancilla nodes (of either the source or sink type) because they correspond to ancilla indices. Applying Step 2 to net $`𝒩_{}^{C}{}_{mod}{}^{}`$ for two body scattering yields $`𝒩^Q`$ defined by Fig.13 and Table 5.2. \| nodes \| states \| amplitudes \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{a}_1`$ \| $`a_1S_{\underset{¯}{a}}`$ \| $`\delta (a_1,0)`$ \| \| \| $`\underset{¯}{a}_2`$ \| $`a_2S_{\underset{¯}{a}}`$ \| $`A\left(a_2\|a_1=0\right)=\sqrt{P_{\underset{¯}{a}}\left(a_2\right)}`$ \| \| \| $`\underset{¯}{a}_3`$ \| $`a_3S_{\underset{¯}{a}}`$ \| $`\delta (a_3,a_2)`$ \| \| \| $`(\underset{¯}{a}_4,\underset{¯}{b}_4,\underset{¯}{x}_{2c},\underset{¯}{x}_{2d})`$ \| $`(a_4,b_4,x_{2c},x_{2d})S_{\underset{¯}{a},\underset{¯}{b},\underset{¯}{x},\underset{¯}{x}}`$ \| $`A(a_4,b_4,x_{2c},x_{2d}\|a_3,b_3,x_{1c}=0,x_{1d}=0)=`$ \| \| \| \| \| $`\sqrt{P_{\underset{¯}{x}\|\underset{¯}{a},\underset{¯}{b}}\left(x_{2c}\|a_3,b_3\right)}\delta (a_4,a_3)\delta (b_4,b_3)\delta (x_{2d},x_{2c})`$ \| \| \| $`\underset{¯}{a}_5`$ \| $`a_5S_{\underset{¯}{a}}`$ \| $`\delta (a_5,a_4)`$ \| \| \| $`\underset{¯}{b}_1`$ \| $`b_1S_{\underset{¯}{b}}`$ \| $`\delta (b_1,0)`$ \| \| \| $`\underset{¯}{b}_2`$ \| $`b_2S_{\underset{¯}{b}}`$ \| $`A\left(b_2\|b_1=0\right)=\sqrt{P_{\underset{¯}{b}}\left(b_2\right)}`$ \| \| \| $`\underset{¯}{b}_3`$ \| $`b_3S_{\underset{¯}{b}}`$ \| $`\delta (b_3,b_2)`$ \| \| \| $`\underset{¯}{b}_5`$ \| $`b_5S_{\underset{¯}{b}}`$ \| $`\delta (b_5,b_4)`$ \| \| \| $`\underset{¯}{c}_1`$ \| $`c_1S_{\underset{¯}{c}}`$ \| $`\delta (c_1,0)`$ \| \| \| $`(\underset{¯}{c}_2,\underset{¯}{x}_{4c})`$ \| $`(c_2,x_{4c})S_{\underset{¯}{c},\underset{¯}{x}}`$ \| $`A(c_2,x_{4c}\|c_1=0,x_{3c})=`$ \| \| \| \| \| $`\sqrt{P_{\underset{¯}{c}\|\underset{¯}{x}}\left(c_2\|x_{3c}\right)}\delta (x_{4c},x_{3c})`$ \| \| \| $`\underset{¯}{c}_3`$ \| $`c_3S_{\underset{¯}{c}}`$ \| $`\delta (c_3,c_2)`$ \| \| \| $`\underset{¯}{d}_1`$ \| $`d_1S_{\underset{¯}{d}}`$ \| $`\delta (d_1,0)`$ \| \| \| $`(\underset{¯}{d}_2,\underset{¯}{x}_{4d})`$ \| $`(d_2,x_{4d})S_{\underset{¯}{d},\underset{¯}{x}}`$ \| $`A(d_2,x_{4d}\|d_1=0,x_{3d})=`$ \| \| \| \| \| $`\sqrt{P_{\underset{¯}{d}\|\underset{¯}{x}}\left(d_2\|x_{3d}\right)}\delta (x_{4d},x_{3d})`$ \| \| \| $`\underset{¯}{d}_3`$ \| $`d_3S_{\underset{¯}{d}}`$ \| $`\delta (d_3,d_2)`$ \| \| \| $`\underset{¯}{x}_{1c}`$ \| $`x_{1c}S_{\underset{¯}{x}}`$ \| $`\delta (x_{1c},0)`$ \| \| \| $`\underset{¯}{x}_{1d}`$ \| $`x_{1d}S_{\underset{¯}{x}}`$ \| $`\delta (x_{1d},0)`$ \| \| \| $`\underset{¯}{x}_{3c}`$ \| $`x_{3c}S_{\underset{¯}{x}}`$ \| $`\delta (x_{3c},x_{2c})`$ \| \| \| $`\underset{¯}{x}_{3d}`$ \| $`x_{3d}S_{\underset{¯}{x}}`$ \| $`\delta (x_{3d},x_{2d})`$ \| \| \| $`\underset{¯}{x}_{5c}`$ \| $`x_{5c}S_{\underset{¯}{x}}`$ \| $`\delta (x_{5c},x_{4c})`$ \| \| \| $`\underset{¯}{x}_{5d}`$ \| $`x_{5d}S_{\underset{¯}{x}}`$ \| $`\delta (x_{5d},x_{4d})`$ \| \| Table 9 $`𝒩^Q`$ looks much more complicated than $`𝒩^C`$, but it really isn’t, since most of its node matrices are delta functions which quickly disappear when adding over node states. According to Table 5.2, the probability amplitude for the external (aka leaf) nodes is given by $$\begin{array}{c}A(a_5,b_5,c_3,d_3,x_{5c},x_{5d})=\hfill \\ =_{}\sqrt{P_{\underset{¯}{a}}(\text{}_2)P_{\underset{¯}{b}}(\text{}_2)P_{\underset{¯}{x}\|\underset{¯}{a},\underset{¯}{b}}(\text{}_{2c}\|\text{}_3,\text{}_3)P_{\underset{¯}{c}\|\underset{¯}{x}}(\text{}_2\|\text{}_{3c})P_{\underset{¯}{d}\|\underset{¯}{x}}(\text{}_2\|\text{}_{3d})}\hfill \\ \theta (\text{}_2=\text{}_3=\text{}_4=a_5)\theta (\text{}_2=\text{}_3=\text{}_4=b_5)\hfill \\ \theta (\text{}_{2c}=\text{}_{3c}=\text{}_{4c}=x_{5c})\theta (\text{}_{2d}=\text{}_{3d}=\text{}_{4d}=x_{5d})\theta (x_{5c}=x_{5d})\hfill \\ \theta (\text{}_2=c_3)\theta (\text{}_2=d_3)\hfill \\ \theta (\text{}_1=\text{}_1=\text{}_1=\text{}_1=\text{}_{1c}=\text{}_{1d}=0)\hfill \end{array},$$ (125) where we have summed over all internal (non-leaf) nodes. Eq.(125) immediately reduces to $$\begin{array}{c}A(a_5,b_5,c_3,d_3,x_{5c},x_{5d})=\hfill \\ =\sqrt{P_{\underset{¯}{a}}(a_5)P_{\underset{¯}{b}}(b_5)P_{\underset{¯}{x}\|\underset{¯}{a},\underset{¯}{b}}(x_{5c}\|a_5,b_5)P_{\underset{¯}{c}\|\underset{¯}{x}}(c_3\|x_{5c})P_{\underset{¯}{d}\|\underset{¯}{x}}(d_3\|x_{5d})}\theta (x_{5c}=x_{5d})\hfill \end{array}.$$ (126) Eq.(126) shows that the net $`𝒩^Q`$ that we constructed from the net $`𝒩^C`$ by following steps 1 and 2 satisfies the definition Eq.(115) that we gave earlier for a q-embedding of $`𝒩^C`$. The probability distribution of the states of the external nodes of the QB net $`𝒩^Q`$ contains all the probabilistic information of the original CB net $`𝒩^C`$. Hurray! From Eq.(126), it is clear that by running $`𝒩^Q`$ on a quantum computer (or similar quantum system), we can calculate any conditional probability that one would want to calculate for $`𝒩^C`$. For example, suppose we wanted to calculate $`P_{\underset{¯}{a},\underset{¯}{d}\|\underset{¯}{x}}`$. Run $`𝒩^Q`$ on the quantum computer several times, each time measuring nodes $`\underset{¯}{a}_5,\underset{¯}{d}_3`$ and $`\underset{¯}{x}_{5d}`$ and not measuring all other external nodes. The resulting measurements will be distributed according to $`P_{\underset{¯}{a},\underset{¯}{d},\underset{¯}{x}}`$. Taking the magnitude squared of the amplitude and summing the result over the states of the un-measured external nodes will be performed automatically by Nature. The laws of quantum mechanics guarantee it. Proceed in the same way to calculate $`P_{\underset{¯}{x}}`$. Run $`𝒩^Q`$ on the quantum computer several times, each time measuring node $`\underset{¯}{x}_{5d}`$ and not measuring all other external nodes. Finally divide $`P_{\underset{¯}{a},\underset{¯}{d},\underset{¯}{x}}`$ by $`P_{\underset{¯}{x}}`$ on a classical (or quantum?) computer. The q-embedding of a CB net, as defined by Eq.(115), is not unique. For example, we could have defined the net $`𝒩^Q`$ given by Fig.13 without nodes $`\underset{¯}{a}_3`$ and $`\underset{¯}{b}_3`$. We chose to include such nodes for pedagogical reasons. To run a QB net on a quantum computer, we need to replace the QB net by an equivalent SEO that a quantum computer can understand. This can be done with the help of a quantum compiler . One could compile individually each node representing a q-embedding, or one could compile whole subgraphs of the QB net all at once. Note that it may suffice to find a SEO that is only approximately (within a certain precision) equivalent instead of exactly equivalent to the QB net. This may be true if, for example, the probabilities associated with the CB net that was q-embedded were not specified too precisely to begin with. Suppose $`a_1,a_2,\mathrm{}a_\nu `$ belong to a finite set $`S_{\underset{¯}{a}}`$, and suppose that they are distributed according to a probability distribution $`P_{\underset{¯}{a}}`$. What number $`\nu `$ of samples $`a_i`$ is necessary to estimate $`P_{\underset{¯}{a}}`$ within a given precision? This question is directly relevant to our method for estimating probabilities by running a QB net on a quantum computer. We will not give a detailed answer to this question here. For an answer, the reader can consult any book on the mathematical theory of Statistics. An imprecise rule of thumb is that if the support of $`P_{\underset{¯}{a}}`$ has $`\nu _0`$ elements, then $`\nu `$ must be at least as large as $`\nu _0`$; i.e., one needs at least “one data point per bin” to estimate $`P_{\underset{¯}{a}}`$ with any decent accuracy. We have given a method for calculating, via a quantum computer, the conditional probabilities associated with a CB net. Does our method have an advantage in time complexity with respect to classical methods for calculating the same probabilities? We will not give a detailed answer to this question here. The answer must be yes, sometimes. After all, our method generalizes the algorithms by Deutsch-Jozsa, Simon, Grover, etc., and these are known to have a complexity advantage. To conclude this section, we will present a second, more complicated example of our method of finding a q-embedding for a CB net. A CB net (first given in Ref.) for lung disease diagnosis is defined by Fig.14 and Table 5.2. \| nodes \| states \| probabilities \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{a}`$ \| $`aBool`$ \| $`P\left(\underset{¯}{a}=1\right)=.01`$ \| Visited Asia? \| \| $`\underset{¯}{b}`$ \| $`bBool`$ \| $`\begin{array}{c}P\left(\underset{¯}{b}=1\|\underset{¯}{s}=1\right)=.60\hfill \\ P\left(\underset{¯}{b}=1\|\underset{¯}{s}=0\right)=.30\hfill \end{array}`$ \| Bronchitis? \| \| $`\underset{¯}{d}`$ \| $`dBool`$ \| $`\begin{array}{c}P(\underset{¯}{d}=1\|\underset{¯}{e}=1,\underset{¯}{b}=1)=.90\hfill \\ P(\underset{¯}{d}=1\|\underset{¯}{e}=1,\underset{¯}{b}=0)=.70\hfill \\ P(\underset{¯}{d}=1\|\underset{¯}{e}=0,\underset{¯}{b}=1)=.80\hfill \\ P(\underset{¯}{d}=1\|\underset{¯}{e}=0,\underset{¯}{b}=0)=.10\hfill \end{array}`$ \| Dyspnea(trouble breathing)? \| \| $`\underset{¯}{e}`$ \| $`eBool`$ \| $`P\left(e\|l,t\right)=\delta (e,lt)`$ \| Either TB or Lung Cancer? \| \| $`\underset{¯}{l}`$ \| $`lBool`$ \| $`\begin{array}{c}P\left(\underset{¯}{l}=1\|\underset{¯}{s}=1\right)=.10\hfill \\ P\left(\underset{¯}{l}=1\|\underset{¯}{s}=0\right)=.01\hfill \end{array}`$ \| Lung Cancer? \| \| $`\underset{¯}{s}`$ \| $`sBool`$ \| $`P\left(\underset{¯}{s}=1\right)=.5`$ \| Smokes? \| \| $`\underset{¯}{t}`$ \| $`tBool`$ \| $`\begin{array}{c}P\left(\underset{¯}{t}=1\|\underset{¯}{a}=1\right)=.05\hfill \\ P\left(\underset{¯}{t}=1\|\underset{¯}{a}=0\right)=.01\hfill \end{array}`$ \| Tuberculosis? \| \| $`\underset{¯}{x}`$ \| $`xBool`$ \| $`\begin{array}{c}P\left(\underset{¯}{x}=1\|\underset{¯}{e}=1\right)=.98\hfill \\ P\left(\underset{¯}{x}=1\|\underset{¯}{e}=0\right)=.05\hfill \end{array}`$ \| Positive X-ray? \| Table 10 If one follows the two steps that were described earlier in this section, one obtains the QB net defined by Fig.15 and Table 5.2. \| nodes \| states \| amplitudes \| \| --- \| --- \| --- \| \| $`\underset{¯}{a}_1`$ \| $`a_1Bool`$ \| $`\delta (a_1,0)`$ \| \| $`\underset{¯}{a}_2`$ \| $`a_2Bool`$ \| $`A\left(a_2\|a_1=0\right)=\sqrt{P_{\underset{¯}{a}}\left(a_2\right)}`$ \| \| $`\underset{¯}{a}_3`$ \| $`a_3Bool`$ \| $`\delta (a_3,a_2)`$ \| \| $`(\underset{¯}{a}_4,\underset{¯}{t}_2)`$ \| $`(a_4,t_2)Bool^2`$ \| $`A\left(a_4,t_2\|a_3,t_1=0\right)=\sqrt{P_{\underset{¯}{t}\|\underset{¯}{a}}\left(t_2\|a_3\right)}\delta (a_4,a_3)`$ \| \| $`\underset{¯}{a}_5`$ \| $`a_5Bool`$ \| $`\delta (a_5,a_4)`$ \| \| $`\underset{¯}{b}_1`$ \| $`b_1Bool`$ \| $`\delta (b_1,0)`$ \| \| $`(\underset{¯}{b}_2,\underset{¯}{s}_{4b})`$ \| $`(b_2,s_{4b})Bool^2`$ \| $`A(b_2,s_{4b}\|b_1=0,s_{3b})=\sqrt{P_{\underset{¯}{b}\|\underset{¯}{s}}\left(b_2\|s_{3b}\right)}\delta (s_{4b},s_{3b})`$ \| \| $`\underset{¯}{b}_3`$ \| $`b_3Bool`$ \| $`\delta (b_3,b_2)`$ \| \| $`(\underset{¯}{b}_4,\underset{¯}{d}_2,\underset{¯}{e}_{4d})`$ \| $`(b_4,d_2,e_{4d})Bool^3`$ \| $`A(b_4,d_2,e_{4d}\|b_3,d_1=0,e_{3d})=\sqrt{P_{\underset{¯}{d}\|\underset{¯}{b},\underset{¯}{e}}\left(d_2\|b_3,e_{3d}\right)}\delta (b_4,b_3)\delta (e_{4d},e_{3d})`$ \| \| $`\underset{¯}{b}_5`$ \| $`b_5Bool`$ \| $`\delta (b_5,b_4)`$ \| \| $`\underset{¯}{d}_1`$ \| $`d_1Bool`$ \| $`\delta (d_1,0)`$ \| \| $`\underset{¯}{d}_3`$ \| $`d_3Bool`$ \| $`\delta (d_3,d_2)`$ \| \| $`\underset{¯}{e}_{1d}`$ \| $`e_{1d}Bool`$ \| $`\delta (e_{1d},0)`$ \| \| $`\underset{¯}{e}_{1x}`$ \| $`e_{1x}Bool`$ \| $`\delta (e_{1x},0)`$ \| \| $`(\underset{¯}{e}_{2d},\underset{¯}{e}_{2x},\underset{¯}{l}_4,\underset{¯}{t}_4)`$ \| $`(e_{2d},e_{2x},l_4,t_4)`$ \| $`A(e_{2d},e_{2x},l_4,t_4\|e_{1d}=0,e_{1x}=0,l_3,t_3)=`$ \| \| \| $`Bool^4`$ \| $`\sqrt{P_{\underset{¯}{e}\|\underset{¯}{l},\underset{¯}{t}}\left(e_{2d}\|l_3,t_3\right)}\delta (e_{2x},e_{2d})\delta (l_4,l_3)\delta (t_4,t_3)`$ \| \| $`\underset{¯}{e}_{3d}`$ \| $`e_{3d}Bool`$ \| $`\delta (e_{3d},e_{2d})`$ \| \| $`\underset{¯}{e}_{3x}`$ \| $`e_{3x}Bool`$ \| $`\delta (e_{3x},e_{2x})`$ \| \| $`(\underset{¯}{e}_{4x},\underset{¯}{x}_2)`$ \| $`(e_{4x},x_2)Bool^2`$ \| $`A\left(e_{4x},x_2\|e_{3x},x_1=0\right)=\sqrt{P_{\underset{¯}{x}\|\underset{¯}{e}}\left(x_2\|e_{3x}\right)}\delta (e_{4x},e_{3x})`$ \| \| $`\underset{¯}{e}_{5d}`$ \| $`e_{5d}Bool`$ \| $`\delta (e_{5d},e_{4d})`$ \| \| $`\underset{¯}{e}_{5x}`$ \| $`e_{5x}Bool`$ \| $`\delta (e_{5x},e_{4x})`$ \| \| $`\underset{¯}{l}_1`$ \| $`l_1Bool`$ \| $`\delta (l_1,0)`$ \| \| $`(\underset{¯}{l}_2,\underset{¯}{s}_{4l})`$ \| $`(l_2,s_{4l})Bool^2`$ \| $`A(l_2,s_{4l}\|l_1=0,s_{3l})=\sqrt{P_{\underset{¯}{l}\|\underset{¯}{s}}\left(l_2\|s_{3l}\right)}\delta (s_{4l},s_{3l})`$ \| \| $`\underset{¯}{l}_3`$ \| $`l_3Bool`$ \| $`\delta (l_3,l_2)`$ \| \| $`\underset{¯}{l}_5`$ \| $`l_5Bool`$ \| $`\delta (l_5,l_4)`$ \| \| $`\underset{¯}{s}_{1b}`$ \| $`s_{1b}Bool`$ \| $`\delta (s_{1b},0)`$ \| \| $`\underset{¯}{s}_{1l}`$ \| $`s_{1l}Bool`$ \| $`\delta (s_{1l},0)`$ \| \| $`(\underset{¯}{s}_{2b},\underset{¯}{s}_{2l})`$ \| $`(s_{2b},s_{2l})Bool^2`$ \| $`A(s_{2b},s_{2l}\|s_{1b}=0,s_{1l}=0)=\sqrt{P_{\underset{¯}{s}}\left(s_{2b}\right)}\delta (s_{2l},s_{2b})`$ \| \| $`\underset{¯}{s}_{3b}`$ \| $`s_{3b}Bool`$ \| $`\delta (s_{3b},e_{2b})`$ \| \| $`\underset{¯}{s}_{3l}`$ \| $`s_{3l}Bool`$ \| $`\delta (s_{3l},s_{2l})`$ \| \| $`\underset{¯}{s}_{5b}`$ \| $`s_{5b}Bool`$ \| $`\delta (s_{5b},s_{4b})`$ \| \| $`\underset{¯}{s}_{5l}`$ \| $`s_{5l}Bool`$ \| $`\delta (s_{5l},s_{4l})`$ \| \| $`\underset{¯}{t}_1`$ \| $`t_1Bool`$ \| $`\delta (t_1,0)`$ \| \| $`\underset{¯}{t}_3`$ \| $`t_3Bool`$ \| $`\delta (t_3,t_2)`$ \| \| $`\underset{¯}{t}_5`$ \| $`t_5Bool`$ \| $`\delta (t_5,t_4)`$ \| \| $`\underset{¯}{x}_1`$ \| $`x_1Bool`$ \| $`\delta (x_1,0)`$ \| \| $`\underset{¯}{x}_3`$ \| $`x_3Bool`$ \| $`\delta (x_3,x_2)`$ \| Table 11 According to Table 5.2, the probability amplitude for the external (aka leaf) nodes is given by $$\begin{array}{c}A(a_5,b_5,d_3,e_{5d},e_{5x},l_5,s_{5b},s_{5l},t_5,x_3)=\hfill \\ =\sqrt{P_{\underset{¯}{a}}(a_5)P_{\underset{¯}{t}\|\underset{¯}{a}}(t_5\|a_5)P_{\underset{¯}{b}\|\underset{¯}{s}}(b_5\|s_{5b})P_{\underset{¯}{d}\|\underset{¯}{b},\underset{¯}{e}}(d_3\|b_5,e_{5d})P_{\underset{¯}{e}\|\underset{¯}{l},\underset{¯}{t}}(e_{5d}\|l_5,t_5)P_{\underset{¯}{x}\|\underset{¯}{e}}(x_3\|e_{5d})P_{\underset{¯}{l}\|\underset{¯}{s}}(l_5\|s_{5l})P_{\underset{¯}{s}}(s_{5b})}\hfill \\ \theta (e_{5d},e_{5x})\theta (s_{5b},s_{5l})\hfill \end{array}.$$ (127) ## 6 Voting Net and Grover’s Microscope In this section we will first present a CB net, call it $`𝒩^C`$, that describes voting. Then we will find a QB net $`𝒩^Q`$ that is a q-embedding of $`𝒩^C`$. In certain cases, the probabilities that we wish to find are too small to be measurable by running $`𝒩^Q`$ on a quantum computer. However, we will show that sometimes it is possible to define a new QB net, call it $`𝒩_{}^{Q}{}_{}{}^{}`$, that magnifies and makes measurable the probabilities that were unmeasurable using $`𝒩^Q`$ alone. We will refer to $`𝒩_{}^{Q}{}_{}{}^{}`$ as Grover’s Microscope for $`𝒩^Q`$, because $`𝒩_{}^{Q}{}_{}{}^{}`$ is closely related to Grover’s algorithm, and it magnifies the probabilities found with $`𝒩^Q`$. Suppose $`yBool`$ and $`\stackrel{}{x}=(x^0,x^1,\mathrm{},x^{N_B1})Bool^{N_B}`$. Consider the CB net (“voting net”) defined by Fig.16 and Table 6. \| nodes \| states \| probabilities \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{x}^i\text{for all}iZ_{0,N_B1}`$ \| $`x^iBool`$ \| $`P(x^i)`$ \| \| \| $`\underset{¯}{y}`$ \| $`yBool`$ \| $`P(y\|\stackrel{}{x})`$ \| \| Table 12 Henceforth, we will abbreviate $`P(y=0\|\stackrel{}{x})=p_i`$ and $`P(y=1\|\stackrel{}{x})=q_i`$, where $`i=dec(\stackrel{}{x})Z_{0,N_S1}`$. Hence $`p_i+q_i=1`$ for all $`iZ_{0,N_S1}`$. In general, the probability matrix $`P(y\|\stackrel{}{x})`$ has $`2^{N_B}`$ free parameters (namely, $`p_i`$ for all $`iZ_{0,N_S1}`$). This number of parameters is a forbiddingly large for large $`N_B`$. To ease the task of specifying $`P(y\|\stackrel{}{x})`$ , it is common to impose additional constraints on $`P(y\|\stackrel{}{x})`$. An interesting special type of $`P(y\|\stackrel{}{x})`$ is deterministic $`pd(Bool\|Bool^{N_B})`$ matrices; that is, those that can be expressed in the form $$P(y\|\stackrel{}{x})=\delta (y,f(\stackrel{}{x})),$$ (128) where $`f:Bool^{N_B}Bool`$. In this case, the voting net can be used to pose the satisfiability problem (SAT): given $`y=0`$, find the most likely $`\stackrel{}{x}Bool^{N_B}`$; in other words, find those $`\stackrel{}{x}`$ for which $`f(\stackrel{}{x})=0`$. We say $`f`$ is AND-like if all $`p_i`$ equal zero except for one $`p_i`$ which equals one. For example, for $`N_B=2`$, if $`f`$ is an AND gate, then $$P(y\|\stackrel{}{x})_{AND}=\{\begin{array}{c}\hfill (x^0,x^1)\\ \hfill \begin{array}{cccccc}& & 00\hfill & 01\hfill & 10\hfill & 11\hfill \\ & & & & & \\ y\hfill & 0\hfill & 1\hfill & 1\hfill & 1\hfill & 0\hfill \\ & 1\hfill & 0\hfill & 0\hfill & 0\hfill & 1\hfill \end{array}\end{array}.$$ (129) A slightly more general type of $`P(y\|\stackrel{}{x})`$ is quasi-deterministic $`pd(Bool\|Bool^{N_B})`$ matrices; that is, those that can be expressed in the form $$P(y\|\stackrel{}{x})=\underset{\stackrel{}{t}}{}\delta (y,f(\stackrel{}{t}))P(t^0\|x^0)P(t^1\|x^1)\mathrm{}P(t^{N_B1}\|x^{N_B1}),$$ (130) where $`f:Bool^{N_B}Bool`$ and we sum over all $`\stackrel{}{t}=(t^0,t^1,\mathrm{},t^{N_B1})Bool^{N_B}`$. When $`f(\stackrel{}{t})=t^0t^1\mathrm{}t^{N_B1}`$, $`P(y\|\stackrel{}{x})`$ is called a noisy-OR. Appendix A discusses how to q-embed deterministic $`pd(Bool\|Bool^{N_B})`$ matrices, and how to express such q-embeddings as a SEO . Appendix B discusses the same thing for quasi-deterministic $`pd(Bool\|Bool^{N_B})`$ matrices. A q-embedding for the CB net defined by Fig.16 and Table 6 is given by the QB net defined by Fig.17 and Table 6. \| nodes \| states \| amplitudes \| comments \| \| --- \| --- \| --- \| --- \| \| $`\underset{¯}{\overset{}{x}}_1`$ \| $`\stackrel{}{x}_1Bool^{N_B}`$ \| $`\delta (\stackrel{}{x}_1,0)`$ \| \| \| $`\underset{¯}{\overset{}{x}}_2`$ \| $`\stackrel{}{x}_2Bool^{N_B}`$ \| $`A\left(\stackrel{}{x}_2\|\stackrel{}{x}_1=0\right)=\sqrt{P_{\underset{¯}{\overset{}{x}}}\left(\stackrel{}{x}_2\right)}`$ \| \| \| $`(\underset{¯}{\overset{}{x}}_3,\underset{¯}{y}_2)`$ \| $`(\stackrel{}{x}_3,y_2)Bool^{N_B+1}`$ \| $`A\left(\stackrel{}{x}_3,y_2\|\stackrel{}{x}_2,y_1=0\right)=\sqrt{P_{\underset{¯}{y}\|\underset{¯}{\overset{}{x}}}\left(y_2\|\stackrel{}{x}_2\right)}\delta (\stackrel{}{x}_3,\stackrel{}{x}_2)`$ \| \| \| $`\underset{¯}{\overset{}{x}}_4`$ \| $`\stackrel{}{x}_4Bool^{N_B}`$ \| $`\delta (\stackrel{}{x}_4,\stackrel{}{x}_3)`$ \| \| \| $`\underset{¯}{y}_1`$ \| $`y_1Bool`$ \| $`\delta (y_1,0)`$ \| \| \| $`\underset{¯}{y}_3`$ \| $`y_3Bool`$ \| $`\delta (y_3,y_2)`$ \| \| Table 13 According to Table 6, the probability amplitude for the leaf (external) nodes is $`A(\stackrel{}{x}_4,y_3)=`$ (131a) $`=`$ $`{\displaystyle \underset{}{}}\sqrt{P_{\underset{¯}{\overset{}{x}}}(\stackrel{}{\text{}}_2)P_{\underset{¯}{y}\|\underset{¯}{\overset{}{x}}}(\text{}_2\|\stackrel{}{\text{}}_2)}\theta (\text{}_2=y_3)\theta (\stackrel{}{\text{}}_2=\stackrel{}{\text{}}_3=\stackrel{}{x}_4)\theta (\stackrel{}{\text{}}_1=\text{}_1=0)`$ $`=`$ $`\sqrt{P_{\underset{¯}{\overset{}{x}}}(\stackrel{}{x}_4)P_{\underset{¯}{y}\|\underset{¯}{\overset{}{x}}}(y_3\|\stackrel{}{x}_4)}.`$ (131b) To fully specify the QB net for voting, we need to extend $`A(\stackrel{}{x}_2\|\stackrel{}{x}_1=0)`$ and $`A(\stackrel{}{x}_3,y_2\|\stackrel{}{x}_2,y_1=0)`$ into unitary matrices by adding columns to them. This can always be accomplished by applying the Gram-Schmidt algorithm. But sometimes one can guess a matrix extension and applying Gram-Schmidt becomes unnecessary. If $`P_{\underset{¯}{\overset{}{x}}}`$ is uniform (i.e., $`P(\stackrel{}{x})=1/N_S`$ for all $`\stackrel{}{x}`$, which means there is no a priori information about $`\stackrel{}{x}`$), then $`A(\stackrel{}{x}_2\|\stackrel{}{x}_1=0)=1/\sqrt{N_S}`$. In this case, we can extend $`A(\stackrel{}{x}_2\|\stackrel{}{x}_1=0)`$ into the unitary matrix $$[A(\stackrel{}{x}_2\|\stackrel{}{x}_1)]=\widehat{H}_{N_B}.$$ (132) (This works because all entries of the first column of $`\widehat{H}_{N_B}`$ are equal to $`1/\sqrt{N_S}`$.) As to extending $`A(\stackrel{}{x}_3,y_2\|\stackrel{}{x}_2,y_1=0)`$, this can be done as follows. Define $$\mathrm{\Delta }_p=diag(\sqrt{p_0},\sqrt{p_1},\mathrm{},\sqrt{p_{N_S1}}),$$ (133) and $$\mathrm{\Delta }_q=diag(\sqrt{q_0},\sqrt{q_1},\mathrm{},\sqrt{q_{N_S1}}).$$ (134) A possible way of extending $`A(\stackrel{}{x}_3,y_2\|\stackrel{}{x}_2,y_1=0)`$ into a unitary matrix is $$[A(\stackrel{}{x}_3,y_2\|\stackrel{}{x}_2,y_1)]=\left(\begin{array}{cc}\mathrm{\Delta }_p& \mathrm{\Delta }_q\\ \mathrm{\Delta }_q& \mathrm{\Delta }_p\end{array}\right).$$ (135) Unitary matrices of this kind are called D-matrices in Ref.. Ref. shows how to decompose any D-matrix into a SEO. Earlier, we explained how to estimate a conditional probability for a CB net by running a QB net $`\nu `$ times on a quantum computer. If we wanted to find $`P(y\|x^0,x^1)`$ for the voting CB net, then the number of runs $`\nu `$ required to estimate $`P(y\|x^0,x^1)`$ with moderate accuracy would not be too onerous, because the domain of $`P(y\|x^0,x^1)`$ is $`Bool^3`$, which contains only 8 points. But what if we wanted to estimate $`P(y\|\stackrel{}{x})`$? For large $`N_B`$, the domain of $`P(y\|\stackrel{}{x})`$ is very large ($`2^{N_B+1}`$ points). If the support of $`P(y\|\stackrel{}{x})`$ occupies a large fraction of this domain, then the number of runs $`\nu `$ required to estimate $`P(y\|\stackrel{}{x})`$ with moderate accuracy is forbiddingly large. However, there are some cases in which “Grover’s Microscope” can come to the rescue, by allowing us to amplify certain salient features of $`P(y\|\stackrel{}{x})`$ so that they become measurable in only a few runs. Next we will discuss Grover’s Microscope for the voting QB net defined by Fig.17 and Table 6. For simplicity, we will assume that $`P_{\underset{¯}{\overset{}{x}}}`$ is uniform. Let $`\stackrel{}{\kappa }=(\kappa _0,\kappa _1,\mathrm{},\kappa _{N_B1})`$ label $`N_B`$ bits and let $`\tau `$ label another bit. Assume that $`\tau `$ and all the $`\kappa _i`$ are distinct. Define $$\|\varphi _p_\stackrel{}{\kappa }=\varphi _p=(\sqrt{p_0},\sqrt{p_1},\mathrm{},\sqrt{p_{N_S1}})^T,$$ (136) $$\|\varphi _q_\stackrel{}{\kappa }=\varphi _q=(\sqrt{q_0},\sqrt{q_1},\mathrm{},\sqrt{q_{N_S1}})^T,$$ (137) and $`\|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_S}}}\left(\|0_\tau \|\varphi _p_\stackrel{}{\kappa }+\|1_\tau \|\varphi _q_\stackrel{}{\kappa }\right)`$ (144) $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_S}}}\left[\left(\begin{array}{c}1\\ 0\end{array}\right)\varphi _p+\left(\begin{array}{c}0\\ 1\end{array}\right)\varphi _q\right]={\displaystyle \frac{1}{\sqrt{N_S}}}\left(\begin{array}{c}\varphi _p\\ \varphi _q\end{array}\right)=\mathrm{\Psi }.`$ Since $`p_i+q_i=1`$ for all $`i`$, $`\varphi _p^T\varphi _p+\varphi _q^T\varphi _q=N_S`$. According to Eq.(131), when $`P_{\underset{¯}{\overset{}{x}}}`$ is uniform, the voting QB net fully specifies a unitary matrix $`U_{net}`$ such that $$\|\mathrm{\Psi }=U_{net}\|0_\stackrel{}{\kappa }\|0_\tau .$$ (145) Define orthonormal vectors $`e_0`$ and $`e_1`$ by $$e_0=\left(\begin{array}{c}\widehat{\varphi }_p\\ 0\end{array}\right),e_1=\left(\begin{array}{c}0\\ \widehat{\varphi }_q\end{array}\right),$$ (146) where $`\widehat{V}`$ is a unit vector in the direction of $`\stackrel{}{V}`$. If $`P(y\|\stackrel{}{x})`$ is deterministic with AND-like $`f`$, then all components of $`e_0`$ are zero except for the one at the target state $`j_{targ}`$. In terms of $`e_0,e_1`$, $`\mathrm{\Psi }`$ can be expressed as $$\mathrm{\Psi }=\frac{1}{\sqrt{N_S}}\left(\begin{array}{c}\varphi _p\\ \varphi _q\end{array}\right)=\frac{1}{\sqrt{N_S}}(\|\varphi _p\|e_0+\|\varphi _q\|e_1).$$ (147) It is convenient to define a vector $`\mathrm{\Psi }_{}`$ orthogonal to $`\mathrm{\Psi }`$: $$\mathrm{\Psi }_{}=\frac{1}{\sqrt{N_S}}(\|\varphi _q\|e_0\|\varphi _p\|e_1).$$ (148) If $`P(y\|\stackrel{}{x})`$ is deterministic with AND-like $`f`$, then $`\|\varphi _p\|=1`$ and $`\|\varphi _q\|=\sqrt{N_S1}`$ so, for large $`N_S`$, $`\mathrm{\Psi }e_1`$ and $`\mathrm{\Psi }_{}e_0`$. For an arbitrary angle $`\alpha `$, let $$\mathrm{\Psi }_{}^{}=\frac{1}{\sqrt{N_S}}\left[(c_{\frac{\alpha }{2}}\|\varphi _q\|+s_{\frac{\alpha }{2}}\|\varphi _p\|)e_0+(s_{\frac{\alpha }{2}}\|\varphi _q\|c_{\frac{\alpha }{2}}\|\varphi _p\|)e_1\right],$$ (149) where $`s_A=\mathrm{sin}A`$ and $`c_A=\mathrm{cos}A`$ for any angle $`A`$. Let $`\mathrm{}(x,y)`$ denote the angle between 2 vectors $`x`$ and $`y`$. Note that $`\mathrm{}(\mathrm{\Psi }_{}^{},\mathrm{\Psi }_{})=\alpha /2`$. We define $`\mathrm{}(e_1,\mathrm{\Psi })=\theta /2`$. Fig.18 portrays various vectors that arise in explaining Grover’s Microscope. Note that $`\mathrm{\Psi }_{}^{}=e_0`$ when $`\alpha =\theta `$. Since we plan to stay within the two dimensional vector space with orthonormal basis $`e_0,e_1`$, it is convenient to switch matrix representations. Within $`span(e_0,e_1)`$, $`e_0,e_1`$ can be represented more simply by: $$e_0=\left(\begin{array}{c}1\\ 0\end{array}\right),e_1=\left(\begin{array}{c}0\\ 1\end{array}\right).$$ (150) If $`e_0,e_1`$ are represented in this way, then $$\mathrm{\Psi }=\frac{1}{\sqrt{N_S}}\left(\begin{array}{c}\|\varphi _p\|\\ \|\varphi _q\|\end{array}\right),$$ (151) $$\mathrm{\Psi }_{}=\frac{1}{\sqrt{N_S}}\left(\begin{array}{c}\|\varphi _q\|\\ \|\varphi _p\|\end{array}\right),$$ (152) and $$\mathrm{\Psi }_{}^{}=W\mathrm{\Psi },\text{where}W=\left(\begin{array}{cc}c_{\frac{\alpha }{2}}& s_{\frac{\alpha }{2}}\\ s_{\frac{\alpha }{2}}& c_{\frac{\alpha }{2}}\end{array}\right)\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$ (153) The matrix $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ is a clockwise rotation by $`\pi /2`$ in space $`span(e_0,e_1)`$. Thus, $`W`$ equals a clockwise rotation by $`\pi /2`$ followed by a counter-clockwise rotation by $`\alpha /2`$. Define the following reflection operators $$R_0=12\mathrm{\Pi }_{\|0_\stackrel{}{\kappa }}\mathrm{\Pi }_{\|0_\tau }=(1)^{\mathrm{\Pi }_{\|0_\stackrel{}{\kappa }}\mathrm{\Pi }_{\|0_\tau }},$$ (154) $$R_\mathrm{\Psi }=U_{net}R_0U_{net}^{},$$ (155) $$R_\mathrm{\Psi }_{}^{}=WR_\mathrm{\Psi }W^{}=WU_{net}R_0U_{net}^{}W^{}.$$ (156) From Eq.(24), it follows that $$R_\mathrm{\Psi }R_\mathrm{\Psi }_{}^{}=c_\alpha \mathrm{\Psi }\mathrm{\Psi }^Ts_\alpha \mathrm{\Psi }\mathrm{\Psi }_{}^T+s_\alpha \mathrm{\Psi }_{}\mathrm{\Psi }^T+c_\alpha \mathrm{\Psi }_{}\mathrm{\Psi }_{}^T.$$ (157) Thus, $`R_\mathrm{\Psi }R_\mathrm{\Psi }_{}^{}`$ rotates vectors in $`span(e_0,e_1)`$, clockwise by an angle $`\alpha `$. Grover’s Microscope can be summarized by the following equation $$(R_\mathrm{\Psi }R_\mathrm{\Psi }_{}^{})^r\mathrm{\Psi }e_0,$$ (158) for some integer $`r`$ to be determined, where “$``$” means approximation at large $`N_S`$. What this means is that our system starts in state $`\mathrm{\Psi }`$ and is rotated consecutively $`r`$ times, each time by a small angle $`\alpha `$, until it arrives at the state $`e_0`$. If $`P(y\|\stackrel{}{x})`$ is deterministic with AND-like $`f`$, then measuring state $`e_0`$ yields the target state $`j_{targ}`$. The optimum number $`r`$ of iterations is $$r\alpha \frac{\pi }{2}(1+2k)$$ (159) for some integer $`k`$. Note that $`\mathrm{cos}(\theta /2)=\mathrm{\Psi }\|e_1=\|\varphi _q\|/\sqrt{N_S}`$ so, in general, $`\theta `$ depends on $`\|\varphi _p\|`$ (or on $`\|\varphi _q\|=\sqrt{N_S\|\varphi _p\|^2}`$). If $`P(y\|\stackrel{}{x})`$ is deterministic with AND-like $`f`$, then $`\|\varphi _p\|=1`$ and $`\|\varphi _q\|=\sqrt{N_S1}`$. In this case, it is convenient to choose $`\alpha =\theta `$, so that $`\mathrm{\Psi }_{}^{}=e_1`$ and Figs.6 and 18 become the same diagram under the mapping $`\mathrm{\Psi }\mu `$ and $`\mathrm{\Psi }_{}^{}\varphi =e_0`$. Then the optimum number $`r`$ of iterations for Grover’s original algorithm and for Grover’s Microscope are equal. If we don’t know ahead of time the value of $`\|\varphi _p\|`$, then setting $`\theta =\alpha `$ will make both $`r`$ and $`\alpha `$ depend on the unknown $`\|\varphi _p\|`$, although the product $`r\alpha `$ will still be independent of it. Let $`U_{\mu scope}`$ $`=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ $`=`$ $`e_1e_0^T+e_0e_1^T`$ $`=`$ $`\mathrm{\Psi }\mathrm{\Psi }_{}^T+\mathrm{\Psi }_{}\mathrm{\Psi }^T.`$ (163) Note that $$U_{\mu scope}\mathrm{\Psi }=\mathrm{\Psi }_{}.$$ (164) From the point of view of quantum compiling, Grover’s Microscope approximates the $`\pi /2`$ rotation $`U_{\mu scope}`$ by the $`r`$-fold product of $`R_\mathrm{\Psi }R_\mathrm{\Psi }_{}^{}`$, where we assume that $`R_\mathrm{\Psi }R_\mathrm{\Psi }_{}^{}`$ can be shown to have a SEO of low (polynomial in $`N_B`$) complexity. (If such a low complexity SEO cannot be found, then it is pointless to divide $`U_{\mu scope}`$ into $`r`$ iterations of $`R_\mathrm{\Psi }R_\mathrm{\Psi }_{}^{}`$, and we might be better off compiling $`U_{\mu scope}`$ all at once.) ## Appendix A Appendix: Deterministic $`pd(Bool\|Bool^{N_B})`$ matrices In this Appendix, we will first define a special kind of probability matrices which we call deterministic $`pd(Bool\|Bool^{N_B})`$ matrices. Then we will show how such probability matrices can be q-embedded, and how their q-embedding can be expressed as a SEO. Suppose $`yBool`$ and $`\stackrel{}{x}=(x^0,x^1,\mathrm{},x^{N_B1})Bool^{N_B}`$. Let $`f:Bool^{N_B}Bool`$. We will say that $`f`$ is AND-like if $`f(\stackrel{}{x})=\theta (\stackrel{}{x}=\stackrel{}{x}_{targ})`$ for some target vector $`\stackrel{}{x}_{targ}Bool^{N_B}`$. An AND-like $`f`$ maps all $`\stackrel{}{x}`$ into zero except for $`\stackrel{}{x}_{targ}`$ which it maps into one. Thus, $`\|f^1(1)\|=1`$. An example of an AND-like $`f`$ is the multiple AND gate $`f(\stackrel{}{x})=x^0x^1\mathrm{}x^{N_B1}`$, which can also be expressed as $`f(\stackrel{}{x})=\theta [\stackrel{}{x}=(1,1,\mathrm{},1)]`$. We will say that $`f`$ is OR-like if $`f(\stackrel{}{x})=\theta (\stackrel{}{x}\stackrel{}{x}_{targ})`$ for some target vector $`\stackrel{}{x}_{targ}Bool^{N_B}`$. An OR-like $`f`$ maps all $`\stackrel{}{x}`$ into one except for $`\stackrel{}{x}_{targ}`$ which it maps into zero. Thus, $`\|f^1(0)\|=1`$. An example of an OR-like $`f`$ is the multiple OR gate $`f(\stackrel{}{x})=x^0x^1\mathrm{}x^{N_B1}`$, which can also be expressed as $`f(\stackrel{}{x})=\theta [\stackrel{}{x}(0,0,\mathrm{},0)]`$. We will say that $`f`$ has a single target if it is either AND-like or OR-like. If $`f`$ has more than one target (i.e., if $`\|f^1(0)\|`$ and $`\|f^1(1)\|`$ are both greater than one), then we will say that $`f`$ has multiple targets. Suppose $`yBool`$ and $`\stackrel{}{x}=(x^0,x^1,\mathrm{},x^{N_B1})Bool^{N_B}`$. Let $`f:Bool^{N_B}Bool`$. In this section, we consider deterministic $`pd(Bool\|Bool^{N_B})`$ matrices; that is, probability matrices of the form $`P(y\|\stackrel{}{x})=\delta (y,f(\stackrel{}{x}))`$. First let us consider the case that $`f`$ has a single target. For example, for $`N_B=2`$, if $`f`$ is an AND gate $$P(y\|\stackrel{}{x})_{AND}=\{\begin{array}{c}\hfill (x^0,x^1)\\ \hfill \begin{array}{cccccc}& & 00\hfill & 01\hfill & 10\hfill & 11\hfill \\ & & & & & \\ y\hfill & 0\hfill & 1\hfill & 1\hfill & 1\hfill & 0\hfill \\ & 1\hfill & 0\hfill & 0\hfill & 0\hfill & 1\hfill \end{array}\end{array},$$ (165) and if $`f`$ is an OR gate $$P(y\|\stackrel{}{x})_{OR}=\{\begin{array}{c}\hfill (x^0,x^1)\\ \hfill \begin{array}{cccccc}& & 00\hfill & 01\hfill & 10\hfill & 11\hfill \\ & & & & & \\ y\hfill & 0\hfill & 1\hfill & 0\hfill & 0\hfill & 0\hfill \\ & 1\hfill & 0\hfill & 1\hfill & 1\hfill & 1\hfill \end{array}\end{array}.$$ (166) Suppose bit value $`y`$ is stored in the bit labelled $`\tau `$. And suppose bit values $`x^0,x^1,\mathrm{},x^{N_B1}`$ are stored in the bits labelled $`\stackrel{}{\kappa }=(\kappa _0,\kappa _1,\mathrm{},\kappa _{N_B1})`$. Define $`e_j`$ for all $`jZ_{0,N_S1}`$ to be the $`N_S`$ dimensional column vector with $`j`$th component equal to one and all other components equal to zero. Let $`\mathrm{\Pi }_j=e_je_j^T`$ and $`\mathrm{\Pi }_{targ}=\mathrm{\Pi }_{j_{targ}}`$, where $`j_{targ}Z_{0,N_S1}`$ is the target state. $`\mathrm{\Pi }_{targ}`$ can expressed as product of number operators. Indeed, if $$j_{targ}=\underset{i=0}{\overset{N_B1}{}}x_{targ,i}2^i,$$ (167) then $$\mathrm{\Pi }_{targ}=\mathrm{\Pi }_{j_{targ}}=\underset{i=0}{\overset{N_B1}{}}\left[n(\kappa _i)\theta (x_{targ,i}=1)+\overline{n}(\kappa _i)\theta (x_{targ,i}=0)\right].$$ (168) For example, if $`j_{targ}=0`$ then $`\mathrm{\Pi }_{targ}=\overline{n}(\kappa _0)\overline{n}(\kappa _1)\mathrm{}\overline{n}(\kappa _{N_B1})`$. An AND-like probability matrix $`P(y\|\stackrel{}{x})`$ is q-embedded within the unitary matrix $$U_{ANDlike}=[A(y,\stackrel{}{\stackrel{~}{x}}\|\stackrel{~}{y},\stackrel{}{x})]=\begin{array}{ccc}& \stackrel{~}{y}=0\hfill & \stackrel{~}{y}=1\hfill \\ & & \\ y=0\hfill & 1\mathrm{\Pi }_{targ}\hfill & \mathrm{\Pi }_{targ}\hfill \\ y=1\hfill & \mathrm{\Pi }_{targ}\hfill & 1\mathrm{\Pi }_{targ}\hfill \end{array}.$$ (169) Note that $`U_{ANDlike}`$ $`=`$ $`1+\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)\mathrm{\Pi }_{targ}`$ (170c) $`=`$ $`1+\mathrm{\Pi }_{targ}(\stackrel{}{\kappa })(i\sigma _y(\tau )1)`$ (170d) $`=`$ $`[i\sigma _y(\tau )]^{\mathrm{\Pi }_{targ}(\stackrel{}{\kappa })}.`$ (170e) Eqs.(168) and (170e) show how to express $`U_{ANDlike}`$ as a qubit rotation with multiple control qubits. Operations of this kind can be decomposed into a SEO using the techniques of Refs. and . An OR-like probability matrix $`P(y\|\stackrel{}{x})`$ is q-embedded within the unitary matrix $$U_{ORlike}=[A(y,\stackrel{}{\stackrel{~}{x}}\|\stackrel{~}{y},\stackrel{}{x})]=\begin{array}{ccc}& \stackrel{~}{y}=0\hfill & \stackrel{~}{y}=1\hfill \\ & & \\ y=0\hfill & \mathrm{\Pi }_{targ}\hfill & 1\mathrm{\Pi }_{targ}\hfill \\ y=1\hfill & 1\mathrm{\Pi }_{targ}\hfill & \mathrm{\Pi }_{targ}\hfill \end{array}.$$ (171) Note that $`U_{ORlike}`$ $`=`$ $`\left(\begin{array}{cc}0& I_{N_S}\\ I_{N_S}& 0\end{array}\right)\left(\begin{array}{cc}1\mathrm{\Pi }_{targ}& \mathrm{\Pi }_{targ}\\ \mathrm{\Pi }_{targ}& 1\mathrm{\Pi }_{targ}\end{array}\right)`$ (172e) $`=`$ $`\sigma _x(\tau )[i\sigma _y(\tau )]^{\mathrm{\Pi }_{targ}(\stackrel{}{\kappa })}.`$ (172f) Finally, let us consider the case when $`f:Bool^{N_B}Bool`$ has multiple targets. Let $`TZ_{0,N_S1}`$ be the set of these targets; i.e., either $`T=f^1(0)`$ or $`T=f^1(1)`$. Define $`\mathrm{\Pi }_{targ}`$ by $$\mathrm{\Pi }_{targ}=\underset{jT}{}\mathrm{\Pi }_j.$$ (173) $`\mathrm{\Pi }_{targ}`$ can be expressed as a product of number operators. Indeed, each $`\mathrm{\Pi }_j`$ on the right hand side of Eq.(173) can be separately expressed, using Eq.(168), as a product of number operators. If $`T=f^1(1)`$, then $`P(y\|\stackrel{}{x})`$ is q-embedded within the unitary matrix $`U_{multitarg}`$ $`=`$ $`[A(y,\stackrel{}{\stackrel{~}{x}}\|\stackrel{~}{y},\stackrel{}{x})]=\begin{array}{ccc}& \stackrel{~}{y}=0\hfill & \stackrel{~}{y}=1\hfill \\ & & \multicolumn{-1}{c}{}\\ y=0\hfill & 1\mathrm{\Pi }_{targ}\hfill & \mathrm{\Pi }_{targ}\hfill \\ y=1\hfill & \mathrm{\Pi }_{targ}\hfill & 1\mathrm{\Pi }_{targ}\hfill \end{array}`$ (177) $`=`$ $`[i\sigma _y(\tau )]^{\mathrm{\Pi }_{targ}(\stackrel{}{\kappa })}.`$ (178) ## Appendix B Appendix: Quasi-deterministic $`pd(Bool\|Bool^{N_B})`$ matrices In this Appendix, we will first define a special kind of probability matrices which we call quasi-deterministic $`pd(Bool\|Bool^{N_B})`$ matrices. Then we will show how such probability matrices can be q-embedded, and how their q-embedding can be expressed as a SEO. Suppose $`yBool`$ and $`\stackrel{}{x}=(x^0,x^1,\mathrm{},x^{N_B1})Bool^{N_B}`$. Let $`f:Bool^{N_B}Bool`$. In the previous appendix, we considered deterministic $`pd(Bool\|Bool^{N_B})`$ matrices; that is, probability matrices of the form $`P(y\|\stackrel{}{x})=\delta (y,f(\stackrel{}{x}))`$. In this section, we will consider quasi-deterministic $`pd(Bool\|Bool^{N_B})`$ matrices; that is, probability matrices of the form $$P(y\|\stackrel{}{x})=\underset{\stackrel{}{t}}{}\delta (y,f(\stackrel{}{t}))P(t^0\|x^0)P(t^1\|x^1)\mathrm{}P(t^{N_B1}\|x^{N_B1}),$$ (179) where we sum over all $`\stackrel{}{t}=(t^0,t^1,\mathrm{},t^{N_B1})Bool^{N_B}`$. Fig.19 shows a CB net representation of Eq.(179). Examples of quasi-deterministic $`pd(Bool\|Bool^{N_B})`$ matrices are: (1)the noisy OR, for which $`f(\stackrel{}{t})=t^0t^1\mathrm{}t^{N_B1}`$; (2)the noisy AND, for which $`f(\stackrel{}{t})=t^0t^1\mathrm{}t^{N_B1}`$; (3)the noisy CNOT, for which $`f(\stackrel{}{t})=t^0t^1\mathrm{}t^{N_B1}`$, etc. For each $`\alpha Z_{0,N_B1}`$, the probabilities $`P(\underset{¯}{t}^\alpha =t\|\underset{¯}{x}^\alpha =x)`$ will be abbreviated by $`p_{t,x}^\alpha `$ for $`t,xBool`$. $`P(\underset{¯}{t}^\alpha =t\|\underset{¯}{x}^\alpha =x)`$ has two independent parameters which we may take to be $`p_{01}^\alpha `$ (the probability of false negatives) and $`p_{10}^\alpha `$ (the probability of false positives). $`p_{00}^\alpha `$ and $`p_{11}^\alpha `$ can be expressed in terms of these independent parameters: $`p_{00}^\alpha =1p_{10}^\alpha `$, $`p_{11}^\alpha =1p_{01}^\alpha `$. Whereas a completely general probability matrix $`P(y\|\stackrel{}{x})pd(Bool\|Bool^{N_B})`$ has $`2^{N_B}`$ free parameters, a quasi-deterministic $`P(y\|\stackrel{}{x})`$ has $`2N_B`$ free parameters. Rather than q-embedding the probability matrix $`P(y\|\stackrel{}{x})`$ as a whole, it is convenient to q-embed separately the probability matrices $`P(y\|\stackrel{}{t})`$ and $`P(t^\alpha \|x^\alpha )`$ for every $`\alpha Z_{0,N_B1}`$. $`P(y\|\stackrel{}{t})=\delta (y,f(\stackrel{}{t}))`$ is a deterministic $`pd(Bool\|Bool^{N_B})`$ matrix so its q-embedding is discussed in Appendix A. As for $`P(t^\alpha \|x^\alpha )`$, it can be easily q-embedded as follows. For each $`\alpha Z_{0,N_B1}`$, let $$\mathrm{\Delta }_p^\alpha =\left(\begin{array}{cc}\sqrt{p_{00}^\alpha }& 0\\ 0& \sqrt{p_{01}^\alpha }\end{array}\right),\mathrm{\Delta }_q^\alpha =\left(\begin{array}{cc}\sqrt{p_{10}^\alpha }& 0\\ 0& \sqrt{p_{11}^\alpha }\end{array}\right).$$ (180) $`P(t^\alpha \|x^\alpha )`$ is q-embedded within the unitary matrix: $$[A(t^\alpha ,\stackrel{~}{x}^\alpha \|\stackrel{~}{t}^\alpha ,x^\alpha )]=\left(\begin{array}{cc}\mathrm{\Delta }_p^\alpha & \mathrm{\Delta }_q^\alpha \\ \mathrm{\Delta }_q^\alpha & \mathrm{\Delta }_p^\alpha \end{array}\right).$$ (181) Unitary matrices of this kind are called D-matrices in Ref.. Ref. shows how to decompose any D-matrix into a SEO.
warning/0004/physics0004015.html	ar5iv	text	# Tumour Therapy with Particle Beams ## Introduction It has been known for a long time that tissue, in particular tumour tissue, is sensitive to ionising radiation. Therefore it is only natural that tumours have been treated with various types of radiation, like $`\gamma `$-rays and electrons. $`\gamma `$-rays are easily available from radioactive sources, like <sup>60</sup>Co, and electrons can be accelerated to $`MeV`$-energies by relatively inexpensive linear accelerators. The disadvantage of $`\gamma `$-rays and electrons is that they deposit most of their energy close to the surface. To reduce the surface dose in tumour treatment requires rotating the source or the patient so that the surface dose is distributed over a larger volume. In contrast, protons and heavy ions deposit most of their energy close to the end of their range (Bragg-peak). The increase in energy loss at the Bragg-peak amounts to a factor of about 5 compared to the surface dose, depending somewhat on the particle’s energy. Heavy ions offer, in addition, the possibility to monitor the destructive power of the beam by observing annihilation radiation by standard positron-emission tomography techniques (PET). The annihilation radiation is emitted by $`\beta ^+`$-active nuclear fragments produced by the incident heavy ion beam itself. ## Energy loss of particles in tissue A photon beam is attenuated in matter according to $$I(x)=I_0e^{\mu x}$$ (1) where $`I_0`$ is the initial intensity and $`I(x)`$ the beam intensity at the depth $`x`$. $`\mu `$ is the linear mass attenuation coefficient which depends on the photon energy $`E`$ and the target charge $`Z`$. $`\mu (E)`$ is shown in figure 1 for a target composed of water, which is essentially equivalent to tissue. The main interaction mechanisms which contribute to $`\mu (E)`$ are the photoelectric effect $`(Z^5/E^{3.5})`$, Compton scattering $`((Z/E)lnE)`$ and pair-production $`(Z^2lnE)`$. For energies typical for radioactive sources $`(MeV)`$ Compton scattering dominates. The absorption profile of photons in matter exhibits a peak close to the surface followed by an exponential decay. Charged particles suffer energy loss by ionisation. This energy loss is described by the Bethe-Bloch formula: $$\frac{dE}{dx}=2\kappa \{ln\frac{E_{\mathrm{kin}}^{\mathrm{max}}}{I}\beta ^2\frac{\delta }{2}\}$$ (2) where $$\kappa =2\pi N_Ar_e^2m_ec^2z^2\frac{Z}{A}\frac{1}{\beta ^2}.$$ (3) \| $`z`$ \| \| charge of the beam particle \| \| --- \| --- \| --- \| \| $`Z`$ \| \| charge of the absorber material \| \| $`A`$ \| \| mass number of the absorber material \| \| $`m_e`$ \| \| electron mass \| \| $`c`$ \| \| velocity of light \| \| $`N_A`$ \| \| Avogadro’s number \| \| $`r_e`$ \| \| classical electron radius \| \| $`\beta `$ \| \| velocity of the particle divided by $`c`$ \| \| $`E_{\mathrm{kin}}^{\mathrm{max}}`$ \| \| maximum transferable energy \| \| \| \| to an atomic electron \| \| $`I`$ \| \| mean excitation energy of the target material \| \| $`\delta `$ \| \| density parameter \| For protons $`(z=1)`$ interacting in water (or tissue) equation (2) can be approximated by $$\frac{dE}{dx}=0,16\frac{1}{\beta ^2}ln\frac{E_{\mathrm{kin}}^{\mathrm{max}}[eV]}{100}\left[\frac{MeV}{cm}\right]$$ (4) where $$E_{\mathrm{kin}}^{\mathrm{max}}2m_ec^2\beta ^2\gamma ^2,$$ (5) which gives an energy loss of 4.2 $`MeV/cm`$ for 200 $`MeV`$ protons at the surface and $`20MeV/cm`$ close to the end of their range. For heavy ions the energy loss is essentially scaled by $`z^2`$. When charged particles reach the end of their range the energy loss first rises like $`1/\beta ^2`$ but when they are very slow they capture electrons from the target material and their effective charge decreases and hence their energy loss rapidly falls to zero. A typical energy loss curve for ions as a function of their energy is sketched in figure 2 Kraft1 . The energy loss of <sup>12</sup>C ions as a function of the depth in water is shown in figure 3 Kraft1 ; Kraft2 . The tail of the energy loss beyond the Bragg-peak originates from fragmentation products of <sup>12</sup>C ions, which are faster than the <sup>12</sup>C ions and have a somewhat longer range. In the ionisation process a generally small fraction of the particle’s energy is transferred to the atomic electrons. In rare cases these electrons can get a larger amount of energy. The $`\delta `$-electrons deviate from the main ionisation trail and produce a fuzzy-like track (figure 4, Kraft1 ). In addition to ionisation light particles, like electrons, can also undergo bremsstrahlung $`(dE/dxz^2Z^2E)`$. Since the probability for this process is inversely proportional to the square of the mass of the beam particle, bremsstrahlung can be neglected for particles heavier than the electron for energies relevant to tumour therapy Grupen1 . The above mentioned fragmentation of heavy ions leads to the production of positron emitters. For the <sup>12</sup>C case, lighter isotopes like <sup>11</sup>C and <sup>10</sup>C are produced. Both isotopes decay with short half-lives $`(T_{1/2}(^{11}C)=20,38min`$; $`T_{1/2}(^{10}C)=19.3s)`$ to boron according to $`{}_{}{}^{11}C`$ $``$ $`{}_{}{}^{11}B+e^++\nu _e`$ (6) $`{}_{}{}^{10}C`$ $``$ $`{}_{}{}^{10}B+e^++\nu _e.`$ The positrons have a very short range, typically below $`1mm`$. After coming to rest they annihilate with electrons of the tissue giving off two monochromatic photons of $`511keV`$ which are emitted back-to-back $$e^++e^{}\gamma +\gamma .$$ (7) These photons can be detected by positron-emission tomography techniques and can be used to monitor the destructive effect of heavy ions on the tumour tissue. ## Production of particle beams The treatment of deep seated tumours requires charged particles of typically 100 to $`400MeV`$ per nucleon, i.e. 100 to $`400MeV`$ protons or 1.2 to $`4.8GeV^{12}`$C ions. These particles are accelerated in either a linear accelerator or in a synchrotron. As an example figure 5 shows a typical set-up for the production of heavy ions. <sup>12</sup>C atoms are evaporated from an ion source and pre-accelerated. Thin foils are used to strip off all electrons from the ions. The <sup>12</sup>C nuclei are then injected into a synchrotron, where they are accelerated by radiofrequency cavities to the desired energy. The ions are kept on track by dipole bending magnets and they are focussed by quadrupoles. After having reached the final energy they are ejected by a kicker magnet, which directs the particles to the treatment room. Their path is monitored by tracking chambers (multi-wire proportional counteres, ion chambers or drift-chambers). If beam losses occur veto-counters (mostly scintillation counters) ensure that only a pencil beam is steered to the treatment room. Nowadays, mainly protons and heavy ions are used for tumour therapy. Other possibilities consist of the use of negative pions Dyson ; Curtis ; Goodman , which are produced by high energy protons in a beam dump according to $$p+\mathrm{nucleus}p+\mathrm{nucleus}+\pi ^{}+\pi ^++\pi ^0$$ (8) where the $`\pi ^{}`$ are momentum selected and collimated. After losing their energy by ionisation the negative pions are captured in the tumour tissue by nuclei at the end of their range and produce so-called ‘stars’ in which neutrons are created. The Bragg-peak of the negative pions along with the local production of neutrons which have a high biological effectiveness leads to an efficient cell killing in the tumour at the end of the pion’s range. Neutrons are also possible candidates for tumour treatment Lennox . For this purpose the tumour is sensitized by a boron compound before neutron treatment. The boron compound must be selected in such a way that it is preferentially deposited in the tumour region. Neutrons are then captured by the boron according to: $$n+^{10}B^7Li+\alpha .$$ (9) The produced $`\alpha `$-particles (He-nuclei) have a very short range ($``$ several $`\mu m`$) and high biological effectiveness. Best results are obtained with epithermal neutrons ($`1keV`$) produced by $`5MeV`$ protons on light targets (e.g. Be). Direct irradiation with neutrons – without sensitizing the tumour – has the disadvantage that neutrons show a similar dose depth curve like <sup>60</sup>Co $`\gamma `$-rays thus producing a high amount of biologically very effective damage in the healthy tissue around the tumour (see figure 6 NAC ). ## Applications in Tumour Therapy The target for cell killing is the DNA in the cell nucleus (see figure 7 (after Kraft1 )). The size of the DNA-molecule compares favorably well with the width of the ionisation track of a heavy ion. The DNA contains two strands containing identical information. A damage of one strand by ionising radiation can easily be repaired by copying the information from the unaffected strand to the damaged one. Therefore the high ionisation density at the end of a particle’s range matches well with the requirement to produce double strand breaks in the DNA, which the cell will not survive. Heavy ions like <sup>12</sup>C seem to be optimal for this purpose. Ions heavier than carbon would even be more powerful in destroying tumour tissue, however, their energy loss in the surrounding tissue and in the entrance region already reaches a level where the fraction of irreparable damage is too high, while for lighter ions (like <sup>12</sup>C) mostly repairable damage is produced in the healthy tissue outside the targeted tumour. The cell killing rate in the tumour region thus benefits from two properties of protons or ions like carbon: * the increased energy loss of protons and ions at the end of their range and * the increased biological effectiveness of double strand breaks at high ionisation density. The cell killing rate is eventually related to the equivalent dose H in the tumour region, which can be expressed by $$H=\frac{1}{m}\frac{dE}{dx}𝑑xRBE$$ (10) where $`m`$ is the tumour mass and RBE the increased relative biological effectiveness. The integral extends over the tumour region. As mentioned above the rate and location of cell killing can be monitored by observing the annihilation photons which result from the $`\beta ^+`$-decay of fragments formed by the beam. These physical and biological principles are employed in an efficient way by the raster scan technique Kraft2 ; Kraft3 ; Kraft4 . A pencil beam of heavy ions (diameter $`1mm`$) is aimed at the tumour. The beam location and spread is monitored by tracking chambers with high spatial resolution. In the treatment planning the tumour is subdivided into three-dimensional pixels (“voxels”). Then the dose required to destroy the tumour, which is proportional to the beam intensity, is calculated for every voxel. For a fixed depth in tissue an areal scan is performed by magnetic deflection sweeping the beam across the area in a similar way as a TV image is produced (see figure 8, Kraft3 ; Kraft4 ). The tumour volume is filled from the back by energy variation ($``$ range variation) of the beam. Typically 50 energy steps are used starting at the rear plane. For a depth profile from $`2cm`$ to $`30cm`$ one has to cover energies from $`80MeV`$/nucleon to $`430MeV`$/nucleon. When the beam energy is reduced the required dose for the plane under irradiation is calculated using the damage that the more energetic beam had already produced in its entrance region. This ensures that the lateral (caused by magnetic deflection) and longitudinal scanning (by energy variation) covers the tumour completely. In figure 9 (after Kraft1 ) the dose distribution for individual energy settings and the resulting total dose is sketched and compared with the damage that X-rays from a <sup>60</sup>Co-source would produce. An artist impression of the dose distribution for a lung and a brain tumour is given in figure 10. ## Treatment facilities Berkeley was the birthplace of therapy with hadrons. Since 1954 protons and later Helium-nuclei were used for treatment. Throughout the world treatment with protons is standard (Sweden, USA, Russia, Japan, Switzerland, England, Belgium, France, South Africa). In some places negative pions have been used in the past (USA, Canada, Switzerland). The most promising results have been obtained with heavy ions (Berkeley, USA; Chiba, Japan; and Darmstadt, Germany). In total $``$ 25000 patients have been treated from 1954 to 1999. ## Summary and Outlook The inverse ionisation dose profile of charged particles has been known for a long time, from nuclear and particle physics. The instrumentation originally developed for elementary particle physics experiments has made it possible to design and monitor particle beams with great precision which can then be used for tumour therapy. Heavy ions seem to be ideal projectiles for tumour treatment. They are suitable for well localized tumours. The availability of treatment facilities is increasing. Naturally such a facility requires an expensive and complex accelerator for the charged particles. For beam steering and control sophisticated particle detectors and interlock systems are necessary to ensure the safety of patients. ## Acknowledgements The author has benefitted a great deal from information provided by G. Kraft from GSI-Darmstadt and from discussions with him. I acknowledge also the help of Mrs. L. Hoppe and C. Haucke for the drawing of the figures, Mrs. A. Wied for typing the text, Mr. Ngac An Bang for giving the paper the final LaTeX-touch, and Mr. D. Robinson for a careful reading of the manuscript.
warning/0004/hep-th0004083.html	ar5iv	text	# Supersymmetric Fluid Mechanics ## I Introduction An isentropic fluid is described by a matter density field $`\rho `$ and a velocity field v. These satisfy the continuity equation, which involves the current $`\text{j}=\text{v}\rho `$, $$\dot{\rho }+\mathbf{}(\rho \text{v})=0$$ (1) and the force equation $$\dot{\text{v}}+\text{v}\mathbf{}\text{v}=\frac{1}{\rho }\mathbf{}P$$ (2) where $`P`$ is the pressure. (Over-dot denotes differentiation with respect to time.) We show that it is possible to supplement the $`(\rho ,\text{v})`$ bosonic/commuting variables with Grassmann/anticommuting variables $`\psi `$ such that the entire system exhibits a centrally extended supersymmetry. Moreover, when the bosonic system is irrotational, so that its vorticity vanishes, $`\omega _{ij}_iv^j_jv^i=0`$, and the velocity is the gradient of a velocity potential $`\text{v}=\mathbf{}\theta `$, the Grassmann variables give rise to nonvanishing vorticity and provide the Gaussian potentials in a Clebsch representation for the total velocity (see below). The specific system that we analyze devolves from the dynamics for a membrane (a 2-dimensional extended object), which propagates in $`(3+1)`$ dimensional space-time. The emergent fluid propagates in two spatial dimensions. When the membrane involves just bosonic variables, the fluid is irrotational . Our supersymmetric fluid is derived by a similar construction, starting from a supermembrane . In the remainder of this section, we review the action/Hamiltonian formulation for the system (1)–(2). In the next section, we directly present the supersymmetric model. Section III is devoted to a derivation of this supersymmetric fluid from a supermembrane, while concluding remarks comprise the last Section IV. For isentropic fluids, the pressure $`P`$ is a function only of the density, and the right side of (2) may also be written as $`\mathbf{}V^{}(\rho )`$, where $`V^{}`$ is the enthalpy, $`V^{\prime \prime }(\rho )=\frac{1}{\rho }P^{}(\rho )`$, and $`\sqrt{P^{}}`$ is the sound speed (dash denotes differentiation with respect to argument). Moreover, equations (1) and (2) can be obtained by (Poisson) bracketing with the Hamiltonian $$H=dr\left(\frac{1}{2}\rho v^2+V(\rho )\right)$$ (3) $`\dot{\rho }`$ $`=`$ $`\{H,\rho \}`$ (5) $`\dot{\text{v}}`$ $`=`$ $`\{H,\text{v}\}`$ (6) provided the nonvanishing brackets of the fundamental variables $`(\rho ,\text{v})`$ are taken to be $`\{v^i(\text{r}),\rho (\text{r}^{})\}`$ $`=`$ $`_i\delta (\text{r}\text{r}^{})`$ (8) $`\{v^i(\text{r}),v^j(\text{r}^{})\}`$ $`=`$ $`{\displaystyle \frac{\omega _{ij}(\text{r})}{\rho (\text{r})}}\delta (\text{r}\text{r}^{}).`$ (9) (The fields in the brackets are at equal times, hence the time argument is suppressed.) An equivalent, more transparent version of the algebra (I) is satisfied by the field momentum density, which also coincides with the current j. $$𝒫=\rho \text{v}$$ (10) As a consequence of (I) we have $`\{𝒫^i(\text{r}),\rho (\text{r}^{})\}`$ $`=`$ $`\rho (\text{r})_i\delta (\text{r}\text{r}^{})`$ (12) $`\{𝒫^i(\text{r}),𝒫^j(\text{r}^{})\}`$ $`=`$ $`𝒫^j(\text{r})_i\delta (\text{r}\text{r}^{})+𝒫^i(\text{r}^{})_j\delta (\text{r}\text{r}^{}).`$ (13) This is the familiar algebra of momentum densities. The Jacobi identity is satisfied by (I) and (I). The above holds in any dimension. One naturally asks whether there is a canonical 1-form that leads to the symplectic structure (I), (I); that is, one seeks a Lagrangian whose canonical variables can be used to derive (I) and (I) from canonical brackets. When the velocity is irrotational, the vorticity vanishes, v can be written as $`\mathbf{}\theta `$, and (I) is satisfied by postulating that $$\{\theta (\text{r}),\rho (\text{r}^{})\}=\delta (\text{r}\text{r}^{})$$ (14) that is, the velocity potential is conjugate to the density, so that the Lagrangian can be taken as $$L\|_{\mathrm{irrotational}}=dr\theta \dot{\rho }H$$ (15) where $`H`$ is given by (3) with $`\text{v}=\mathbf{}\theta `$. With nonvanishing vorticity, the canonical formulation is more indirect. One writes the velocity in a Clebsch decomposition, which in two and three spatial dimensions reads $$\text{v}=\mathbf{}\theta +\alpha \mathbf{}\beta .$$ (16) Then $$L=dr\rho (\dot{\theta }+\alpha \dot{\beta })H.$$ (17) Here $`\alpha `$ and $`\beta `$ are the “Gauss potentials”, and from (17) is seen that $`\{\theta ,\rho \}`$ as well as $`\{\beta ,\rho \alpha \}`$ are canonically conjugate. It then follows that v, given by (16), and $`\rho `$ satisfy (I).<sup></sup><sup></sup>*Some more observations on the Clebsch decomposition of the vector field v: In three dimensions, Eq. (16) involves the same number of functions on the left and right sides of the equality: three. Nevertheless the Gauss potentials are not uniquely determined by v. The following is the reason why a canonical formulation of (I) requires using the Clebsch decomposition (16). Although the algebra (I) is consistent in that the Jacobi identity is satisfied, it is degenerate in that the kinematic helicity $`h`$ $$h\frac{1}{2}d^3r\text{v}(\mathbf{}\times \text{v})=\frac{1}{2}d^3r\text{v}𝝎$$ ($`\omega ^i=\frac{1}{2}ϵ^{ijk}\omega _{jk}`$) has vanishing bracket with $`\rho `$ and v. (Note that $`h`$ is just the Abelian Chern-Simons term of v.) Consequently, a canonical formulation requires eliminating the kernel of the algebra, that is, neutralizing $`h`$. This is achieved by the Clebsch decomposition: $`\text{v}=\mathbf{}\theta +\alpha \mathbf{}\beta `$, $`𝝎=\mathbf{}\alpha \times \mathbf{}\beta `$, $`\text{v}𝝎=\mathbf{}\theta (\mathbf{}\alpha \times \mathbf{}\beta )=\mathbf{}(\theta \mathbf{}\alpha \times \mathbf{}\beta )`$. Thus in the Clebsch parameterization the helicity is given by a surface integral $`h=\frac{1}{2}d\text{S}\theta (\mathbf{}\alpha \times \mathbf{}\beta )`$ — it possesses no bulk contribution, and the obstruction to a canonical realization of (I) is removed. In two spatial dimensions, the Clebsch parameterization is redundant, involving three functions to express the two velocity components. Moreover, the kernel of (I) in two dimensions comprises an infinite number of quantities $$k_n=d^2r\rho \left(\frac{\omega }{\rho }\right)^n$$ for which the Clebsch parameterization offers no simplification. (Here $`\omega `$ is the two-dimensional vorticity $`\omega _{ij}=ϵ_{ij}\omega `$.) Nevertheless, a canonical formulation in two dimensions also uses Clebsch variables to obtain an even-dimensional phase space. ## II Supersymmetric Fluid Mechanics ### A The Model The bosonic fluid model in two spatial dimensions, which descends from a bosonic Nambu-Goto action, is supplemented by Grassmann variables $`\psi _a`$ that are Majorana spinors (real, two-component: $`\psi _a^{}=\psi _a`$, $`a=1,2`$). The Lagrange density reads $$=\rho (\dot{\theta }\frac{1}{2}\psi \dot{\psi })\frac{1}{2}\rho (\mathbf{}\theta \frac{1}{2}\psi \mathbf{}\psi )^2\frac{\lambda }{\rho }\frac{\sqrt{2\lambda }}{2}\psi 𝜶\mathbf{}\psi .$$ (18) Here $`\alpha ^i`$ are two ($`i=1,2`$), $`2\times 2`$, real symmetric Dirac “alpha” matrices; in terms of Pauli matrices we can take $`\alpha ^1=\sigma ^1`$, $`\alpha ^2=\sigma ^3`$. Note that the matrices satisfy the following relations, which are needed to verify subsequent formulas $`ϵ_{ab}\alpha _{bc}^i`$ $`=`$ $`ϵ^{ij}\alpha _{ac}^j`$ (19) $`\alpha _{ab}^i\alpha _{bc}^j`$ $`=`$ $`\delta ^{ij}\delta _{ac}ϵ^{ij}ϵ_{ac}`$ (20) $`\alpha _{ab}^i\alpha _{cd}^i`$ $`=`$ $`\delta _{ac}\delta _{bd}\delta _{ab}\delta _{cd}+\delta _{ad}\delta _{bc}`$ (21) $`ϵ_{ab}`$ is the $`2\times 2`$ antisymmetric matrix $`ϵi\sigma ^2`$. In (18) $`\lambda `$ is a coupling strength, taken positive. The density-dependent potential $`V(\rho )=\lambda /\rho `$ corresponds to a negative pressure $`P=2\lambda /\rho `$ and to sound velocity $`\sqrt{2\lambda }/\rho `$. These describe the “Chaplygin gas”. The Grassmann term enters with coupling $`\sqrt{2\lambda }`$, so chosen to ensure supersymmetry (see below). It is evident that the velocity should be defined as $$\text{v}=\mathbf{}\theta \frac{1}{2}\psi \mathbf{}\psi .$$ (22) The Grassmann variables directly give rise to a Clebsch formula for v, and provide the Gauss potentials. The two-dimensional vorticity reads $`\omega =ϵ^{ij}_iv^j=\frac{1}{2}ϵ^{ij}_i\psi _j\psi =\frac{1}{2}\mathbf{}\psi \times \mathbf{}\psi `$. The variables $`\{\theta ,\rho \}`$ remain a canonical pair, while the canonical 1-form in (18) indicates that the canonically independent Grassmann variables are $`\sqrt{\rho }\psi `$ so that the antibracket of the $`\psi `$’s is $$\{\psi _a(\text{r}),\psi _b(\text{r}^{})\}=\frac{\delta _{ab}}{\rho (\text{r})}\delta (\text{r}\text{r}^{}).$$ (23) One verifies that the algebra (I) or (10) is satisfied, and further, one has $`\{\theta (\text{r}),\psi (\text{r})\}`$ $`=`$ $`{\displaystyle \frac{1}{2\rho (\text{r})}}\psi (\text{r})\delta (\text{r}\text{r}^{})`$ (25) $`\{\text{v}(\text{r}),\psi (\text{r}^{})\}`$ $`=`$ $`{\displaystyle \frac{\mathbf{}\psi (\text{r})}{\rho (\text{r})}}\delta (\text{r}\text{r}^{})`$ (26) $`\{𝓟(\text{r}),\psi (\text{r}^{})\}`$ $`=`$ $`\mathbf{}\psi (\text{r})\delta (\text{r}\text{r}^{}).`$ (27) The equations of motion read $`\dot{\rho }+\mathbf{}(\rho \text{v})`$ $`=`$ $`0`$ (29) $`\dot{\theta }+\text{v}\mathbf{}\theta `$ $`=`$ $`\frac{1}{2}v^2+{\displaystyle \frac{\lambda }{\rho ^2}}+{\displaystyle \frac{\sqrt{2\lambda }}{2\rho }}\psi 𝜶\mathbf{}\psi `$ (30) $`\dot{\psi }+\text{v}\mathbf{}\psi `$ $`=`$ $`{\displaystyle \frac{\sqrt{2\lambda }}{\rho }}𝜶\mathbf{}\psi `$ (31) and together with (22) they imply $$\dot{\text{v}}+\text{v}\mathbf{}\text{v}=\mathbf{}\frac{\lambda }{\rho ^2}+\frac{\sqrt{2\lambda }}{\rho }\mathbf{}\psi 𝜶\mathbf{}\psi .$$ (32) All these equations may be obtained by bracketing with the Hamiltonian $$H=d^2r\left(\frac{1}{2}\rho v^2+\frac{\lambda }{\rho }+\frac{\sqrt{2\lambda }}{2}\psi 𝜶\mathbf{}\psi \right)$$ (33) when (I), (14) and (23) are used. We record the components of the energy-momentum tensor, and the continuity equations they satisfy. The energy density $`=T^{oo}`$, given by $$T^{oo}=\frac{1}{2}\rho v^2+\frac{\lambda }{\rho }+\frac{\sqrt{2\lambda }}{2}\psi 𝜶\mathbf{}\psi $$ (34) satisfies a continuity equation with the energy flux $`T^{oj}`$. $$\dot{T}^{oo}+_jT^{oj}=0$$ (36) $$T^{oj}=\frac{1}{2}\rho v^2v^j\frac{\lambda v^j}{\rho }+\frac{\sqrt{2\lambda }}{2}\psi \alpha ^j\text{v}\mathbf{}\psi \frac{\lambda }{\rho }\psi _j\psi +\frac{\lambda }{\rho }ϵ^{jk}\psi ϵ_k\psi $$ (37) This ensures that the total energy, that is, the Hamiltonian, is time-independent. Conservation of the total momentum $$\text{P}=d^2r𝓟$$ (38) follows from the continuity equation satisfied by the momentum density $`𝒫^i=T^{io}`$ $$\dot{T}^{io}+_jT^{ij}=0$$ (40) $$T^{ij}=\rho v^iv^j\delta ^{ij}\left(\frac{2\lambda }{\rho }+\frac{\sqrt{2\lambda }}{2}\psi 𝜶\mathbf{}\psi \right)+\frac{\sqrt{2\lambda }}{2}\psi \alpha ^j_i\psi $$ (41) but the momentum flux $`T^{ij}`$, that is, the stress tensor, is not symmetric in its spatial indices, owing to the presence of spin in the problem. However, rotational symmetry makes it possible to effect an “improvement”, which modifies the momentum density by a total derivative term, leaving the integrated total momentum unchanged (provided surface terms can be ignored) and rendering the stress tensor symmetric. The improved quantities are $$𝒫_I^i=T_I^{io}=\rho v^i+\frac{1}{8}ϵ^{ij}_j(\rho \psi ϵ\psi )$$ (42) $$\dot{T}_I^{io}+_jT_I^{ij}=0$$ (44) $`T_I^{ij}`$ $`=`$ $`\rho v^iv^j\delta ^{ij}\left({\displaystyle \frac{2\lambda }{\rho }}+{\displaystyle \frac{\sqrt{2\lambda }}{2}}\psi 𝜶\mathbf{}\psi \right)+{\displaystyle \frac{\sqrt{2\lambda }}{4}}\left(\psi \alpha ^i_j\psi +\psi \alpha ^j_i\psi \right)`$ (46) $`\frac{1}{8}_k\left[(ϵ^{ki}v^j+ϵ^{kj}v^i)\rho \psi ϵ\psi \right].`$ It immediately follows that the angular momentum $$M=d^2rϵ^{ij}r^i𝒫_I^j=d^2r\rho ϵ^{ij}r^iv^j+\frac{1}{4}dr\rho \psi ϵ\psi $$ (47) is conserved. The first term is clearly the orbital part (which still receives a Grassmann contribution through v), whereas the second, coming from the improvement, is the spin part. Indeed, since $`\frac{i}{2}ϵ=\frac{1}{2}\sigma ^2\mathrm{\Sigma }`$, we recognize this as the spin matrix in (2+1) dimensions. The extra term in the improved momentum density $`\frac{1}{8}ϵ^{ij}_j(\rho \psi ϵ\psi )`$ can then be readily interpreted as an additional localized momentum density, generated by the nonhomogeneity of the spin density. This is analogous to the magnetostatics formula giving the localized current density $`\text{j}_m`$ in a magnet in terms of its magnetization m: $`\text{j}_m=\mathbf{}\times \text{m}`$. All in all, we are describing a fluid with spin. Also the total number $$N=d^2r\rho $$ (48) is conserved by virtue of the continuity equation (29) satisfied by $`\rho `$. Finally, the theory is Galileo invariant, as is seen from the conservation of the Galileo boost, $$\text{B}=t\text{P}d^2r\text{r}\rho $$ (49) which follows from (29) and (38). The generators $`H,\text{P},M,\text{B}`$ and $`N`$ close on the (extended) Galileo group. \[The theory is not Lorentz invariant in $`(2+1)`$-dimensional space-time, hence the energy flux $`T^{oi}`$ does not coincide with the momentum density, improved or not.\] We observe that $`\rho `$ can be eliminated from (18) so that $``$ involves only $`\theta `$ and $`\psi `$. From (30) and (31) it follows that $$\rho =\left(\dot{\theta }\frac{1}{2}\psi \dot{\psi }+\frac{1}{2}v^2\right)^{{\scriptscriptstyle \frac{1}{2}}}.$$ (50) Substituting into (18) leaves $$=\sqrt{2\lambda }\left\{\sqrt{2\dot{\theta }\psi \dot{\psi }+(\mathbf{}\theta \frac{1}{2}\psi \mathbf{}\psi )^2}+\frac{1}{2}\psi 𝜶\mathbf{}\psi \right\}.$$ (51) Note that the coupling strength has disappeared from the dynamical equations, remaining only as a normalization factor for the Lagrangian. Consequently the above elimination of $`\rho `$ cannot be carried out in the free case, $`\lambda =0`$. ### B Supersymmetry The theory also possesses supersymmetry. This can be established, first of all, by verifying that the following two-component dynamics-dependent supercharges are time-independent Grassmann quantities. $$Q_a=d^2r\left[\rho \text{v}(𝜶_{ab}\psi _b)+\sqrt{2\lambda }\psi _a\right].$$ (52) Taking a time derivative and using the evolution equations (II A) establishes that $`\dot{Q}_a=0`$. Next, the transformation rule for the dynamical variables is found by considering the Grassmann charge contracted with a constant Grassmann parameter $`\eta ^a`$, giving a bosonic symmetry generator $`Q=\eta ^aQ_a`$. Using the canonical brackets one verifies the field transformation rules $`\delta \rho =\{Q,\rho \}`$ $`=`$ $`\mathbf{}\rho (\eta 𝜶\psi )`$ (54) $`\delta \theta =\{Q,\theta \}`$ $`=`$ $`\frac{1}{2}(\eta 𝜶\psi )\mathbf{}\theta \frac{1}{4}(\eta 𝜶\psi )\psi \mathbf{}\psi +{\displaystyle \frac{\sqrt{2\lambda }}{2\rho }}\eta \psi `$ (55) $`\delta \psi =\{Q,\psi \}`$ $`=`$ $`(\eta 𝜶\psi )\mathbf{}\psi \text{v}𝜶\eta {\displaystyle \frac{\sqrt{2\lambda }}{\rho }}\eta `$ (56) $`\delta \text{v}=\{Q,\text{v}\}`$ $`=`$ $`(\eta 𝜶\psi )\mathbf{}\text{v}+{\displaystyle \frac{\sqrt{2\lambda }}{\rho }}\eta \mathbf{}\psi .`$ (57) Supersymmetry is reestablished by determining the variation of the action $`dtL`$, consequent to the above field variations: the action is invariant. One then reconstructs the supercharges (52) by Noether’s theorem. Finally, upon computing the bracket of two supercharges, one finds $$\{\eta _1^aQ_a,\eta _2^bQ_b\}=2(\eta _1\eta _2)H$$ (58) which again confirms that the charges are time-independent: $$\{H,Q_a\}=0.$$ (59) Additionally a further, kinematical, supersymmetry can be identified. According to the equations of motion the following two supercharges are also time-independent: $$\stackrel{~}{Q}_a=d^2r\rho \psi _a.$$ (60) $`\stackrel{~}{Q}=\stackrel{~}{\eta }^a\stackrel{~}{Q}_a`$ effects a shift of the Grassmann field: $`\stackrel{~}{\delta }\rho =\{\stackrel{~}{Q},\rho \}`$ $`=`$ $`0`$ (62) $`\stackrel{~}{\delta }\theta =\{\stackrel{~}{Q},\theta \}`$ $`=`$ $`\frac{1}{2}(\stackrel{~}{\eta }\psi )`$ (63) $`\stackrel{~}{\delta }\psi =\{\stackrel{~}{Q},\psi \}`$ $`=`$ $`\stackrel{~}{\eta }`$ (64) $`\stackrel{~}{\delta }\text{v}=\{\stackrel{~}{Q},\text{v}\}`$ $`=`$ $`0.`$ (65) This transformation leaves the Lagrangian invariant, and Noether’s theorem reproduces (60). The algebra of these charges closes on the total number $`N`$. $$\{\stackrel{~}{\eta }_1^a\stackrel{~}{Q}_a,\stackrel{~}{\eta }_2^b\stackrel{~}{Q}_b\}=(\stackrel{~}{\eta }_1\stackrel{~}{\eta }_2)N$$ (66) while the algebra with the generators (52), closes on the total momentum, together with a central extension, proportional to volume of space $`\mathrm{\Omega }=d^2r`$ $$\{\stackrel{~}{\eta }^a\stackrel{~}{Q}_a,\eta ^bQ_b\}=(\stackrel{~}{\eta }𝜶\eta )\text{P}+\sqrt{2\lambda }(\stackrel{~}{\eta }ϵ\eta )\mathrm{\Omega }.$$ (67) The supercharges $`Q_a,\stackrel{~}{Q}_a`$, together with the Galileo generators ($`H`$, P, $`M`$, and B), with $`N`$ form a superextended Galileo algebra. The additional, nonvanishing brackets are $`\{M,Q_a\}`$ $`=`$ $`\frac{1}{2}ϵ^{ab}Q_b`$ (68) $`\{M,\stackrel{~}{Q}_a\}`$ $`=`$ $`\frac{1}{2}ϵ^{ab}\stackrel{~}{Q}_b`$ (69) $`\{\text{B},Q_a\}`$ $`=`$ $`𝜶_{ab}\stackrel{~}{Q}_b.`$ (70) ## III Membrane Connection The equations for a supersymmetric Chaplygin fluid devolve from the supermembrane Lagrangian, $`L_M`$. We shall give two different derivations of this result, which make use of two different parameterizations for the parameterization-invariant membrane action and give rise, respectively, to (18) and (51). We work in a light-cone gauge-fixed theory: The membrane in 4-dimensional space-time is described by coordinates $`x^\mu `$ $`(\mu =0,1,2,3)`$, which are decomposed into light-cone components $`x^\pm =\frac{1}{\sqrt{2}}(x^0\pm x^3)`$ and transverse components $`x^i`$ $`\{i=1,2\}`$. These depend on an evolution parameter $`\tau `$ and two space-like parameters $`\varphi ^r`$ $`\{r=1,2\}`$. Additionally there are two-component, real Grassmann spinors $`\psi `$, which also depend on $`\tau `$ and $`\varphi ^r`$. In the light-cone gauge, $`x^+`$ is identified with $`\tau `$, $`x^{}`$ is renamed $`\theta `$, and the supermembrane Lagrangian is $$L_M=d^2\varphi _M=d^2\varphi \{\sqrt{G}\frac{1}{2}ϵ^{rs}_r\psi 𝜶_s\psi \text{x}\}$$ (71) where $`G=detG_{\alpha \beta }`$; $`G_{\alpha \beta }`$ $`=`$ $`\left(\begin{array}{cc}G_{oo}& \hfill G_{os}\\ G_{ro}& \hfill g_{rs}\end{array}\right)`$ (74) $`=`$ $`\left(\begin{array}{cc}2_\tau \theta _\tau \text{x}_\tau \text{x}\psi _\tau \psi & \hfill u_s\\ u_r& \hfill g_{rs}\end{array}\right)`$ (77) $`G`$ $`=`$ $`g\mathrm{\Gamma }`$ (78) $`\mathrm{\Gamma }`$ $``$ $`2_\tau \theta _\tau \text{x}_\tau \text{x}\psi _\tau \psi +g^{rs}u_ru_s`$ (79) $`g_{rs}`$ $``$ $`_r\text{x}_s\text{x},g=detg_{rs}`$ (80) $`u_r`$ $``$ $`_r\theta \frac{1}{2}\psi _r\psi _\tau \text{x}_r\text{x}.`$ (81) Here $`_\tau `$ signifies differentiation with respect to the evolution parameter $`\tau `$, while $`_r,_s`$ differentiate with respect to the space-like parameters $`(\varphi ^r,\varphi ^s)`$, and $`g^{rs}`$, the inverse of $`g_{rs}`$, is used to move the $`(r,s)`$ indices. Note that the dimensionality of the transverse coordinates $`x^i`$ is the same as of the parameters $`\varphi ^r`$, namely two. ### A First Derivation To give our first derivation, we rewrite the Lagrangian in canonical, first-order form, with the help of canonical momenta defined by $`{\displaystyle \frac{_M}{_\tau \text{x}}}`$ $`=`$ $`\text{p}=\mathrm{\Pi }_\tau \text{x}\mathrm{\Pi }u^r_r\text{x}`$ (83) $`{\displaystyle \frac{_M}{_\tau \theta }}`$ $`=`$ $`\mathrm{\Pi }=\sqrt{g/\mathrm{\Gamma }}`$ (84) $`_M`$ $`=`$ $`\text{p}_\tau \text{x}+\mathrm{\Pi }_\tau \theta \frac{1}{2}\mathrm{\Pi }\psi _\tau \psi +{\displaystyle \frac{1}{2\mathrm{\Pi }}}(p^2+g)+\frac{1}{2}ϵ^{rs}_r\psi 𝜶_s\psi \text{x}`$ (86) $`+u^r\left(_r\text{x}\text{p}+\mathrm{\Pi }_r\theta \frac{1}{2}\mathrm{\Pi }\psi _r\psi \right).`$ In (86) $`u^r`$ serves as a Lagrange multiplier enforcing a subsidiary condition on the canonical variables. The equations that follow from (86) coincide with the Euler-Lagrange equations for (71)–(81). The theory still possesses an invariance against redefining the spatial parameters with a $`\tau `$-dependent function of the parameters. This freedom may be used to set $`u_r`$ to zero and fix $`\mathrm{\Pi }`$ at $`1`$. Next we introduce the hodographic transformation , whereby independent-dependent variables are interchanged, namely we view the $`\varphi ^r`$ to be functions of $`x^i`$. It then follows that the constraint on (86), which with $`\mathrm{\Pi }=1`$ reads $$_r\text{x}\text{p}_r\theta +\frac{1}{2}\psi _r\psi =0$$ (88) becomes $$_r\text{x}\left(\text{p}\mathbf{}\theta +\frac{1}{2}\psi \mathbf{}\psi \right)=0.$$ (89) Here p, $`\theta `$ and $`\psi `$ are viewed as functions of x, renamed r, with respect to which acts the gradient $`\mathbf{}`$. Also we rename p as v, which according to (89) is $$\text{v}=\mathbf{}\theta \frac{1}{2}\psi \mathbf{}\psi .$$ (90) From the chain rule, it follows that $$_\tau =_t+_\tau \text{x}\mathbf{}$$ (91) and according to (83) (at $`\mathrm{\Pi }=1`$, $`u^r=0`$) $`_\tau \text{x}=\text{p}=\text{v}`$. Finally, the measure transforms according to $`\mathrm{d}^2\varphi \mathrm{d}^2r\frac{1}{\sqrt{g}}`$. Thus the Lagrangian for (86) becomes, after setting $`u^r`$ to zero and $`\mathrm{\Pi }`$ to $`1`$, $$L_M=\frac{\mathrm{d}^2r}{\sqrt{g}}\left(v^2\dot{\theta }\text{v}\mathbf{}\theta +\frac{1}{2}\psi (\dot{\psi }+\text{v}\mathbf{}\psi )\frac{1}{2}(v^2+g)\frac{1}{2}ϵ^{rs}\psi \alpha ^i_j\psi _sx^j_rx^i\right).$$ (92) But $`ϵ^{rs}_sx^j_rx^i=ϵ^{ij}det_rx^i=ϵ^{ij}\sqrt{g}`$. After $`\sqrt{g}`$ is renamed $`\sqrt{2\lambda }/\rho `$, (92) finally reads $$L_M=\left(\frac{1}{\sqrt{2\lambda }}\right)d^2r\left(\rho (\dot{\theta }\frac{1}{2}\psi \dot{\psi })\frac{1}{2}\rho (\mathbf{}\theta \frac{1}{2}\psi \mathbf{}\psi )^2\frac{\lambda }{\rho }\frac{\sqrt{2\lambda }}{2}\psi 𝜶\times \mathbf{}\psi \right).$$ (93) Upon replacing $`\psi `$ by $`\frac{1}{\sqrt{2}}(1ϵ)\psi `$, this is seen to reproduce the Lagrange density (18), apart from an overall factor. ### B Second Derivation For our second derivation, we return to (71)–(81) and use the remaining reparameterization freedom to equate the two $`x^i`$ variables with the two $`\varphi ^r`$ variables, renaming both as $`r^i`$ . Also $`\tau `$ is renamed as $`t`$. In (71)–(81) $`g_{rs}=\delta _{rs}`$, and $`_\tau \text{x}=0`$, so that (81) becomes simply $`G=\mathrm{\Gamma }`$ $`=`$ $`2\dot{\theta }\psi \dot{\psi }+u^2`$ (94) u $`=`$ $`\mathbf{}\theta \frac{1}{2}\psi \mathbf{}\psi .`$ (95) Therefore the Nambu-Goto action (71) reads $$L_M=d^2r\left\{\sqrt{2\dot{\theta }\psi \dot{\psi }+\left(\mathbf{}\theta \frac{1}{2}\psi \mathbf{}\psi \right)^2}+\frac{1}{2}\psi 𝜶\times \mathbf{}\psi \right\}.$$ (96) Again a replacement of $`\psi `$ by $`\frac{1}{\sqrt{2}}(1ϵ)\psi `$ demonstrates that the integrand coincides with the Lagrange density in (51) (apart from a normalization factor). ### C Further Consequences of the Supermembrane Connection The supermembrane dynamics is Poincaré invariant in (3+1)-dimensional space-time. This invariance is hidden by the choice of light-cone parameterization: only the light-cone subgroup of the Poincaré group is left as a manifest invariance. This is just the $`(2+1)`$ Galileo group generated by $`H`$, P, $`M`$, B, and $`N`$. (The light-cone subgroup of the Poincaré group is isomorphic to the Galileo group in one lower dimension.) The Poincaré generators not included in the above list correspond to Lorentz transformations in the “$``$” direction. We expect therefore that these generators are “dynamical”, that is, hidden and unexpected conserved quantities of our supersymmetric Chaplygin gas, similar to the situation with the purely bosonic model . One verifies that the following quantities $`D`$ $`=`$ $`tH{\displaystyle d^2r\rho \theta }`$ (97) G $`=`$ $`{\displaystyle d^2r(\text{r}\theta 𝓟_I\frac{1}{8}\psi 𝜶𝜶𝓟_I\psi )}`$ (98) $`=`$ $`{\displaystyle d^2r(\text{r}\theta 𝓟\frac{1}{4}\psi 𝜶𝜶𝓟\psi )}`$ (99) are time-independent by virtue of the equations of motion (II A), and they supplement the Galileo generators to form the full $`(3+1)`$ Poincaré algebra, which becomes the super-Poincaré algebra once the supersymmetry is taken into account. ## IV Conclusion We have shown how fluid dynamics can be extended to include Grassmann variables, which also enter in a supersymmetry-preserving interaction. Since our construction is based on a supermembrane in (3+1)-dimensional space-time, the fluid model is necessarily a planar Chaplygin gas. It remains to be shown how this construction could be extended to arbitrary dimensions and to different interactions. Note that Grassmann Gauss potentials can be used even in the absence of supersymmetry. For example, our theory (18), with the last term omitted, posseses a conventional, bosonic Hamiltonian without supersymmetry, while the Grassmann variables are hidden in v and occur only in the canonical 1-form. In a related investigation, conventional fluid mechanics is generalized, so that it possesses a non-Abelian gauge symmetry . #### Note Added: J. Hoppe has informed us that some of the above results were obtained by him in unpublished research: Karlsruhe preprints KA-THEP-6-93 and KA-THEP-9-93 (hep-th/9311059).
warning/0004/astro-ph0004312.html	ar5iv	text	# Circumbinary disks around T Tauri stars: HST/NICMOS near-infrared images and polarimetric maps ## 1. Introduction The presence of circumstellar disks around T Tauri stars has been suspected for a long time, but it is only very recently that these disks were directly detected, using high-angular millimetre imaging. These images revealed extended gas structures which appeared to be in Keplerian rotation around the central object. Among the rare detections so far, two disks were found around binary T Tauri stars: GG Tau (Dutrey, Guilloteau & Simon 1994) and UY Aur (Duvert et al. 1998). The separations of the binaries are 0$`.^{\prime \prime }`$25 and 0$`.^{\prime \prime }`$89 respectively, which correspond to projected physical separations of 35 and 125 AU at the distance of the Taurus star-forming region (140 pc). In both cases, light scattered off the surface of the disks has been detected afterwards with adaptive optics imaging at near-infrared wavelengthes. Roddier et al. (1996) found that the GG Tau ring has a clumpy appearance and that several radial spokes of material extend from the ring onto the central stars. The ring is brighter in its northern part, but is detected in all directions. They interpret this brightness difference as being due to the scattering geometry. The UY Aur case is very different, as Close et al. (1998) only detected the disk on one side of the binary. Furthermore, they found evidences that a “spiral arm” splits from the main disk and gets closer to the star. Deconvolution processes were applied in both studies to retrieve the highest spatial resolution allowed by adaptive optics devices, and this may lead to some artifacts in the final images. More recently, the first visible wavelength images of UY Aur were obtained by Ménard et al. (1999) at 600 and 800 nm with HST/WFPC2. The PSF-subtracted images revealed a more complicated structure that was found by Close et al. (1998): a large “clump” appears to be independent from the disk itself. If true, this implies that the inclination of the system to the line-of-sight is larger than was first thought (about 60 instead of about 40). To improve our knowledge of these two circumbinary disks, we have performed new observations at 1$`\mu `$m and 2$`\mu `$m of these systems with HST/NICMOS. We used the polarimetric modes, and we obtained both intensity and polarization maps, which do not need to be deconvolved. The GG Tau polarization maps are the first ever obtained of this system, while Potter et al. (1998) already presented a deconvolved J-band polarization map of UY Aur which revealed a nice centrosymetric pattern. Polarization maps are powerfull tools to investigate the dust grain properties and the geometry and structure of the disks. In section 2, we summarize our observations and data processing steps, and the maps of both systems are presented and commented in section 3. Section 4 describes some implications of our results on the properties of these disks. ## 2. Observations and data processing The 1$`\mu `$m and 2$`\mu `$m images were obtained with Camera 1 and Camera 2 respectively, providing pixel scales of 0$`.^{\prime \prime }`$043 and 0$`.^{\prime \prime }`$075. Both binaries were observed through the three polarizers at each wavelength, during three 96 seconds exposures for each filter. The regular NICMOS data reduction pipeline prooved to be unsatisfying, and we had to re-reduced all data, with specific care to the so-called “pedestal effect”, to obtain final images where the sky level is flat all over the detector. To allow clear detections of the disks, it is mandatory to remove the bright stellar point spread funtions (PSFs). We first tried Tinytim PSFs, but it appeared that their match with the real ones is quite poor, so we turned to a “natural” star, i.e. a bright single star observed through the same filters. The diffraction spikes subtraction, though unperfect, is quite good, and the optical ghosts induced by some polarizers are naturally removed. Some residuals in the core of the PSFs, however, are still large, and nothing can be securely detected in the inner 0$`.^{\prime \prime }`$5 at 1$`\mu `$m. At 2$`\mu `$m, some fringing can be seen at separations as large as 3$`.^{\prime \prime }`$5. No deconvolution process was applied to our images, which allows an easier interpretation. ## 3. Results ### 3.1. GG Tau The new 1$`\mu `$m image of the GG Tau ring is presented in Fig. 1. Its overall geometry is in good agreement with Roddier et al. (1996)’s images, though with a higher signal-to-noise ratio. However, there are some noticeable features. First, the ring does not appear clumpy in our image. This property was likely an artifact introduced by the deconvolution process applied to the adaptive optics images. Fitting an ellipse onto the ring, we find a semi-major axis, a position angle and an inclination in excellent agreement with the millimetre results of Guilloteau et al. (1999). It is noticeable, however, that this ellipse is not centered on the center of mass of the binary. Our image does not allow us to confirm the existence of the spokes of material discovered by Roddier et al. (1996), because of the large PSF subtraction residuals inside the ring. Finally, a significant east-west asymetry in the northern part of the ring is seen in our intensity map. The polarization vectors are strikingly well organized in a centrosymetric pattern, which is symetric about the semi-minor axis of the ring. The brightest part of the ring, which is the closest to the observer, displays a lower polarization level than the faintest side, typically 20% as opposed to 50–60%. At 2$`\mu `$m, the disk is too close to the stars, and the large subtraction residuals prevent us from obtaining a clear image of the ring. However, we calculated the polarization map at this wavelength and, though the image is strongly dominated by the unpolarized stellar fluxes, a centrosymetric pattern is found in the polarization vectors, with a typical level of 5–10%, indicating that the intrinsic polarization level of the light scattered by the ring is high. ### 3.2. UY Aur The morphology of the UY Aur circumbinary disk in our new 1$`\mu `$m image is in good agreement with Ménard et al. (1999)’s optical image, though the former suffer from a poor signal-to-noise ratio. As an be seen in Fig. 2, the disk appears as an unresolved arc to the Southwest of the binary at both wavelengthes, while a bright clump to the southeast appears to be unrelated to this structure. Noticeably, the bright arc seems to widen to the West of the binary in our image. This can be interpreted as the arc breaking into two separate arcs: the disk itself remaining at least 2$`.^{\prime \prime }`$5 away from the stars, and an inner arc, getting closer to the stars. The latter would correspond to the “spiral arm” described by Close et al. (1998). A second feature in our map is an inner arc which is much brighter than the disk itslef. It lies about 1<sup>′′</sup> to the Southwest of the secondary. This may be a PSF artifact, but it can be seen using both a natural or a Tinytim-built PSF. Furthermore, its coincidence with a similar arc detected in the WFPC2 images is suggestive (see Fig. 2). At 2$`\mu `$m, the only area which is clearly separated from the PSF residuals is found to the southeast of the binary. It can be traced up to 5<sup>′′</sup> away from the stars. The back side of the disk remains undetected at both wavelengthes, which provides strong constraints on the dust grain properties. We note however, that the small arc seen to the Northeast of the primary at 1$`\mu `$m can be seen in all indivual images. The polarization pattern at 2$`\mu `$m is well organized in its southeastern part, with all vectors well aligned, in a fashion consistent with centrosymetric. The typical polarization level is about 40%. At 1$`\mu `$m, however, the picture is much different: though the vectors are basically aligned in the southwestern part of the disk with a typical level of 20%, again in a more or less centrosymetrical pattern, they are quite randomly oriented in the southeastern clump. Whether this is due to our reduction pipeline, to the low signal-to-noise ratio or reflects the intrinsic pattern is unclear. It must be noted that Potter et al. (1998) found a very different behaviour. Their data reduction process, however, included a deconvolution step, which impact on the polarization is unknwon. ## 4. Implications and open questions ### 4.1. GG Tau As already pointed out by Guilloteau et al. (1999), the shift between the apparent center of the ring and the center of mass of the binary is naturally explained by a thick ring geometry. This is related to the fact that, in the Mie theory, forwards scattering is strongly favoured. Hence, most of the light scattered towards the observer comes from the upper part of the disk’s inner edge, whose projection onto the sky is not symetric about the physical center of the ring. The quantitative model proposed by Guilloteau et al. (1999) is a ring with a half-thickness of 60 AU at its inner radius (180 AU) and a very sharp edge. The observed location of the ring appears to be in excellent agreement with the prediction from this model. Roddier et al. (1996) suggested that the thickness-to-radius ratio was about one tenth at the inner edge of the ring, which seems imcompatible with our results. The origin of the east-west asymetry is unclear. It may be due to the presence of two illuminating stars, or be related to the slight asymetry found in the millimetre wavelength image, which itself may reveal internal structures differences. ### 4.2. UY Aur Fitting ellipses on the large southwest unresolved arc both at 1$`\mu `$m and 600 nm yields very similar figures and, especially, an inclination to the line-of-sight of about 60 in both cases. This is much larger than the 42 estimated by Close et al. (1998). The reason for that discrepancy is that they assumed that the southeast clump belongs to the disk, which now seems unlikely. Duvert et al. (1998) pointed out that a larger inclination (60–70) is in better agreement with the millimetre observations. If the inner arc close to the secondary star proves to be real, it will represent a challenge for theoretical studies, as most of them predict that structures at such a location should be highly unstable due to the interaction with the binary system. ### 4.3. Monte Carlo modelling The polarization levels, as well as the front-to-back side flux ratios, are tightly linked to the disk geometry and to the dust grain properties. Our next step in this study is to run Monte Carlo simulations to investigate these properties. For instance, we may determine whether the dust grain size distribution in the disks is compatible with that of interstellar grains. In principle, we will also constrain the amount of “flaring” in both disks, as well as their geometrical height and optical depth. ## References Close, L., Dutrey, A., Roddier, F., Guilloteau, S., Roddier, C., Northcott, M., Ménard, F., Duvert, G., Graves, J., Potter, D. 1998, ApJ, 499, 833 Dutrey, A., Guilloteau, S., Simon, M. 1994, A&A, 286, 149 Duvert, G., Dutrey, A., Guilloteau, S., Ménard, F., Schuster, K., Prato, L., Simon, M. 1998, A&A, 332, 867 Guilloteau, S., Dutrey, A., Simon, M. 1999, A&A, 348, 570 Ménard, F., Stapelfeldt, K., Krist, J., Duvert, G., Padgett, D., Burrows, C. 1999, BAAS, 194, 6811 Potter, D., Close, L., Roddier, F., Roddier, C., Graves, J., Northcott, M. 1998, in ESO Conf proceedings n56, Astronomy with Adaptive Optics, ed. D. Bonaccini, 353 Roddier, C., Roddier, F., Northcott, M., Graves, J., Jim, K. 1996, ApJ, 463, 326
warning/0004/hep-ph0004162.html	ar5iv	text	# Neurino propagation in matter using the wave packet approach ## 1 Introduction Recent experimental discoveries seem to suggest that neutrino oscillations really exist. Since the observation of neutrino oscillation may provide valuable information on the basic properties of neutrinos, e.g. masses and mixing angles, it is important to know the underlying physics also on the conceptual level. It was pointed out e.g. in Refs. and that the standard quantum mechanical treatment of neutrino oscillations using plane waves is not completely satisfactory for many reasons. The wave packet approach , provides a more physical picture which is particularly adapted for describing phenomena localized in space and time. This formalism also elegantly accounts for the loss of coherence by the separation of wave packets. However, some authors (e.g. ) have been skeptical about the use of wave packets, and Refs. conclude that the concept of wave packet is unnecessary for all the relevant physical cases. A bunch of other methods has also been discussed (to name a few). In this paper we will consider neutrino oscillations and other phenomena in matter. We have decided to use wave packets because the calculations of the kind presented here have never been carried out before. We will first focus on the equation of motion for neutrinos propagating in matter of variable density. Generally the equation is modified in matter due to coherent forward scatterings , and an exact solution can be found only for few special cases. Here we address one special case where the density of matter changes slowly enough, so that the situation is said to be adiabatic. Now the eigensolutions for the equation of motion are found trivially, for arbitrary number of relativistic neutrinos. To describe neutrino oscillations in matter, we apply the method of Ref. in connection with these solutions. We can express the results formally using effective oscillation and coherence lengths which are not local quantities anymore. Generally neutrinos may propagate nonadiabatically, and solving the complete equation of motion even for two neutrino flavors is usually far from trivial (see e.g. for two specific density profiles). In this paper we show that the solution for “arbitrary” density profile can be constructed by using infinite integral series. Some supplementary calculations which may be of formal interest are enclosed in Appendix B. Nonadiabaticity is also related to the fact that so-called level crossings (or hoppings) between matter “eigenstates” take place ( and references therein). Our calculations show that the level crossing probabilities are modified when neutrinos are described by wave packets. ## 2 Neutrino propagation in matter in the adiabatic limit ### 2.1 Equation of motion Barring the details of the production and detection mechanisms of neutrinos, we we will focus on the propagation. Traditionally the propagation of neutrinos is modeled by solving the relativistic Schrödinger equation with the effective Hamiltonian $$\widehat{H}=\sqrt{\widehat{p}^2+m^2}+V(\widehat{x},t)\widehat{p}+\frac{1}{2}m^2\widehat{p}^1+V(\widehat{x},t),$$ (1) where $`\widehat{p}`$ is the momentum operator, $`m`$ is the neutrino mass matrix and $`V(x,t)`$ is a semiclassical potential due to the presence of medium. In this work we assume that $`V(x)`$ does not depend on time, and that it is a matrix diagonal in the weak interaction basis or flavor basis. This is well justified when considering neutrinos from the Sun, from supernovae, or on Earth. On the other hand, this assumption excludes neutrinos in the early universe that must be treated otherwise. Instead of decomposing the wave function in momentum space, we prefer to consider the Fourier transform in energy space, $$\psi (x,t)=\frac{1}{\sqrt{2\pi }}𝑑Ee^{iEt}\psi (x,E).$$ (2) This leads to a more physical picture when the Hamilton operator is a function of $`x`$. Particularly, now the energy $`E`$ is a constant, while the momentum $`p`$ is a function of $`x`$. Throughout this work we will assume the propagation to be a one-dimensional phenomenon. The resulting time independent Schrödinger equation can be written in the relativistic limit as $$i_x\psi (x,E)=\left(E+\frac{m^2}{2E}+V(x)\right)\psi (x,E)+O\left(\frac{m^4+m^2EV}{E^3}\right)\psi (x,E),$$ (3) where we have also assumed that $`EV`$. The right-hand side can be locally diagonalized. We call the respective eigenvectors as local matter eigenstates, and the respective eigenvalues can be written as $$p_a(x,E)E\frac{\mu _a^2(x,E)}{2E},$$ (4) where we introduced the effective mass $`\mu _a(x,E)`$ that simplifies the notation. In this paper the Greek indices refer to the flavor basis and the Latin indices to the matter basis. In the adiabatic limit, or for negligible mixing, Eq. (3) can be written as $$i_x\psi _a(x,E)\left(E+\frac{\mu _a^2(x,E)}{2E}\right)\psi _a(x,E)=p_a(x,E)\psi _a(x,E),$$ (5) Outside the adiabatic limit the complete equation of motion includes also nondiagonal terms (compare to Eq. (11) of Ref. ). From now on we follow the treatment presented in Ref. . Eq. (5) can be solved easily, and one has $$\psi _a(x,t)=\frac{1}{\sqrt{2\pi }}𝑑E\mathrm{exp}\left[i_0^x𝑑x^{}p_a(x^{},E)iEt\right]\psi _a(0,E),$$ (6) where Eq. (2) was used. The initial value of $`x`$ is put to zero since we consider neutrinos produced at the origin. ### 2.2 Wave packet solutions As initial condition, we assume that the neutrino wave function is a linear combination of matter eigenstates, $$\psi (0,E)=\underset{a}{}C_a(E)\psi _a(0,E).$$ (7) In the relativistic limit the coefficients $`C_a`$ match the respective row in the mixing matrix, $`C_a(x,E)=U_{\alpha a}^{}(x,E)`$, but for less relativistic neutrinos correction factors related to energy conservation must be taken into account . We assume further that each component of the wave function $`\psi _a(0,E)`$ has a Gaussian form $$\psi _a(0,E)=(2\pi \sigma _{EP}^2)^{1/4}\mathrm{exp}\left[\frac{(EE_a)^2}{4\sigma _{EP}^2}\right],$$ (8) where $`E_a`$ is the average energy of the corresponding matter eigenstate and $`\sigma _{EP}`$ is the energy width related to the production process fulfilling the uncertainty relation $`\sigma _{EP}\sigma _{tP}=1/2`$. In the adiabatic limit we obtain for the partial wave packet the solution $$\psi _a(x,t)=C_1𝑑E\mathrm{exp}\left[i_0^x𝑑x^{}p_a(x^{},E)iEt\frac{(EE_a)^2}{4\sigma _{EP}^2}\right],$$ (9) where $`C_1`$ is a numerical factor. We assume that the wave packet is sufficiently narrow in energy space, $`\sigma _{EP}E_a`$. This helps us to integrate the above integral and we also avoid considering the negative energies. This assumption is justified e.g. for solar neutrinos . One can now expand the momentum as $$p_a(x,E)p_a(x,E_a)+\frac{p_a}{E_a}(EE_a)=p_a(x,E_a)+\frac{EE_a}{v_a(x,E_a)},$$ (10) where $`v_a(x,E)`$ is the group velocity of each wave packet. Note that within this framework the group velocity depends on $`x`$ via the potential $`V(x)`$, unlike in some alternative works. Writing $`\sigma _{EP}v_a(0)\sigma _{pP}=v_a(0)/(2\sigma _{xP})`$, one has after a simple integration $$\psi _a(x,t)=C_2\mathrm{exp}\left[i_0^x𝑑x^{}p_a(x^{},E_a)iE_at\frac{v_a^2(0)}{4\sigma _{xP}^2}\left(_0^x\frac{dx^{}}{v_a(x^{})}t\right)^2\right],$$ (11) where $`v_a(0)`$ is the group velocity at the origin and $`\sigma _{pP}(\sigma _{xP})`$ is the momentum (spatial) width related to the production process. Eq. (11) implies that the wave function of a flavor neutrino $`\nu _\alpha `$ is given by $$\|\nu _\alpha (x,t)=C_2\underset{a}{}U_{\alpha a}^{}(0)\mathrm{exp}\left[i_0^x𝑑x^{}p_a(x^{},E_a)iE_at\frac{v_a^2(0)}{4\sigma _{xP}^2}\left(_0^x\frac{dx^{}}{v_a(x^{})}t\right)^2\right]\|\nu _a,$$ (12) where $`U(x)`$ represents the effective mixing matrix and the states $`\|\nu _a`$ are orthonormal. To get out relevant physics we need to know the size of the wave packet. This is very non-trivial and lots of rather confusing estimates have been presented. Here we assume that the width of the neutrino wave is mostly related to the spatial details of the production process. In principle it also depends on temporal properties, like the stability of the state producing the neutrino, but this is relevant only for very short-lived particles. It has been estimated ( and relevant references therein) that $`\sigma _{xP}`$ is of the order $`10^9\text{m}`$ for solar neutrinos, $`\sigma _{xP}10^{11}\text{m}`$ at the neutrino sphere of a supernova, and $`\sigma _{xP}10^6\text{m}`$ for reactor neutrinos, for example. ### 2.3 Observation of neutrino oscillation The quantum state of the neutrino is measured by an appropriate reaction. Assuming the respective process being $`\nu _\alpha +X`$ (something visible), we can relate the quantum mechanical uncertainty of the detection process to the quantum mechanical state of the particle $`X`$ before the collision. Here we assume that the relevant wave functions are Gaussians, centered at a distance $`L`$ from the origin (source), with a spatial width $`\sigma _{xD}`$. Hence the detection can be described by $$\|\nu _\beta (xL)=C_2^{}\underset{a}{}U_{\beta a}^{}(L)\mathrm{exp}\left[i_L^x𝑑x^{}p_a(x^{},E_a)\frac{v_a^2(L)}{4\sigma _{xD}^2}\left(_L^x\frac{dx^{}}{v_a(x^{})}\right)^2\right]\|\nu _a,$$ (13) where $`v_a(L)`$ is the group velocity at $`L`$. Note that this is independent of time . We emphasize that the spatial uncertainty, $`\sigma _{xD}`$, arises for purely quantum-mechanical reasons. One may be apt to believe that it is at most of the order of atomic distances, i.e. $`\sigma _{xD}10^9\text{m}10^{10}\text{m}`$ or less. In reality there are also other instrumental uncertainties related to the resolution or dimension of the detector. These can be accounted for by taking an average over the relevant length scale. The amplitude of the process $`\nu _\alpha \nu _\beta `$ is now simply $`A_{\alpha \beta }(L,T)`$ $`=`$ $`{\displaystyle 𝑑x\nu _\beta (xL)\|\nu _\alpha (x,T)}`$ (14) $`=`$ $`C_2^{\prime \prime }{\displaystyle \underset{a}{}}U_{\alpha a}^{}(0)U_{\beta a}(L){\displaystyle }dx\mathrm{exp}[i{\displaystyle _0^L}dx^{}p_a(x^{},E_a)iE_aT`$ $`{\displaystyle \frac{v_a^2(0)}{4\sigma _{xP}^2}}({\displaystyle _0^x}{\displaystyle \frac{dx^{}}{v_a(x^{})}}T)^2{\displaystyle \frac{v_a^2(L)}{4\sigma _{xD}^2}}\left({\displaystyle _L^x}{\displaystyle \frac{dx^{}}{v_a(x^{})}}\right)^2].`$ The first two terms in the exponential are independent of $`x`$, but the integration of the two other terms is far from obvious. Notice that the analytical form of $`v_a(x)`$ is unknown in general, and that we have not even chosen any specific density profile so far. It turns out, however, that the integral can be evaluated with a saddle point method to a sufficiently good approximation. After a straightforward, but rather lengthy calculation, presented in Appendix A, we obtain the amplitude (with Eqs. (A6) and (A10)) $$A_{\alpha \beta }(L,T)=C_3\underset{a}{}U_{\alpha a}^{}(0)U_{\beta a}(L)\mathrm{exp}\left[i_0^L𝑑xp_a(x,E_a)iE_aT\frac{1}{4}\frac{\left(T_0^L\frac{dx}{v_a(x)}\right)^2}{\left(\frac{\sigma _{xP}}{v_a(0)}\right)^2+\left(\frac{\sigma _{xD}}{v_a(L)}\right)^2}\right].$$ (15) The respective probability of the process $`\nu _\alpha \nu _\beta `$ is then $`P_{\alpha \beta }(L,T)`$ $`=`$ $`\|A_{\alpha \beta }(L,T)\|^2`$ (16) $``$ $`{\displaystyle \underset{a,b}{}}U_{\alpha a}^{}(0)U_{\beta a}(L)U_{\alpha b}(0)U_{\beta b}^{}(L)e^{G(L,T)},`$ where $`G(L,T)`$ $`=`$ $`i{\displaystyle _0^L}𝑑x\left(p_a(x,E_a)p_b(x,E_b)\right)i(E_aE_b)T`$ (17) $`{\displaystyle \frac{1}{4\chi _a}}\left(T{\displaystyle _0^L}{\displaystyle \frac{dx}{v_a(x)}}\right)^2{\displaystyle \frac{1}{4\chi _b}}\left(T{\displaystyle _0^L}{\displaystyle \frac{dx}{v_b(x)}}\right)^2`$ with $$\chi _a\left(\frac{\sigma _{xP}}{v_a(0)}\right)^2+\left(\frac{\sigma _{xD}}{v_a(L)}\right)^2.$$ (18) We still have to perform an integration over time in $`P_{\alpha \beta }(L,T)`$ since we are mainly interested in $`P_{\alpha \beta }(L)`$ . A straightforward Gaussian integration of the exponential of Eq. (16) yields $`{\displaystyle 𝑑Te^{G(L,T)}}`$ $``$ $`\mathrm{exp}\left\{i{\displaystyle _0^L}𝑑x(p_a(x,E_a)p_b(x,E_b))\right\}`$ (19) $`\mathrm{exp}\left\{{\displaystyle \frac{1}{4\chi _a}}\left({\displaystyle _0^L}{\displaystyle \frac{dx}{v_a(x)}}\right)^2{\displaystyle \frac{1}{4\chi _b}}\left({\displaystyle _0^L}{\displaystyle \frac{dx}{v_b(x)}}\right)^2\right\}`$ $`\times \mathrm{exp}\{{\displaystyle \frac{\chi _a\chi _b}{\chi _a+\chi _b}}[(E_aE_b)^2+{\displaystyle \frac{1}{4}}({\displaystyle \frac{1}{\chi _a}}{\displaystyle _0^L}{\displaystyle \frac{dx}{v_a(x)}}+{\displaystyle \frac{1}{\chi _b}}{\displaystyle _0^L}{\displaystyle \frac{dx}{v_b(x)}})^2`$ $`i(E_aE_b)({\displaystyle \frac{1}{\chi _a}}{\displaystyle _0^L}{\displaystyle \frac{dx}{v_a(x)}}+{\displaystyle \frac{1}{\chi _b}}{\displaystyle _0^L}{\displaystyle \frac{dx}{v_b(x)}})]\}`$ without the prefactor, which is a constant in the relativistic limit. A detailed analysis renders Eq. (19) to $`{\displaystyle 𝑑Te^{G(L,T)}}`$ $``$ $`\mathrm{exp}\{i{\displaystyle _0^L}dx(p_a(x,E_a)p_b(x,E_b))i(E_aE_b)L`$ (20) $`{\displaystyle \frac{1}{8\sigma _x^2}}\left[{\displaystyle _0^L}dx(v_a(x)v_b(x))\right]^2{\displaystyle \frac{(E_aE_b)^2}{8\sigma _p^2}}\},`$ where $`\sigma _x^2\sigma _{xP}^2+\sigma _{xD}^2,\sigma _p1/(2\sigma _x)`$ and the relativistic limit is once again considered. In order to simplify Eq. (20), we approximate (compare to Ref. ) $$E_aE_0+\xi \frac{\mu _a^2(0)}{2E_0},$$ (21) where $`\mu _a(0)`$ is the effective mass at the origin, $`E_0`$ is the central energy in the limit of zero neutrino masses and $`\xi `$ is a parameter of order unity, related to the energy-momentum conservation of the production process. The corresponding momentum is given by $$p_a(x,E_a)E_a\frac{\mu _a^2(x)}{2E_a}E_0+\xi \frac{\mu _a^2(0)}{2E_0}\frac{\mu _a^2(x)}{2E_0}.$$ (22) Notice that $`E_a`$ is a constant, but $`p_a(x,E_a)`$ depends on $`x`$, thus this approximation does not contradict the chosen convention. One may also point out that the initial dependence $`\mu _a=\mu _a(x,E_a)`$ has reduced to $`\mu _a=\mu _a(x,E_0)`$. Using Eqs. (21) and (22), Eq. (20) gives $$𝑑Te^{G(L,T)}\mathrm{exp}\left[\frac{i}{2E_0}_0^L𝑑x\mathrm{\Delta }\mu _{ab}^2(x)\frac{1}{8\sigma _x^2}\left(_0^L𝑑x\mathrm{\Delta }v_{ab}(x)\right)^2\frac{(E_aE_b)^2}{8\sigma _p^2}\right],$$ (23) where $`\mathrm{\Delta }\mu _{ab}^2(x)\mu _a^2(x)\mu _b^2(x),\mathrm{\Delta }v_{ab}(x)v_a(x)v_b(x)`$ and the last term has not been altered on purpose. Combining the above results we obtain a compact formula for the probability to observe a neutrino in a state $`\nu _\beta `$ at a distance $`L`$ (with the correct normalization $`_\beta P_{\alpha \beta }(L)=1`$), $`P_{\alpha \beta }(L)`$ $`=`$ $`{\displaystyle \underset{a,b}{}}U_{\alpha a}^{}(0)U_{\beta a}(L)U_{\alpha b}(0)U_{\beta b}^{}(L)`$ (24) $`\times \mathrm{exp}\left[2\pi i{\displaystyle \frac{L}{L_{ab}^{osc}(L)}}\left({\displaystyle \frac{L}{L_{ab}^{coh}(L)}}\right)^2{\displaystyle \frac{(E_aE_b)^2}{8\sigma _p^2}}\right],`$ where the effective oscillation and coherence lengths are defined by $$L_{ab}^{osc}(L)\frac{4\pi E_0L}{_0^L𝑑x\mathrm{\Delta }\mu _{ab}^2(x)},L_{ab}^{coh}(L)\frac{2\sqrt{2}\sigma _xL}{\left\|_0^L𝑑x\mathrm{\Delta }v_{ab}(x)\right\|}.$$ (25) The problem has reduced to computing the integrals over the effective mass and group velocity differences along the neutrino path. Let us comment on the following issues in connection with Eqs. (24) and (25): 1. The fact that the effective mixing matrices depend on the production and detection location is one part of the well-known MSW effect . 2. Compared to the usual oscillation and coherence lengths, $`L^{osc}=4\pi E_0/\mathrm{\Delta }_0,L^{coh}=2\sqrt{2}\sigma _x/\|\mathrm{\Delta }v_0\|`$ (where $`\mathrm{\Delta }_0(\mathrm{\Delta }v_0)`$ is the mass squared (velocity) difference in vacuum), we see that the corresponding effective lengths take into account the changes of $`\mathrm{\Delta }\mu _{ab}^2(x)`$ and $`\mathrm{\Delta }v_{ab}(x)`$ over the whole path of propagation. In this sense the oscillation and coherence lengths are not local quantities anymore. 3. The physical coherence length, corresponding to a length scale where the oscillation ceases, can be obtained by solving the equation $`L=L_{ab}^{coh}(L)`$. The physical oscillation length for variable density lacks a clear definition. 4. The last term of the exponential in Eq. (24) is related to the energy conservation within the uncertainty $`\sigma _p`$. Its physical meaning is easy to understand: if e.g. $`\|E_1E_2\|\sigma _p`$, only one of the states $`\nu _1`$ or $`\nu _2`$ is “allowed”, i.e. there is no oscillation. 5. The treatment presented here is not limited to any specific density profile if only Eq. (5) is valid (i.e. the adiabaticity is in effect, the matter density is not too high etc.). The discussion of Ref. is in many respects applicable also here. We finally point out that the calculation of $`_0^L𝑑x\mathrm{\Delta }v_{ab}(x)`$ is trivial in the relativistic limit: the group velocity is given by definition by $$v_a(x)=\frac{E_a}{p_a}1\frac{\mu _a^2(x)}{2p_a^2(x)}+\frac{1}{2p_a(x)}\frac{\mu _a^2}{p_a},$$ (26) and taking into account the approximation of Eq. (22) $$v_a(x)1\frac{\mu _a^2(x)}{2E_0^2}+\frac{_{E_0}\mu _a^2}{2E_0}.$$ (27) Hence $$_0^L𝑑x\mathrm{\Delta }v_{ab}(x)\frac{1}{2E_0^2}(1+E_0_{E_0})_0^L𝑑x\mathrm{\Delta }\mu _{ab}^2(x),$$ (28) i.e. we see that $`_0^L𝑑x\mathrm{\Delta }v_{ab}(x)`$ can be obtained easily if $`_0^L𝑑x\mathrm{\Delta }\mu _{ab}^2(x)`$ is known. ## 3 Examples As an example we calculate the effective oscillation and coherence lengths for two specific density profiles. For simplicity only two neutrino flavors, $`\nu _e`$ and $`\nu _\mu `$, are considered. The effective mass squared difference in matter is given by $$\mathrm{\Delta }_m(x)=\sqrt{(\mathrm{\Delta }_0\mathrm{cos}2\theta 2\sqrt{2}E_0G_FN_e(x))^2+\mathrm{\Delta }_0^2\mathrm{sin}^22\theta },$$ (29) where $`\mathrm{\Delta }_0`$ is the mass squared difference in vacuum, $`\theta `$ is the vacuum mixing angle, $`G_F`$ is the Fermi constant and $`N_e(x)`$ is the number density of electrons. (For other neutrino flavors or exotic matter contents $`N_e(x)`$ should be replaced by an appropriate combination of the particle densities of the matter.) The resonance, where the effective mixing is maximal, is reached at the density $$N_e(x_R)=\frac{\mathrm{\Delta }_0\mathrm{cos}2\theta }{2\sqrt{2}E_0G_F}.$$ (30) The passage through the resonance is adiabatic when $$Q\frac{\mathrm{\Delta }_0\mathrm{sin}^22\theta }{E_0\mathrm{cos}2\theta }\left\|\frac{N_e(x_R)}{N_e^{}(x_R)}\right\|1,$$ (31) where $`Q`$ is the adiabaticity parameter. It is also assumed that $`\|G_FN_e(x)\|E_0`$. One should notice that the matter may affect considerably the oscillation of neutrinos propagating in constant density. This is the case e.g. if $`N_e`$ is very close to its resonance value and $`\theta `$ is small. Then (see Eq. (29)) $`\mathrm{\Delta }_m\mathrm{\Delta }_0`$ and the oscillation length may be highly longer than in the vacuum case. Since at the surface of Earth $`\rho 3\text{g}/\text{cm}^3`$ and $`Y_e1/2`$, and hence $`G_FN_e10^{13}`$ eV, the effect might be observable in long baseline experiments e.g. for $`E_010\text{GeV}`$, $`\mathrm{\Delta }_010^3\text{eV}^2`$ (and small $`\theta `$). Also the group velocity difference, and consequently the coherence length, can be modified for suitable parameter values (cf. Eqs. (15) and (16) in Ref. ). At the resonance the coherence length is usually increased. ### 3.1 Linear density profile The linear profile is by far the most important example, because many actual profiles can be locally approximated by it. Let us parameterize the density as $$N_e(x)=\lambda (\kappa x),$$ (32) where $`\lambda `$ and $`\kappa `$ are parameters. The adiabaticity condition, Eq. (31), leads to $$\frac{(\mathrm{\Delta }_0\mathrm{sin}2\theta )^2}{E_0^2G_F\lambda }1.$$ (33) We first calculate $$_0^L𝑑x\mathrm{\Delta }_m(x)=\frac{1}{4c}[I(\kappa )I(\kappa L)],$$ (34) with (see e.g. ) $$I(x)=(2cx+b)\sqrt{a+bx+cx^2}+\frac{4acb^2}{2\sqrt{c}}\mathrm{ln}(2\sqrt{c(a+bx+cx^2)}+2cx+b),$$ (35) and $`a=\mathrm{\Delta }_0^2,b=4\sqrt{2}E_0G_F\lambda \mathrm{\Delta }_0\mathrm{cos}2\theta ,c=8(E_0G_F\lambda )^2`$. A direct substitution to Eq. (25) then yields the effective oscillation length, in principle. Using Eqs. (25), (28) and (34), the effective coherence length could also be calculated, but here that tedious calculation is omitted. Since Eq. (34) is not very illustrative, we will next consider the low density limit (i.e. small $`\lambda `$). One has $$_0^L𝑑x\mathrm{\Delta }_m(x)=L\mathrm{\Delta }_0+\sqrt{2}E_0G_F\lambda \mathrm{cos}2\theta (L^22\kappa L)+O(\lambda ^2)$$ (36) yielding $$L^{osc}\frac{4\pi E_0}{\mathrm{\Delta }_0}\left[1\frac{\sqrt{2}E_0G_F\lambda \mathrm{cos}2\theta }{\mathrm{\Delta }_0}(L2\kappa )\right],$$ (37) where the second term in the brackets can be regarded as the first order correction due to matter. The coherence length is to this order (with Eqs. (25) and (28)) $$L^{coh}\frac{4\sqrt{2}\sigma _xE_0^2}{\|\mathrm{\Delta }_0\|},$$ (38) i.e. the same as in vacuum. The lowest order correction of the coherence length is in fact always of the second order: the linear term $`E_0N_e(x)`$ (see Eq. (29)) is wiped out by the “operator” $`1+E_0_{E_0}`$ in Eq. (28). ### 3.2 Exponential density profile Let the density of electrons be given by $$N_e(x)=\lambda e^{\kappa x},$$ (39) where $`\lambda `$ and $`\kappa `$ are parameters. Eq. (31) yields in this case $$\frac{\mathrm{\Delta }_0\mathrm{sin}^22\theta }{E_0\mathrm{cos}2\theta \kappa }1.$$ (40) Now we can write $$_0^L𝑑x\mathrm{\Delta }_m(x)=\frac{1}{\kappa }[I(1)I(e^{\kappa L})],$$ (41) where $`I(x)`$ $`=`$ $`\sqrt{a+bx+cx^2}\sqrt{a}\mathrm{ln}\left({\displaystyle \frac{2a+bx+2\sqrt{a(a+bx+cx^2)}}{x}}\right)`$ (42) $`+{\displaystyle \frac{b}{2\sqrt{c}}}\mathrm{ln}(2\sqrt{c(a+bx+cx^2)}+2cx+b),`$ and $`a=\mathrm{\Delta }_0^2,b=4\sqrt{2}E_0G_F\lambda \mathrm{\Delta }_0\mathrm{cos}2\theta ,c=8(E_0G_F\lambda )^2`$. Once again, this result with Eqs. (25) and (28) would allow us to obtain the effective oscillation and coherence lengths in principle. In the low density limit we can expand $$_0^L𝑑x\mathrm{\Delta }_m(x)=L\mathrm{\Delta }_0+\frac{2\sqrt{2}E_0G_F\lambda \mathrm{cos}2\theta }{\kappa }(e^{\kappa L}1)\frac{2(E_0G_F\lambda \mathrm{sin}2\theta )^2}{\mathrm{\Delta }_0\kappa }(e^{2\kappa L}1)+O(\lambda ^3).$$ (43) Now we get the first order correction to the effective oscillation length $$L^{osc}\frac{4\pi E_0}{\mathrm{\Delta }_0}\left[1\frac{2\sqrt{2}E_0G_F\lambda \mathrm{cos}2\theta }{L\mathrm{\Delta }_0\kappa }(e^{\kappa L}1)\right].$$ (44) Eqs. (25) and (28) yield the correction to the effective coherence length $$L^{coh}\frac{4\sqrt{2}\sigma _xE_0^2}{\|\mathrm{\Delta }_0\|}\left[1\frac{2(E_0G_F\lambda \mathrm{sin}2\theta )^2}{L\mathrm{\Delta }_0^2\kappa }(e^{2\kappa L}1)\right],$$ (45) which is of the second order, in accordance with the remark made in the previous example. We conclude with a simple numerical application on solar neutrinos. The electron number density in the Sun is approximately $$N_e(x)=245N_A\mathrm{exp}\left(10.54\frac{x}{R_{}}\right)\text{cm}^3,$$ (46) where $`N_A`$ is Avogadro’s number. Hence (Eq. (39)) $`\lambda 10^{12}\text{eV}^3`$ and $`\kappa L10^3`$ if neutrinos are produced in the center of the Sun and detected on Earth. We use the values $`\mathrm{\Delta }_010^4\text{eV}^2`$ and $`E_01\text{MeV}`$, which fulfill the adiabaticity condition, Eq. (40), even for rather small values of $`\theta `$. Eq. (44) yields the correction of the effective oscillation length $$\frac{E_0G_F\lambda \mathrm{cos}2\theta }{L\mathrm{\Delta }_0\kappa }10^4\mathrm{cos}2\theta ,$$ (47) where $`G_F10^{23}\text{eV}^2`$. The correction of the effective coherence length is similarly $$\frac{(E_0G_F\lambda \mathrm{sin}2\theta )^2}{L\mathrm{\Delta }_0^2\kappa }10^5\mathrm{sin}^22\theta ,$$ (48) where Eq. (45) was used. It is thus seen that the matter effect modifies the oscillation and coherence lengths of solar neutrinos only slightly at least within the framework of this example. This fact is, of course, due to the insignificance of the density outside the Sun, i.e. the neutrinos propagate mainly in vacuum between the Sun and Earth. Similarly, it is to be believed that the matter effect is unimportant for the oscillation of all extraterrestrial neutrinos (excluding, of course, MSW effect and parametric resonance). We finally call your attention to the fact that even though both the linear and exponential density profiles violate the condition $`\|G_FN_e(x)\|E_0`$ for some values of $`x`$, the sharpness of the wave packets forces $`x`$ to be situated in harmless region (see Eqs. (12) and (13)). It is due to this same reason that the profile of Eq. (46), written in spherical coordinates, is directly applicable. A more obvious reason is seen in the integrals in Eq. (25). ## 4 Nonadiabatic neutrino propagation in matter ### 4.1 Solution of the equation of motion for two neutrino flavors In the previous sections we have assumed that the neutrino propagation is adiabatic. In most physical environments, however, the matter density may change so rapidly that the nonadiabaticity must be taken into account. When neutrinos propagate nonadiabatically the local matter “eigenstates” are not anymore eigenstates but they become mixed. This means that transitions (i.e. level crossings) between the matter states may occur at certain probability. Mathematically this fact manifests itself in the complete equation of motion (in the matter basis), which includes also nondiagonal terms, whereas in the adiabatic limit the corresponding equation is diagonal (cf. Eq. (5)). In this section we look for a formal solution to the two flavor equation of motion $$i_x\left(\begin{array}{c}\psi _1(x,E)\\ \psi _2(x,E)\end{array}\right)=\left(\begin{array}{cc}\frac{\mathrm{\Delta }(x)}{4E}& i\theta _m^{}(x)\\ i\theta _m^{}(x)& \frac{\mathrm{\Delta }(x)}{4E}\end{array}\right)\left(\begin{array}{c}\psi _1(x,E)\\ \psi _2(x,E)\end{array}\right),$$ (49) where $`\mathrm{\Delta }(x)\mu _2^2(x)\mu _1^2(x)`$ (this is not necessarily the same as in Eq. (29) if the intermediate matter contains e.g. muons), $`\theta _m^{}(x)=_x\theta _m(x)`$ and $`\theta _m(x)`$ is the mixing angle in matter. In the limit $`4E\|\theta _m^{}(x)\|\mathrm{\Delta }(x)`$, equivalent to the previously considered adiabaticity condition, Eq. (31) (for $`\nu _e\nu _\mu `$ oscillation), the nondiagonal terms can be neglected, and Eqs. (49) and (5) (with $`a=1,2`$) correspond to each other perfectly (discarding terms proportional to the identity matrix). Level crossings, on the other hand, are obviously caused by the nondiagonal terms $`\pm i\theta _m^{}(x)`$. These issues, as well as exact calculations of the level crossing probabilities for specific density profiles, have been extensively discussed by numerous authors; see e.g. ,,. One should remember that Eq. (49) describes the physical situation accurately enough if the following conditions hold (see e.g. Ref. ): 1. Neutrinos are relativistic. 2. The density is low enough, i.e. $`\|G_FN(x)\|E`$, where $`N(x)`$ is the number density of particle species relevant for a case under consideration (mentioned already in the previous section for $`N(x)=N_e(x)`$). 3. The density must not change appreciably over a length scale equal to the neutrino’s de Broglie wavelength. The expression “arbitrary density profile”, used in the following, is to be understood in the context of the abovementioned limitations. We now demonstrate how Eq. (49) can be solved. Defining (we omit here the energy dependence of the wave functions) $$\left(\begin{array}{c}\psi _1(x)\\ \psi _2(x)\end{array}\right)=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)\left(\begin{array}{c}\varphi _1(x)\\ \varphi _2(x)\end{array}\right)$$ (50) one has $$i_x\left(\begin{array}{c}\varphi _1(x)\\ \varphi _2(x)\end{array}\right)=\left(\begin{array}{cc}0& B(x)\\ B^{}(x)& 0\end{array}\right)\left(\begin{array}{c}\varphi _1(x)\\ \varphi _2(x)\end{array}\right),$$ (51) where $`B(x)\frac{\mathrm{\Delta }(x)}{4E}i\theta _m^{}(x)`$ (and $`B^{}(x)=\frac{\mathrm{\Delta }(x)}{4E}+i\theta _m^{}(x)`$). The solution of Eq. (51) is obtainable after some effort: $`\varphi _1(x)`$ $`=`$ $`C_1\mathrm{\Xi }_{}(B,B^{})+C_2\mathrm{\Xi }_+(B,B^{}),`$ $`\varphi _2(x)`$ $`=`$ $`C_1\mathrm{\Xi }_{}(B^{},B)C_2\mathrm{\Xi }_+(B^{},B),`$ (52) where $`C_1`$ and $`C_2`$ are constants, and $$\mathrm{\Xi }_\pm (Y,Z)1\pm iYYZiYZY+YZYZ\pm \mathrm{}$$ (53) with e.g. $`YZY_{x_0}^x𝑑x_1Y(x_1)_{x_0}^{x_1}𝑑x_2Z(x_2)_{x_0}^{x_2}𝑑x_3Y(x_3)`$. Eq. (50) yields finally $`\psi _1(x,E)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\psi _1(x_0,E)(\mathrm{\Xi }_{}(B,B^{})+\mathrm{\Xi }_{}(B^{},B))`$ $`{\displaystyle \frac{1}{2}}\psi _2(x_0,E)(\mathrm{\Xi }_+(B,B^{})\mathrm{\Xi }_+(B^{},B)),`$ $`\psi _2(x,E)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\psi _1(x_0,E)(\mathrm{\Xi }_{}(B,B^{})+\mathrm{\Xi }_{}(B^{},B))`$ (54) $`+{\displaystyle \frac{1}{2}}\psi _2(x_0,E)(\mathrm{\Xi }_+(B,B^{})+\mathrm{\Xi }_+(B^{},B)).`$ A straightforward substitution shows that this is really the solution of Eq. (49). As far as we know, no such complete solution has previously been presented in the literature. It is to be emphasized that Eq. (4.1) applies to arbitrary density profile, but on the other hand it has the undeniable deficiency of being somewhat formal and perhaps not very useful for practical calculations. Further discussion on the solution of Eq. (49) and on the corresponding equation in the flavor basis is found in Appendix B. Let us show briefly that our solution gives a meaningful result in two specific cases. In the adiabatic limit $`B(x)\frac{\mathrm{\Delta }(x)}{4E}B^{}(x)`$ and hence $$\psi _1(x,E)=\psi _1(x_0,E)\mathrm{\Xi }_{}(B,B)=\psi _1(x_0,E)\mathrm{exp}\left(i_{x_0}^x𝑑x^{}\frac{\mathrm{\Delta }(x)}{4E}\right),$$ (55) and similarly $$\psi _2(x,E)=\psi _2(x_0,E)\mathrm{exp}\left(i_{x_0}^x𝑑x^{}\frac{\mathrm{\Delta }(x)}{4E}\right),$$ (56) i.e. the correct result is obtained. In the extremely nonadiabatic limit<sup>1</sup><sup>1</sup>1Here we consider the usual textbook example where the matter states are related to electron and muon flavor states, and the matter does not contain muons. Remember that Eq. (49) is not necessarily limited to this standard case., on the other hand, $`B(x)i\theta _m^{}(x)B^{}(x)`$ (in the resonance region), yielding $`\psi _1(x,E)`$ $`=`$ $`\psi _1(x_0,E)\left(1+{\displaystyle BB}+\mathrm{}\right)\psi _2(x_0,E)\left(i{\displaystyle B}+i{\displaystyle BBB}+\mathrm{}\right),`$ $`\psi _2(x,E)`$ $`=`$ $`\psi _1(x_0,E)\left(i{\displaystyle B}+i{\displaystyle BBB}+\mathrm{}\right)+\psi _2(x_0,E)\left(1+{\displaystyle BB}+\mathrm{}\right).`$ (57) One may assume that $`\theta _m(x)`$ changes abruptly from $`\pi /2`$ to $`\theta `$ (the vacuum mixing angle) in the resonance, i.e. $$_{x_0}^xBi\left(\theta \frac{\pi }{2}\right)$$ (58) if $`x_0(x)`$ is before (after) the resonance. Hence $`\psi _1(x,E)`$ $`=`$ $`\psi _1(x_0,E)\mathrm{cos}\left(\theta {\displaystyle \frac{\pi }{2}}\right)\psi _2(x_0,E)\mathrm{sin}\left(\theta {\displaystyle \frac{\pi }{2}}\right),`$ $`\psi _2(x,E)`$ $`=`$ $`\psi _1(x_0,E)\mathrm{sin}\left(\theta {\displaystyle \frac{\pi }{2}}\right)+\psi _2(x_0,E)\mathrm{cos}\left(\theta {\displaystyle \frac{\pi }{2}}\right).`$ (59) Putting e.g. $`\psi _1(x_0,E)=0,\psi _2(x_0,E)=1`$, we see that the level crossing probability is $$\|\psi _1(x,E)\|^2=\mathrm{sin}^2\left(\theta \frac{\pi }{2}\right)=\mathrm{cos}^2\theta ,$$ (60) as it should in this limit . ### 4.2 Level crossing probabilities and wave packets We now present how the existing results on the level crossing probabilities can be combined with the wave packet description in a consistent manner. As in the previous section, we restrict to two flavors which is sufficient for understanding the relevant phenomena. Consider a neutrino that propagates initially as $`\psi _2`$ and has the usual Gaussian form (Eq. (8)), i.e. $$\psi _1(x_0,E)=0\text{and}\psi _2(x_0,E)=N\mathrm{exp}\left[\frac{(EE_2)^2}{4\sigma _{EP}^2}\right],$$ (61) where $`N=(2\pi \sigma _{EP}^2)^{1/4}`$. From Eq. (4.1) it follows that $`\psi _1(x,E)`$ $`=`$ $`{\displaystyle \frac{N}{2}}\mathrm{exp}\left[{\displaystyle \frac{(EE_2)^2}{4\sigma _{EP}^2}}\right]f(x_0,x,E),`$ $`\psi _2(x,E)`$ $`=`$ $`{\displaystyle \frac{N}{2}}\mathrm{exp}\left[{\displaystyle \frac{(EE_2)^2}{4\sigma _{EP}^2}}\right]h(x_0,x,E),`$ (62) where $`f(x_0,x,E)`$ and $`h(x_0,x,E)`$ represent the corresponding expansions in Eq. (4.1). Hence $$\psi _1(x,t)=\frac{N}{2\sqrt{2\pi }}𝑑E\mathrm{exp}\left[iEt\frac{(EE_2)^2}{4\sigma _{EP}^2}\right]f(x_0,x,E),$$ (63) where Eq. (2) was used. Defining $$\mathrm{exp}\left[\frac{(EE_2)^2}{4\sigma _{EP}^2}\right]f(x_0,x,E)g(x_0,x,E)$$ (64) one has $$\psi _1(x,t)=\frac{N}{2}\widehat{g}(x_0,x,t),$$ (65) where the circumflex stands for the Fourier transform. The level crossing probability is $`\|\psi _1(x,t)\|^2`$ but we integrate over time (cf. Sec. 2) and have $$\|\psi _1(x)\|^2=𝑑t\|\psi _1(x,t)\|^2=\frac{N^2}{4}𝑑t\|\widehat{g}(x_0,x,t)\|^2=\frac{N^2}{4}𝑑E\|g(x_0,x,E)\|^2,$$ (66) where Parseval’s identity was used. We can express the level crossing probability for wave packets using the respective probability for a plane wave, $`P_{lc}(E)=\frac{1}{4}\|f(x_0,x,E)\|^2`$. Hence $$P_{lc}(E_2,\sigma _{EP})\|\psi _1(x)\|^2=N^2𝑑E\mathrm{exp}\left[\frac{(EE_2)^2}{2\sigma _{EP}^2}\right]P_{lc}(E),$$ (67) where $`P_{lc}(E_2,\sigma _{EP})`$ is the generalized level crossing probability that takes into account the energy width of the wave packet. The wave packet effects can be seen more clearly by expanding $`P_{lc}(E)`$ in series. Assuming $`\sigma _{EP}`$ to be small (Sec. 2), it is sufficient to take the lowest terms, and Eq. (67) gives (with $`E_2E`$) $$P_{lc}(E,\sigma _{EP})=P_{lc}(E)+\frac{\sigma _{EP}^2}{2}\frac{^2P_{lc}(E)}{E^2}+O\left(\sigma _{EP}^4\frac{^4P_{lc}(E)}{E^4}\right).$$ (68) This equation clearly indicates that the use of the wave packets modifies the usual level crossing probabilities, $`P_{lc}(E)`$. It also shows that the wave packet correction is easily calculable if $`P_{lc}(E)`$ is known. Let us consider two simple examples: 1. Linear density profile. The well-known Landau-Zener probability is $$P_{LZ}(E)=\mathrm{exp}\left(\frac{\pi }{4}Q\right),$$ (69) where $`Q`$, given in Eq. (31), should not be too small. Using Eq. (68) and remembering that $`Q\frac{1}{E^2}`$ (see Eq. (33)) one has $$P_{LZ}(E,\sigma _{EP})=\left[1+\left(\frac{\sigma _{EP}}{2E}\right)^2\left(\frac{\pi ^2}{2}Q^23\pi Q\right)\right]\mathrm{exp}\left(\frac{\pi }{4}Q\right).$$ (70) For the most relevant cases (i.e. $`P_{LZ}(E)`$ not too small) $`(\pi Q)^2/23\pi Q`$ is of the order of unity. 2. Exponential density profile ($`N_e(x)e^{\kappa x}`$). Now $$P_{lc}(E)=\frac{e^{\pi \delta (1\mathrm{cos}2\theta )}e^{2\pi \delta }}{1e^{2\pi \delta }}=\frac{\mathrm{sinh}(BA)}{\mathrm{sinh}B}e^A,$$ (71) where $$\delta =\frac{\mathrm{\Delta }_0}{2E\kappa },A=\frac{\pi \delta }{2}(1\mathrm{cos}2\theta ),B=\pi \delta .$$ (72) A tedious calculation gives $$P_{lc}(E,\sigma _{EP})=P_{lc}(E)\left[1+\left(\frac{\sigma _{EP}}{E}\right)^2\mathrm{\Gamma }\right],$$ (73) where $$\mathrm{\Gamma }=A(1+\mathrm{coth}(BA))(A1+B(\mathrm{coth}B1))+B(\mathrm{coth}(BA)\mathrm{coth}B)(1B\mathrm{coth}B).$$ (74) Again, for the interesting range of parameters $`\mathrm{\Gamma }O(1)`$ at most when e.g. solar neutrinos are considered (with $`E=1\text{MeV}`$, $`\mathrm{\Delta }_0=10^4\text{eV}^2`$, and the density profile as given in Eq. (46)). The wave packet corrections turn out to be negligible for small $`\sigma _{EP}/E`$. This happens to be true in most physical environments, e.g. for solar neutrinos $`\sigma _{EP}/E10^410^5`$ . If $`\sigma _{EP}E`$, on the other hand, the simple wave packet treatment is inaccurate . This fact is manifest also in our calculation: if $`\sigma _{EP}`$ is not small enough, the integration in Eq. (67) becomes problematic since $`P_{lc}(E)`$ is not necessarily meaningfully defined for negative $`E`$ values. Anyway, there might be at least in principle a situation where the energy distribution of the neutrino wave function deviates considerably from a plane wave. Our calculation suggests that then the usual level crossing probabilities would not be totally reliable. Finally, it is to be noted that the use of some specific $`P_{lc}(E)`$ restricts the location of $`x_0`$ and $`x`$ in Eqs. (4.1) and (62); $`P_{LZ}(E)`$, for example, is valid only if $`x_0(x)`$ is situated well before (after) the resonance. If, on the contrary, $`x_0`$ and/or $`x`$ are/is in the resonance region, $`P_{LZ}(E)`$ cannot be used. ## 5 Summary We applied the wave packet formalism in order to study neutrino oscillations in matter in the adiabatic limit, and found out that the effective oscillation and coherence lengths take into account the whole path the neutrino has traversed. Results for the linear and exponential density profiles were briefly presented. The corrections for the predictions of observable fluxes of solar neutrinos seem to be quite small. On the other hand, the matter may affect significantly the oscillation of the neutrinos e.g. in long baseline experiments for suitable values of parameters. We then considered the equation of motion for two neutrino flavors, and managed to solve it formally for arbitrary density profile. Our method clearly applies to any differential equation of the same form. Finally, we showed that the level crossing probabilities of the wave packets differ from those of plane waves. The difference is practically equivalent to a simple average over energy (cf. Eq. (67)), not essentially related to the separation of wave packets. The finite width of the wave packet does not result in any observable effect for the physical situations we have considered. We could have continued our work by combining the wave packet treatment of Sec. 2 with the complete solution of the two flavor equation in Sec. 4.1. In that case the solution of Eq. (5) (for $`a=1,2`$) should be replaced by Eq. (4.1) in order to correctly take into account the effects due to nonadiabaticity. However, the presented results suggest that this would not reveal any new physics, so the calculation has been omitted. In this work we used a model that in principle is more physical and hence more accurate than the models used normally. On the other hand, since many of the calculations in this framework are quite complicated, we considered the limits of validity of the simpler plane wave approaches. Our results show that the present neutrino observations and the phenomena behind them can be described by a plane wave model accurately enough, regarding the current precision of the experiments. ## Acknowledgements We are indebted to J. Maalampi for discussions and careful reading of the manuscript. V.S. wishes to thank the Graduate School of Particle Physics (Finland) for financial support. ## Appendix A: Integration over $`x`$ in Eq. (14) We can write Eq. (14) as $$A_{\alpha \beta }(L,T)C(L,T)𝑑xe^{F(x)},$$ (A1) where $$F(x)\frac{v_a^2(0)}{4\sigma _{xP}^2}\left(_0^x\frac{dx^{}}{v_a(x^{})}T\right)^2\frac{v_a^2(L)}{4\sigma _{xD}^2}\left(_L^x\frac{dx^{}}{v_a(x^{})}\right)^2.$$ (A2) The definition of the saddle point is in turn $`F^{}(x_0)=0(F^{}(x)=\frac{d}{dx}F(x))`$, leading to $$\sigma _{xD}^2\left(_0^{x_0}\frac{dx^{}}{v_a(x^{})}T\right)v_a^2(0)+\sigma _{xP}^2_L^{x_0}\frac{dx^{}}{v_a(x^{})}v_a^2(L)=0.$$ (A3) From this one can solve that $$_0^{x_0}\frac{dx^{}}{v_a(x^{})}=\frac{\sigma _{xD}^2v_a^2(0)T+\sigma _{xP}^2v_a^2(L)_0^L\frac{dx^{}}{v_a(x^{})}}{\sigma _{xD}^2v_a^2(0)+\sigma _{xP}^2v_a^2(L)}$$ (A4) and $$_L^{x_0}\frac{dx^{}}{v_a(x^{})}=\frac{\sigma _{xD}^2v_a^2(0)\left(T_0^L\frac{dx^{}}{v_a(x^{})}\right)}{\sigma _{xD}^2v_a^2(0)+\sigma _{xP}^2v_a^2(L)},$$ (A5) and finally after some algebra $$F(x_0)=\frac{1}{4}\frac{\left(T_0^L\frac{dx}{v_a(x)}\right)^2}{\left(\frac{\sigma _{xP}}{v_a(0)}\right)^2+\left(\frac{\sigma _{xD}}{v_a(L)}\right)^2}.$$ (A6) Let us now examine the higher derivatives of $`F(x)`$. The second derivative is $`F^{\prime \prime }(x)`$ $`=`$ $`{\displaystyle \frac{1}{2\sigma _{xP}^2}}\left({\displaystyle \frac{v_a(0)}{v_a(x)}}\right)^2+{\displaystyle \frac{1}{2\sigma _{xP}^2}}\left({\displaystyle _0^x}{\displaystyle \frac{dx^{}}{v_a(x^{})}}T\right)\left({\displaystyle \frac{v_a(0)}{v_a(x)}}\right)^2v_a^{}(x)`$ (A7) $`{\displaystyle \frac{1}{2\sigma _{xD}^2}}\left({\displaystyle \frac{v_a(L)}{v_a(x)}}\right)^2+{\displaystyle \frac{1}{2\sigma _{xD}^2}}\left({\displaystyle _L^x}{\displaystyle \frac{dx^{}}{v_a(x^{})}}\right)\left({\displaystyle \frac{v_a(L)}{v_a(x)}}\right)^2v_a^{}(x).`$ In the relativistic limit $`v_a(x)=1+O(\mu _a^2/E^2)`$, and Eq. (A7) reduces to<sup>2</sup><sup>2</sup>2The adiabaticity further reinforces the smallness of $`v_a^{}(x)`$. $$F^{\prime \prime }(x_0)=\frac{1}{2\sigma _{xP}^2}\frac{1}{2\sigma _{xD}^2}+O\left(\frac{LT}{\sigma _{xP,D}^2}\frac{d}{dx_0}\frac{\mu _a^2(x_0)}{E^2}\right),$$ (A8) where effectively $`\sigma _{xP,D}^2=\mathrm{max}\{\sigma _{xP}^2,\sigma _{xD}^2\}`$. From Eq. (A7) one sees also that the third derivative of $`F(x)`$ includes terms proportional to $`v_a^{}(x),(v_a^{}(x))^2`$ and $`v_a^{\prime \prime }(x)`$. It is thus obvious that $`F^{\prime \prime \prime }(x)`$ and all the higher derivatives are negligible in the relativistic limit, and we can approximate (with $`F^{}(x_0)=0`$) $$F(x)F(x_0)+\frac{1}{2}F^{\prime \prime }(x_0)(xx_0)^2.$$ (A9) The relevant part of the integral of Eq. (14) becomes $$𝑑xe^{F(x)}e^{F(x_0)}𝑑x\mathrm{exp}\left[\frac{1}{4}\left(\frac{1}{\sigma _{xP}^2}+\frac{1}{\sigma _{xD}^2}\right)(xx_0)^2\right],$$ (A10) where the final Gaussian integral gives only an unimportant numerical factor. ## Appendix B: Additional calculations in connection with the two flavor equation of motion We present another way of solving Eq. (49). Instead of Eq. (50), we now define (compare to Ref. ) $$\left(\begin{array}{c}\psi _1(x)\\ \psi _2(x)\end{array}\right)=\left(\begin{array}{cc}\mathrm{exp}\left(i_{x_0}^x𝑑x^{}\frac{\mathrm{\Delta }(x^{})}{4E}\right)& 0\\ 0& \mathrm{exp}\left(i_{x_0}^x𝑑x^{}\frac{\mathrm{\Delta }(x^{})}{4E}\right)\end{array}\right)\left(\begin{array}{c}\varphi _1(x)\\ \varphi _2(x)\end{array}\right),$$ (B1) leading to $$i_x\left(\begin{array}{c}\varphi _1(x)\\ \varphi _2(x)\end{array}\right)=\left(\begin{array}{cc}0& iD(x)\\ iD^{}(x)& 0\end{array}\right)\left(\begin{array}{c}\varphi _1(x)\\ \varphi _2(x)\end{array}\right),$$ (B2) where $$D(x)\theta _m^{}(x)\mathrm{exp}\left(i_{x_0}^x𝑑x^{}\frac{\mathrm{\Delta }(x^{})}{2E}\right).$$ (B3) Eq. (B2) yields easily $`\varphi _1(x)`$ $`=`$ $`C_1\left(1{\displaystyle D}{\displaystyle DD^{}}+{\displaystyle DD^{}D}+{\displaystyle DD^{}DD^{}}+\mathrm{}\right)`$ $`+C_2\left(1+{\displaystyle D}{\displaystyle DD^{}}{\displaystyle DD^{}D}+{\displaystyle DD^{}DD^{}}+\mathrm{}\right),`$ $`\varphi _2(x)`$ $`=`$ $`C_1\left(1+{\displaystyle D^{}}{\displaystyle D^{}D}{\displaystyle D^{}DD^{}}+{\displaystyle D^{}DD^{}D}+\mathrm{}\right)`$ $`C_2\left(1{\displaystyle D^{}}{\displaystyle D^{}D}+{\displaystyle D^{}DD^{}}+{\displaystyle D^{}DD^{}D}+\mathrm{}\right),`$ where the notation is as before (see Eq. (53)). Finally one has $`\psi _1(x)`$ $`=`$ $`\mathrm{exp}\left(i{\displaystyle _{x_0}^x}𝑑x^{}{\displaystyle \frac{\mathrm{\Delta }(x^{})}{4E}}\right)`$ $`\times [\psi _1(x_0)(1{\displaystyle }D{\displaystyle }D^{}+{\displaystyle }D{\displaystyle }D^{}{\displaystyle }D{\displaystyle }D^{}+\mathrm{})`$ $`+\psi _2(x_0)({\displaystyle }D+{\displaystyle }D{\displaystyle }D^{}{\displaystyle }D+\mathrm{})],`$ $`\psi _2(x)`$ $`=`$ $`\mathrm{exp}\left(i{\displaystyle _{x_0}^x}𝑑x^{}{\displaystyle \frac{\mathrm{\Delta }(x^{})}{4E}}\right)`$ (B5) $`\times [\psi _1(x_0)({\displaystyle }D^{}{\displaystyle }D^{}{\displaystyle }D{\displaystyle }D^{}+\mathrm{})`$ $`+\psi _2(x_0)(1{\displaystyle }D^{}{\displaystyle }D+{\displaystyle }D^{}{\displaystyle }D{\displaystyle }D^{}{\displaystyle }D+\mathrm{})],`$ where the energy dependence is not written down explicitly. We thus see that Eq. (4.1) is not the only form in which the solution of Eq. (49) can be expressed. The equation of motion in the flavor basis is (and many others) $$i_x\left(\begin{array}{c}\psi _e(x)\\ \psi _\mu (x)\end{array}\right)=\frac{1}{4E}\left(\begin{array}{cc}\mathrm{\Delta }_0\mathrm{cos}2\theta +A_c(x)& \mathrm{\Delta }_0\mathrm{sin}2\theta \\ \mathrm{\Delta }_0\mathrm{sin}2\theta & \mathrm{\Delta }_0\mathrm{cos}2\theta A_c(x)\end{array}\right)\left(\begin{array}{c}\psi _e(x)\\ \psi _\mu (x)\end{array}\right),$$ (B6) where $`A_c(x)=2\sqrt{2}EG_FN_e(x)`$ and other notations are obvious. Since the calculation proceeds as above, we omit the details and content ourselves with giving the final answer: $`\psi _e(x)`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{i}{4E}}{\displaystyle _{x_0}^x}𝑑x^{}(\mathrm{\Delta }_0\mathrm{cos}2\theta A_c(x^{}))\right]`$ $`\times [\psi _e(x_0)(1{\displaystyle }F{\displaystyle }F^{}+{\displaystyle }F{\displaystyle }F^{}{\displaystyle }F{\displaystyle }F^{}+\mathrm{})`$ $`+\psi _\mu (x_0)(i{\displaystyle }F+i{\displaystyle }F{\displaystyle }F^{}{\displaystyle }F+\mathrm{})],`$ $`\psi _\mu (x)`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{i}{4E}}{\displaystyle _{x_0}^x}𝑑x^{}(\mathrm{\Delta }_0\mathrm{cos}2\theta A_c(x^{}))\right]`$ (B7) $`\times [\psi _e(x_0)(i{\displaystyle }F^{}+i{\displaystyle }F^{}{\displaystyle }F{\displaystyle }F^{}+\mathrm{})`$ $`+\psi _\mu (x_0)(1{\displaystyle }F^{}{\displaystyle }F+{\displaystyle }F^{}{\displaystyle }F{\displaystyle }F^{}{\displaystyle }F+\mathrm{})],`$ where $$F(x)\frac{\mathrm{\Delta }_0\mathrm{sin}2\theta }{4E}\mathrm{exp}\left[\frac{i}{2E}_{x_0}^x𝑑x^{}(\mathrm{\Delta }_0\mathrm{cos}2\theta A_c(x^{}))\right].$$ (B8)
warning/0004/astro-ph0004019.html	ar5iv	text	# Imaging of the protoelliptical NGC 1700 and its globular cluster system ## 1 Introduction Merging is now thought to be a key process in the evolution of galaxies. The hypothesis that two colliding spiral galaxies will eventually form an elliptical galaxy (?) has gained much observational and theoretical support over the years. The ‘smoking gun’ of this type of merger is the presence of two tidal tails, formed from the progenitor’s discs. A number of well known examples are found in the local Universe, such as The Mice, The Antennae, Arp 220, NGC 3256 and NGC 7252. For such galaxies a variety of methods are available (e.g. spectroscopy of the stellar populations, dynamical measurements, comparison with models etc.) to estimate the time since the merger occurred. The derived ages, for these classic tidal tail systems, is up to 1–2 Gyr since nuclear coalescence. It was Ivan King in 1977 who first pointed out the general lack of obvious candidates for older merger remnants (i.e. 2–5 Gyr old). These galaxies have been referred to as ‘King gap objects’ or protoellipticals. Identifying, and age dating, these protoellipticals could provide the missing link between late stage spiral mergers and elliptical galaxies. Indeed, a crucial step in testing the merger hypothesis would be to show an ‘evolutionary consistency’, i.e. that spiral mergers and protoellipticals evolve to have the same energetic, structural, dynamical and chemical properties as normal, old ellipticals. This has been problematic due to the difficulty of estimating the age of old stellar populations, without telltale morphological signatures such as tidal tails. Recently, several different methods have become available. These include breaking the age–metallicity degeneracy with spectral synthesis (e.g. ?) and new models (e.g. ?), quantifying optical fine structure (?) and using the colours of globular clusters (e.g. ?). An ideal candidate for a nearby protoelliptical is NGC 1700. It possesses a kinematically distinct core (Franx, Illingworth & Heckman 1998a), reveals evidence for extensive morphological disturbance (?) and it has a high rotational velocity to velocity dispersion ratio, relative to other ellipticals (Bender, Burstein & Faber 1992). As well as being a possible protoelliptical, it offers the opportunity to compare various age estimates. For example, age estimates can be made from its two faint tidal tails, optical fine structure, globular cluster colours, kinematic structure, and its stellar component. We have adopted the same distance to NGC 1700 as ?, i.e. 51.4 Mpc (which includes a correction for Virgocentric infall and assumes $`H_{}=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup>). This corresponds to 249 pc per arcsec. The total $`B`$ band magnitude for NGC 1700 is $`M_\mathrm{B}=21.56`$ (RC3:?). It is classified as an E4 elliptical in the RC3 and E3 in the RSA. Photometric studies have been carried out by Franx, Illingworth & Heckman (1998b) and ?. A recent kinematic study is that of Bender, Saglia & Gerhard (1994). In this paper we reanalyse Hubble Space Telescope (HST) WFPC2 $`V`$, $`I`$ images and present new Keck $`B`$, $`V`$ and $`I`$ images. We focus on the morphology of the galaxy and photometry of the globular cluster system, both of which provide new age estimates. We compare our age estimates with a variety of alternative methods and discuss the implications of our results for the formation and evolution of elliptical galaxies. ## 2 Observations and initial data reduction ### 2.1 Keck data New $`B`$, $`V`$ and $`I`$ images of NGC 1700 were obtained using the 10 m Keck–I telescope at the W. M. Keck Observatory, Mauna Kea, Hawaii. The observations were carried out on 1997 September 30th, with seeing in each filter of $`1.3`$ arcsec. The LRIS instrument was used with the Tek $`2048\times 2048`$ CCD, giving a rectangular field–of–view of $`6\times 8`$ arcmin and a pixel scale of 0.215 arcsec pixel<sup>-1</sup>. A single long (600 s) exposure was taken in $`V`$ and $`I`$ filters and two 600 s exposures were taken in the $`B`$ band. We also took a short (10 s) exposure in each filter. For the purpose of photometric calibration we obtained exposures in each filter of the field SA98 from ?. The data were reduced in a standard manner using iraf software. After bias subtraction there remained a small offset between the two halves of the CCD. This is due to the dual amplifier mode of the CCD read-out. This effect is reasonably well approximated by a step function which varies perpendicular to the read-out direction. The bulk of this offset was successfully removed after division by the normalised flat-field. A small offset ($`13`$ per cent) remained, which was corrected for by multiplying one side by a factor determined by taking the mean values of several sky regions either side of the offset. This successfully removed the offset for most of the image. However, an offset remained in the bright regions at the very centre of the galaxy ($``$ few arcsec). This may be due to some non-linearity in the gains of the amplifiers at high count rates (as suggested in the LRIS manual). This problem clearly affects galaxy photometry in the inner regions but as described below, does not affect standard star or globular cluster (GC) photometry. Photometric calibration was performed using our standard star exposures of SA98 and the photometric measurements of ?. The resulting rms error in the final photometry is $`\pm 0.02`$ mag in all three bands. K-correction and Galactic extinction were included using the same values for NGC 1700 as in ?, 0.18, 0.11 and 0.07 in the $`B`$, $`V`$ and $`I`$ bands respectively. In order to investigate the random errors on our zeropoints we compared the measurements for the same stars that appear in each pointing. We found a random error of $`0.02`$ mag. To better quantify the systematic error due to the dual amplifier offset, we compared the mean zeropoints for the stars on the left and stars on the right of the offset in the same frame. We found a systematic error due to the offset of $`0.03`$ mag. ## 3 Data analysis ### 3.1 HST data The $`V`$ and $`I`$ band HST data of NGC 1700 were obtained from proposal G0–5416 in the HST archive. Photometry of the GCs in NGC 1700 has been discussed by ?. Here we reanalyse the HST data, and use the results to aid with GC selection in our Keck data. We performed the photometry in 2 pixel radius apertures with sextractor (?), applying the aperture correction given by ? and the calibration and transformation to Johnson magnitudes following ?. From the initial object list of 383 sources detected by sextractor we applied the following selection criteria: The magnitude range was restricted to be $`21.5<V<26.5`$, the faint limit chosen to avoid introducing colour bias effects due to lack of completeness. Galactic GCs at the distance of NGC 1700 are expected to be at magnitudes of $`V>22.5`$. However, in order to account for the possible presence of a younger population of GCs that might be up to $`1`$ mag brighter than those in the Milky Way, we relaxed our bright limit to be $`V=21.5`$. The object colours were restricted to $`0.6<VI<1.7`$ (which is roughly equivalent to $`2.5<[\mathrm{𝐹𝑒}/H]<+1`$). Our typical error is $`\pm 0.1`$ in magnitude and $`\pm 0.15`$ in colour. Finally we carried out a careful visual check of the objects on the image display in order to remove point-like or galaxy-like objects. The final list contains a total of 146 GCs from all 4 WFPC2 chips. ### 3.2 Keck data #### 3.2.1 Galaxy morphology In order to better reveal the fine structure present within NGC 1700 it is necessary to subtract off a model of the underlying galaxy. All modelling was performed using the isophote package in stsdas (e.g. ?). We initially constructed a model for the purpose of sky subtraction. Elliptical isophotes were fitted with a fixed centre and employing a 3-$`\sigma `$ clipping algorithm. The ellipticity and position angle of the isophotes were allowed to vary freely. A pixel mask file was also constructed to explicitly mask out stars and pixel defects. We were unable to perform fits for the long exposure images at small radii because the central regions were saturated. For these images we set the initial fit radius beyond the saturated region. The fits were allowed to progress outwards until constraints on signal to noise $`(S/N)`$ terminated the process. Sky subtraction was performed in a similar way to that detailed in ?, i.e. the model intensity profiles in the outer parts of the galaxy were fitted by a power-law in order to accurately determine the sky level. The appropriate level was then subtracted from each image. The sky-subtracted images were then re-modelled in an identical way to before, to reveal the structure within the body of the galaxy. The models typically reached out to a distance of $`180`$ arcsec (45 kpc) from the galaxy centre. Our fits extend significantly further from the centre than the imaging of either ? or ?. Using the short exposure image, we estimate a total $`B`$ band magnitude for the galaxy of $`12.00\pm 0.05`$, in excellent agreement with that given by the RC3 of 12.01. We also estimate a total $`V`$ magnitude of $`11.18\pm 0.05`$ which agrees reasonably well with the RC3 value of 11.10, and a total $`I`$ magnitude of $`9.80\pm 0.10`$ compared with 9.87 from ?. We show in Fig. 1 the radial profiles of the ellipticity, position angle (PA), third cosine term (C3) and the fourth cosine term (C4) of the fitted isophotes for the short $`B`$ band image. All radii are expressed in terms of ‘equivalent radius’, $`r_{\mathrm{eq}}=a\sqrt{1ϵ}`$, where $`a`$ is the major axis radius and $`ϵ`$ is the isophote ellipticity. The C4 term is perhaps the most interesting of a Fourier series which expresses the deviations from a pure ellipse at a given radius. This term describes whether the isophotes are discy (C4 is positive) with excess light along the major and/or minor axes, or boxy (C4 is negative), representing excess light at 45° with respect to these axes. There is good agreement between the shape and absolute values of our data and that of ? and ?. The slight discrepancies at small radii are almost certainly due to seeing effects. Our profiles in $`V`$ and $`I`$ were similar in both shape and absolute value to that in the $`B`$ band. #### 3.2.2 Globular cluster sample selection and photometry As with the HST data, potential GC candidates were detected in the Keck images using the sextractor program. sextractor provides measures of FWHM, ellipticity and an indication of whether an object was ‘pointy’ (star-like) or extended (galaxy-like). We found 615 candidates in common between the $`B`$, $`V`$ and $`I`$ band images. Using this object list, we performed photometry with the iraf task phot . After a curve of growth analysis on several sources we determined an optimum aperture size of 8 pixels with a suitable background annulus of 15 to 20 pixels. In order to check the photometry and to determine suitable selection criteria, we examined the 34 sources in common between the $`V`$ band Keck image and our final $`V`$ band HST globular cluster list. Similarly, 14 candidates were found in common between our final HST list and the Keck $`I`$ band image. Although there was large scatter between the HST and Keck photometry at faint magnitudes, there was no obvious systematic bias. Our selection criteria are determined from the sample of 34 $`V`$ band sources in common between the HST and the Keck frames. Due to HST’s high spatial resolution, we are confident that the vast majority of sources that appear in the final HST list are bona fide GCs and hence use them to refine our Keck candidate list. The following criteria were applied to the 615 objects in common between the Keck $`B`$, $`V`$ and $`I`$ frames: (a) $`21.5<V<25.0`$, (b) $`2.5<\mathrm{FWHM}<12.0`$ (c) ellipticity $`<0.7`$. These selection cuts are shown in Fig. 2. As with our HST magnitude selection, the bright limit was chosen to include a population of young GCs while excluding bright foreground stars. Our faint magnitude cut-off was chosen in order to avoid introducing any colour bias into our results. After magnitude selection, 507 candidate GCs remained in our sample. The selection using the sextractor FWHM and ellipticity measures excluded a further 30 sources. These selection cuts were based on the shape parameters observed for the GCs in common with the HST list. We also excluded a further 25 sources because their $`x`$ or $`y`$ centre positions as determined by phot deviated from those originally found by sextractor by more than 2 pixels, and were thus liable to be misidentifications. The full range of expected GC metallicities is $`2.5<[\mathrm{𝐹𝑒}/H]<+1.0`$. Using the Galactic colour–metallicity relation of Couture, Harris & Allwright (1990), this corresponds to a range in colour of $`1.2<BI<2.5`$. The errors in our photometry lead to an average error in colour of $`\pm 0.4`$ mag. We thus decided to relax our final colour selection to $`0.8<BI<3.0`$. Our final colour selection is shown by the dotted box in Fig. 3. Here the range in $`VI`$ is defined as linear functions of $`BI`$. This colour selection reduced our sample size down to a list of 352 objects. As a final stage in selection we performed a visual check of all the 352 candidate GCs in our colour selected list. From this inspection we rejected a further 40 sources that resembled galaxies (i.e. appeared diffuse and/or ellipsoidal) leaving a final list of 312 GCs. Although these selection cuts are designed to reduce contamination to a minimum, there is always the concern that our final GC sample will still be contaminated by background galaxies and foreground stars. An estimate of the number of contaminating sources can be made by taking an image of a nearby ‘blank’ field. However, as we did not have such an image, we had to use an alternative method. ? give the predicted stellar densities in 17 fields based on their model of the Milky Way. Using the star densities in the field closest in direction to our observations of NGC 1700 (field 13), and correcting for our field–of–view, we predict only 41 stars for our magnitude range of $`21.5<V<25.0`$. Similarly we predict a total of 54 stars in our $`I`$ band field (using limiting magnitudes of $`20.5<I<23.5`$). The equivalent stellar density for the $`B`$ band predicts only 29 stars in our field (assuming $`B`$ band limiting magnitudes of $`22.0<B<25.5`$). As our selection criteria requires sources to be in each of our $`B`$, $`V`$ and $`I`$ images, we conclude that $`29`$ foreground stars that could be present do not make up a significant contribution to our candidate sample. Differential galaxy counts in the $`B`$ band are given in fig. 1 of ?. From this we estimate a total of 1343 background galaxies in our $`B`$ band frame down to our limiting magnitude of $`B=25.5`$. Using ? we estimate a total of 1130 galaxies in our $`I`$ band frame (to $`I=23.5`$). Once again, as we selected only those sources in common between the $`B`$, $`V`$ and $`I`$ frames we can take the lower of these estimates as an indication of the number of contaminating galaxies in our list. These numbers are high compared to the 615 objects detected in common between our $`B`$, $`V`$ and $`I`$ frames, though the number of background galaxies actually detected will be significantly lower due to the incompleteness of the sample at faint magnitudes. Although potentially large in number, we show in Section 4.2.2 that the additional constraint of colour effectively excludes the majority of galaxies from our candidate GC list. ## 4 Results and discussion ### 4.1 HST data Two $`VI`$ histograms for the GCs in our HST sample are shown in Fig. 4. The lower panel shows the colour histogram for the whole 146 GC sample. There is a single peak at around $`VI=1.07`$, which is consistent with that of ?, i.e. $`VI=1.05`$. The upper panel shows the $`VI`$ colour distribution for GCs in our sample that are brighter than $`V=24.5`$. At $`V=24.5`$, our typical error in colour is $`0.1`$ mag. Any fainter than this and the two distributions merge into a single broad peak due to photometric errors. For the bright sample, there is a suggestion of the presence of two peaks separated by $`\mathrm{\Delta }(VI)=0.30\pm 0.07`$, with one peak at $`VI=0.85\pm 0.05`$ and one at $`VI=1.15\pm 0.05`$. As we have a small number of GCs brighter than $`V=24.5`$ we use the dip statistic (see ??) in order to confirm this bimodality. For our bright sample of HST GCs, the dip statistic indicates a probability of $`>90`$ per cent that the distribution is not unimodal. The two peaks appear to have been ‘washed out’ at fainter magnitudes due to the photometric errors. The shaded region in this figure represents the $`(VI)_0`$ histogram of Galactic GCs. This distribution shows a sharp peak at $`(VI)_00.9`$, i.e. slightly bluer than the peak of the full sample but consistent with the blue peak of the $`V<24.5`$ population. An HST study of the NGC 1700 GC system prior to that of ? was performed by ?. Although only 39 GCs were detected, there was a hint of possible bimodality in the $`VI`$ histogram with peaks at $`VI0.9`$ and $`VI1.2`$. Another HST observation of NGC 1700 was carried out by Richstone et al. (see ?). An analysis of these observations detected 27 GCs, again with two possible peaks in the colour distribution at $`VI0.9`$ and $`VI1.2`$. However, a statistical analysis of the $`VI`$ distribution detected no significant bimodality (?). The positions of the peaks in Fig. 4 are consistent with those suggested by the previous observations mentioned above. ### 4.2 Keck data #### 4.2.1 Galaxy morphology and age estimates In Fig. 5 we show the ‘residual image’ of NGC 1700 produced by subtracting the model from the 600 s $`V`$ band image. Here the fine structure within the galaxy is better revealed. Most notable are the two broad, faint tidal tails or plumes visible to the North-West and South-East of the galaxy. In addition many GCs are apparent. The projected extent of the tidal features is about 165 arcsec (41 kpc) from the galaxy centre. A faint shell system is just visible in the central region (i.e. within the central $`25`$ arcsec) of the galaxy, as found in ?. In Fig. 6 we show the profile of the 4th cosine (C4) parameter in the $`B`$, $`V`$ and $`I`$ bands for the long exposure images. The C4 term is generally positive at radii $`\stackrel{<}{}30`$ arcsec, representing discy isophotes. At greater than 30 arcsec the isophotes become strongly boxy, indicating an excess of light at $`45\mathrm{°}`$ from the major and minor axes. Fig. 5 shows that this is very probably due to some light from the tail-like structures which remains in the galaxy model. To investigate this further, the ‘tails’ were masked out and the galaxy remodelled. Fig. 6 also shows the resulting C4 profiles after masking the tails. Note that the boxiness in the original fits has been greatly reduced in all cases (to less than $`2`$ per cent deviation), indicating that the dominant cause of the boxiness at these radii is the tails. The tails themselves appear clearly brighter in the residual maps produced by subtraction of these models from the original images. This shows that without masking, some of the light from the tails is included in the model, modifying the fit parameters. ? constructed a ‘merger sequence’ from a sample of galaxies that are good candidates for ongoing mergers and remnants of mergers between two approximately equal mass disc galaxies. They assigned a ‘merger stage’ to each galaxy based upon dynamical crossing times, with zero age defined to be the point of nuclear coalescence. They noted that the fraction of galaxy starlight contained within the tails roughly anti-correlates with this merger stage. This allows us make a very rough estimate of the time-scale since the merger event that created the tails within NGC 1700. This analysis was performed on the long exposure residual images. As we wanted to include as much of the tail light in the residual images as possible, we used the residual images produced by excluding the tails from the elliptical fit of the galaxy. As mentioned above, if the tails are not masked out during the fit, some of their light is included in the model, modifying the fit parameters and resulting in the subtraction of a significant amount of tail light from the galaxy image. The total flux contained within $`300`$ small circular apertures positioned on the tails was measured using the iraf utility imexamine. This was used to derive a mean surface brightness for the tails, which was then multiplied by their total area (i.e. including the area missed due to contaminating bright sources). The resulting surface brightnesses were 26.6, 25.9 and 24.4 mag arcsec<sup>-2</sup> for the $`B`$, $`V`$ and $`I`$ bands respectively. The total light in the tails was divided by the total galaxy light to give the ‘tail fraction’. We also roughly estimated the total tail light using polygon-shaped apertures. This method gave a similar result to the mean surface brightness method but included the light from contaminating point sources and was thus less reliable. The tail fractions in the $`B`$, $`V`$ and $`I`$ bands were found to be $`1.84`$, $`1.64`$ and $`1.72`$ per cent respectively. The tail fractions in ? were derived in the $`V`$ band, with a few exceptions. We thus use our $`V`$ band tail fraction (which is similar to the $`B`$ and $`I`$ values) for NGC 1700 in the following analysis. We fitted a linear least-square to the points obtained from the ? data. This fit is shown as a dashed line in Fig. 7. The solid horizontal line represents the measured $`V`$ band tail fraction for NGC 1700. If we extrapolate the fit we find that the $`V`$ band tail fraction measured for NGC 1700 corresponds to a merger stage of $`8.1\pm 3.8`$. The large error on this stage estimate is due to the scatter of the data points of ?. We also performed a quadratic fit on the Keel & Wu data. Extrapolating this fit yielded a stage comparable with the result using a linear least-square fit. By comparing the spectroscopic (i.e. central starburst) ages of several galaxies in the ? sample with their assigned merger stage, we determined the approximate time since nuclear coalescence. For the stage of NGC 1700, the time since the central starburst and hence tail formation is approximately $`3.2\pm 1.5`$ Gyr. A lower limit for the age of the tidal tails can be estimated from the dynamical time-scale, i.e. $`t_{\mathrm{dyn}}(R/V)`$. The projected distance of the tidal tails from the galaxy centre is about 165 arcsec (41 kpc), and the rotation velocity in the outer parts is $`50`$ km s<sup>-1</sup> (?). This gives a dynamical time of 0.8 Gyr, which is an underestimate if the tails do not lie in the plane of the sky; thus the tidal tails are $`>0.8`$ Gyr old. This lower limit is consistent with the age derived from the fraction of galaxy light contained within the tails. The presence of two symmetric tidal tails is generally taken to be a signature of a recent major merger involving two, approximately equal mass spiral galaxies (see e.g. ?). If the tidal features in Fig. 5 are indeed genuine tidal tails then we could conclude that NGC 1700 has experienced a major merger during its recent history. Alternatively, if the features are merely plumes of tidally disturbed material, the situation is less clear. While a major merger could not be ruled out, the situation of a disc galaxy merging into an existing elliptical would be possible. #### 4.2.2 Age estimates from globular cluster colours and magnitudes We present in Fig. 8 the $`BI`$ colour histogram for NGC 1700 GCs from the Keck data. The histogram appears bimodal with a blue peak at $`BI=1.54\pm 0.05`$ and a second peak $`0.44\pm 0.07`$ magnitudes redder at $`BI=1.98\pm 0.05`$. This bimodality does not appear to be an artifact of the data binning and is still present if the histogram bin boundaries are changed. For the subsequent discussion and analysis we define the blue population as those GCs possessing $`0.8<BI1.75`$ and the red population with colours $`1.75<BI<3.0`$. A statistical analysis using the kmm algorithm (Ashman, Bird & Zepf 1994) detects bimodality in the distribution with $`>99`$ per cent confidence. The kmm algorithm assigned a colour cut between the blue and red populations of $`BI=1.8`$, thus confirming our initial visual estimate. The $`BI`$ distribution of GCs in the Milky Way is also shown in Fig. 8 (shaded area). The peak of this distribution is at $`(BI)_01.5`$ which is similar to the peak of the blue population of NGC 1700. Is the bimodality of the $`BI`$ histogram evidence for two distinct GC populations in NGC 1700 ? In order to address this question, one has to consider the expected number of contaminating sources within our final sample. As the number of predicted stars before colour selection is small, we can immediately conclude that the contamination by foreground stars in our final sample is negligible. Another source of contamination is background galaxies. Our automatic and visual checks have removed obvious galaxies but it is possible that small, unresolved background galaxies remain. For our mean $`B`$ magnitude of 24.5, we expect to be detecting sources out to a mean redshift of $`z0.8`$ (?). At this redshift all morphological types of galaxies (with the exception of irregular types) have typical $`BI`$ colours in excess of 2.5. This is significantly redder than the peak of the red population and we can thus be fairly confident that the two peaks are real and due to two distinct populations of GCs. To support this view we performed photometry on 13 galaxy-like objects that appeared in our initial list but were rejected during the selection process. The mean colour of these objects was $`BI=2.7\pm 0.2`$; significantly redder than the peak of our red population of objects in our final sample. Moreover, the mean FWHM value measured for these objects was $`12.3\pm 1.2`$ pixels; greater than the upper cut-off used in our selection based on FWHM. The study of globular cluster systems can reveal important information regarding the formation and evolution of the parent galaxy. The merger model of galaxy formation (?; ?) makes a number of predictions about the properties of the GC system of the resultant galaxy. In this scenario, an elliptical galaxy is formed by the gas rich merger of two spiral galaxies. If our interpretation of two, roughly symmetric tidal tails in NGC 1700 is correct, then the galaxy has probably undergone such an event. If this is the case, the merger model predicts the formation of a new, metal rich population of GCs during the merger, which then reddens and fades with time. Post-merger ellipticals would thus be expected to possess a ‘new’ metal rich population and an old, metal poor population originating in the progenitor disc galaxies. Fig. 8 reveals evidence for two GC populations in NGC 1700. We next use the colours and magnitudes of these two populations to determine their age and metallicity properties. The blue peak of our $`BI`$ histogram is roughly consistent with the peak of the $`BI`$ distribution of Galactic GCs (see Fig. 8). The Milky Way is likely to be a typical example of a progenitor in a present day gas-rich merger, thus indicating that we may assume that our blue peak in $`BI`$ is due to an old, metal poor population of GCs. Here we use the stellar population models of both ? and ? to predict the difference in magnitude and colour for a population of young clusters, relative to an old, metal poor population with an age of 15 Gyr and $`[\mathrm{𝐹𝑒}/H]=1.5`$. Fig. 9 shows the predicted differences in $`\mathrm{\Delta }(BI)`$ vs. $`\mathrm{\Delta }B`$ and $`\mathrm{\Delta }I`$ based on the models of ?. Fig. 10 shows the equivalent colour and magnitude changes based on the models of ?. In each case, three different metallicity tracks are shown. From Fig. 8, $`\mathrm{\Delta }(BI)=0.44\pm 0.07`$. The $`\mathrm{\Delta }B`$ and $`\mathrm{\Delta }I`$ values correspond to the magnitude offset in the $`B`$ and $`I`$ bands between the young and old populations. These offsets were estimated from the cumulative $`B`$ and $`I`$ band GC luminosity functions for both the blue and red populations separately. These luminosity functions are shown in Fig. 11. To reduce the effect of photometric errors we considered only those GCs at least 0.5 mag brighter than the limiting magnitude in each band. The data has been normalised to the faint limit of our luminosity functions. The median magnitudes of the red and blue GCs considered were then calculated to allow us to determine the separation of the blue and red populations. The magnitude differences between the red and blue populations were determined to be $`\mathrm{\Delta }B=0.0\pm 0.2`$ and $`\mathrm{\Delta }I=0.4\pm 0.2`$ in the $`B`$ and $`I`$ bands respectively. The observed values with their associated errors are indicated by a cross on Fig. 9 and Fig. 10. Reference to the Bruzual & Charlot models shows that the red population is consistent with a young, metal rich population of globular clusters, i.e. with an age of 2.5–5.0 Gyr and super-solar metallicity ($`[\mathrm{𝐹𝑒}/H]+0.1`$ to +0.6). A plot of $`\mathrm{\Delta }V`$ vs. $`\mathrm{\Delta }(BI)`$, where $`\mathrm{\Delta }V=0.1\pm 0.2`$ predicts a similar age and metallicity. It is interesting to note that the Worthey models predict significantly different values for the age and metallicity of the red GC population in the sense that they are older and more metal poor. From the Worthey models the age of the red GCs is 5.0–8.0 Gyr and $`[\mathrm{𝐹𝑒}/H]0.2`$. It thus appears that ages derived from GC colours and magnitudes are highly model dependent. However both models suggest that the red GCs are both younger and more metal rich than the blue population. We may suspect that $`[\mathrm{𝐹𝑒}/H]`$ should be at least solar metallicity. In the merger model, the young population is thought to have formed from the relatively enriched gas in the spiral discs, in contrast to the old population which formed from metal poor gas. Moreover, the line strengths plotted in Fig. 13 (see Section 4.3 for details) imply for both models that the stellar metallicity of the last major starburst is $`[\mathrm{𝐹𝑒}/H]+0.5`$. If the young GCs formed from the same gas, we would expect them to be metal rich as well, which in turn would favour the younger age of the Bruzual & Charlot models. However, if the young GCs were formed at an early stage of a merger-induced starburst then they could in principle be more metal poor than the present young stellar population (?). To summarise, if we assume that the blue GCs are an old metal poor population, the red GCs are consistent with an age of 2.5–5.0 Gyr, assuming they have a metallicity of $`0.0[\mathrm{𝐹𝑒}/H]+0.5`$. In the case that the red GCs are relatively metal poor, as suggested by the Worthey model tracks in Fig. 10 then a slightly older (5.0–8.0 Gyr) age is indicated. Infrared photometry or good spectra would help to resolve the age and metallicity independently. From Section 4.1 we found evidence for bimodality in $`VI`$ with peaks at $`VI=0.85\pm 0.05`$ and $`VI=1.15\pm 0.05`$. A GC population with an age of 15 Gyr, $`[\mathrm{𝐹𝑒}/H]=1.5`$ and $`VI=0.85`$ will have $`BI1.5`$. This is consistent with the blue population seen in our $`BI`$ colour histogram. Similarly, a population with an age of 3 Gyr, $`[\mathrm{𝐹𝑒}/H]=+0.5`$ and $`VI=1.15`$ is expected to have $`BI1.9`$, i.e. consistent with the red peak of our $`BI`$ distribution. It thus appears that the colours of the red and blue peaks in our $`BI`$ histogram are consistent with the the peaks hinted at in the HST $`VI`$ colour distribution. The bimodality seen in independent data sets confirms the presence of two GC populations in NGC 1700. We show in Fig. 12 the mean $`BI`$ colour of our GCs vs. galactocentric radius. The data are radially binned into 10 annuli, each of width 5.5 kpc (22 arcsec). The GC system of NGC 1700 does not show any significant correlation between colour and galactocentric distance. In particular the red (i.e. young) GCs are not preferentially concentrated towards the centre of the galaxy, as might have been expected from the merger origin for new GCs (??). A linear least-square fit to the radial colour bins gives a gradient consistent with zero. We do not detect any significant radial colour trend in the HST sample of GCs either, agreeing with the results of ?. This lack of an obvious radial trend is potentially a problem for the merger origin interpretation of the GCs, and deserves further consideration in any future studies. It is also interesting to note that the GC systems of several other ellipticals do show radial colour gradients in the sense that the red GCs are more centrally concentrated than the blue (see e.g. Geisler, Lee & Kim 1996; Forbes, Brodie & Grillmair 1997). ### 4.3 Other age estimates Until recently, direct age dating of the stars in an elliptical was near-impossible due to the well known degeneracy of age–metallicity effects in old stellar populations. Spectral synthesis methods have now been developed (e.g. ??) that show that combinations of certain spectral line indices can efficiently disentangle the effects of age and metallicity for young to intermediate age stellar populations. As well as stellar age dating, a number of other methods are available for NGC 1700. These estimates are shown for comparison in Table 1. We show a mean spectroscopic age based on the line strength data of Fisher, Franx & Illingworth (1996) and ?, the age derived from the globular clusters by ?, the ‘fine structure age’ of ?, the age from stellar dynamics (Statler, Smecker–Hane & Cecil 1996) and the age derived from the galaxy’s deviation from the Fundamental Plane (see Forbes, Ponman & Brown 1998). #### 4.3.1 Ages from stellar population synthesis models The spectroscopic age estimates in Table 1 are based on comparisons of observed absorption line strengths with single-burst stellar population models of ? and ?. From these model grids, observed line strengths can be used to give estimates of the age and metallicity of the stellar component of a galaxy. We have used the $`H\beta `$ and $`[\mathrm{𝑀𝑔𝐹𝑒}]`$ line strengths from two sources, i.e. ? and ?. Fig. 13 shows the $`[\mathrm{𝑀𝑔𝐹𝑒}]H\beta `$ model grids with the observed values for NGC 1700 shown. In both panels, the solid circle indicates the observed line strengths from ?, whereas the open circle shows the line strengths from ?. The apertures used in each case are $`r_\mathrm{e}/10`$ and $`r_\mathrm{e}/8`$ respectively because we are primarily concerned with the age of the central starburst. Using the Worthey models, the age derived from the Fisher et al. data is $`2.7\pm 0.4`$ Gyr and the age from the Gonzalez data is $`2.3\pm 0.3`$ Gyr. In both cases, the metallicity, $`[\mathrm{𝐹𝑒}/H]`$ is $`+0.5`$. Similar results are obtained by using the Bruzual & Charlot models. In this case the Fisher et al. data gives an estimated age of $`3.5\pm 0.5`$ Gyr and the data of Gonzalez suggests an age of $`2.9\pm 0.3`$. Combining these results gives an average of $`2.9\pm 0.3`$ Gyrs. Again both sets of data suggest a metallicity of $`[\mathrm{𝐹𝑒}/H]+0.5`$. Note that this method measures a ‘luminosity weighted mean age’ of the central stellar population and is thus likely to be dominated by the young stellar population associated with the last episode of star formation. This may in turn be the result of a merger induced starburst. #### 4.3.2 Ages from globular clusters ? employed a method similar to that used in Section 4.2.2 to derive GC ages. This method relies on the observed colour difference between two detected populations of GCs. ? detected only one population (at $`VI=1.05`$). Their age of ‘$`4\pm 2`$ Gyr ?’ was derived from the non-detection of two distinct populations, though it is consistent with the other age estimates discussed herein. The non-detection appears to have arisen due to the small colour separation of the GC populations. Using the same data but restricting the sample to bright GCs, we have detected two separate populations as shown in Fig. 4. By including the whole sample (i.e. down to $`V=26.5`$) the two peaks merge into a single distribution. #### 4.3.3 Fine structure ages Another method of estimating galaxy age was developed by ?. They defined a fine structure parameter, $`\mathrm{\Sigma }`$, based on the amount of optical ‘fine structure’ present in a galaxy. This included a measure of the maximum boxiness of the galaxy isophotes, the number and strength of shells and tails and the presence or absence of ‘X-structure’. The values of $`\mathrm{\Sigma }`$ for the galaxies in their sample ranged from 0 for galaxies with no fine structure to 7.6 indicating the largest amount of fine structure observed. NGC 1700 was assigned a value of 3.70. They found that $`\mathrm{\Sigma }`$ correlated with a galaxy’s residual from the mean colour–magnitude relation. Using this fact and relating galaxy colours to ages via a star formation model, they estimated the time since the merger event. We quote their most representative age of 6.0 Gyr with an error of $`\pm 2.3`$ Gyr in Table 1. We note however , that the models used in ? assumed a solar metallicity starburst. If the stellar population is super-solar, as suggested by Fig. 13, then the ? age is expected to be a slight over-estimate. #### 4.3.4 Dynamical considerations The age estimate of ? comes from measurements of the stellar velocity field of NGC 1700. They found that within $`2.5r_\mathrm{e}`$ (8.8 kpc) the galaxy is kinematically well mixed. This constrains the time since the last major merger event to be $`2.7`$ Gyr to allow sufficient time for phase mixing and differential precession. Their observations also define an upper limit for the time since the merger. The asymmetric photometric and kinematic signatures at larger radii preclude a merger age greater than 5.3 Gyr, otherwise these features would have relaxed and disappeared. Thus we quote an age of 2.7–5.3 Gyr. The same paper argues that the counter rotating core and boxy features (the latter of which we attribute to the tidal tail-like structures) could not have been created by the same merger event and that the observed form of NGC 1700 must have arisen from the merger of at least three separate stellar systems. However, it was not clear whether these events occurred sequentially or simultaneously. If the two tail-like structures are indeed genuine tidal tails this would suggest that the history of NGC 1700 has included at least one major merger event involving two approximately equal mass disc galaxies. If the kinematically distinct core (KDC) was formed prior to this event it would have probably been disrupted during the ensuing violent relaxation processes. If the KDC was indeed formed by a separate process it must have resulted from a subsequent minor merger (e.g. ?) or an interaction (see ?) some time after the major merger event that created NGC 1700 and its tails. Neither of these suggestions are very appealing. Alternatively, the tidal structures seen in Fig. 5 could be interpreted as plumes which have arisen from the infall of a small disc galaxy into a pre-existing elliptical. #### 4.3.5 Scatter from the Fundamental Plane A recent study by ? showed that a galaxy’s deviation from the Fundamental Plane (FP) correlated with its age, albeit with large scatter. The FP residual is defined as $`R(\sigma _0,M_\mathrm{B},\mu _\mathrm{e})=2\mathrm{log}(\sigma _0)+0.286M_\mathrm{B}+0.2\mu _\mathrm{e}3.101`$ (?). Young ellipticals fall below the FP (i.e. have negative residuals) and evolve towards it until they lie on the FP at an age of about 10 Gyr. Older ellipticals tend to lie above the FP. It is thus possible to use the FP residual to approximately age date an elliptical galaxy. The FP residual for NGC 1700 is $`R(\sigma _0,M_\mathrm{B},\mu _\mathrm{e})=0.37`$. This corresponds to an age of $`1.2\pm 2.0`$ Gyr based on the fit in ?. This is slightly younger but comparable to the other age estimates discussed in this section. #### 4.3.6 Comparison of various age estimates Before directly comparing the different age estimates, one should bear in mind that the dating methods may be measuring different time-scales and may have large errors associated with them. As mentioned above, the stellar spectroscopy methods measure the central ‘luminosity weighted average age’ which means they are dominated by the last burst of star formation, although the old stellar population also contributes. Thus the true age of the starburst may be slightly less than the spectroscopic age. The young GCs were possibly formed in the same star formation event as the galaxy starburst, and as such should give a similar age as the stellar spectroscopy. The age estimate from the ? trend is also based on a star formation time-scale. The stellar dynamics, fine structure and tail fraction ages estimate the time since the last merger event. The ‘dynamical/structural’ estimates are reasonably consistent with the ‘starburst’ ages, indicating a relatively young age for NGC 1700, and suggesting that the merger was a gaseous one. ? have shown that in a merger between two spirals with bulges the main starburst occurs at the time of nuclear coalescence while the tails form $`0.5`$ Gyr earlier. If this is the case, we might expect the ‘dynamical/structural’ age estimates to be slightly higher than the ‘starburst’ ones. We have decided to adopt an age for NGC 1700 of $`3.0\pm 1.0`$ Gyr as our best estimate. This may also correspond to the time since the nuclei of the progenitors merged. ### 4.4 Globular cluster spatial properties and specific frequency #### 4.4.1 Spatial distribution Given the field–of–view, our Keck observations are ideal for defining the outer reaches of the GC system. We have calculated the surface density (SD) profile for the GC system in the Keck images within 9 annuli centred on the galaxy. For objects in the corner of the CCD, the SD was calculated by taking into account the area of the annulus ‘missing’ off the edges of the chip. This method allowed us to calculate the density out to a radius of 232 arcsec ($`60`$ kpc). The resulting surface density profile is shown in Fig. 14. The error bars simply reflect the Poisson errors on the number of GCs in each bin. At large radii the surface density decreases like a power-law with radius (open squares). At radii $`<140`$ arcsec however, there appears to be a significant deficit in the number of GCs detected in our Keck images (shown by open circles). The most likely cause for this is the fact that GCs at small radii are superimposed on the bright body of the galaxy which also possesses a steep radial gradient at these distances. Although an elliptical model of the galaxy was subtracted from each of the initial images, this effect seems to have caused a reduction in our ability to detect GCs as our HST surface density continues to rise towards the centre at these radii as shown in Fig. 15. In addition, we were unable to detect GCs in the Keck images within the central 33 arcsec due to the large saturated region at the centre of the $`I`$ band image. To quantify the outer SD profile of the GCs detected in our Keck images, we fitted the outer-most points in Fig. 14 with a function of the form $`\rho =\rho _0r^\alpha `$. This fit is shown by the dashed line in Fig. 14. We find that $`\rho _0=0.5\pm 3.3`$ and $`\alpha =1.07\pm 0.25`$. Background contamination of the sample would add a constant SD level to the profile. Although this contamination is likely to be small, ideally it would be taken into consideration when computing the fit. However, our data are not sufficient to calculate the level of contamination and the calculated slope of the profile should be treated with caution. The solid line in Fig. 14 represents the stellar profile of the galaxy. The galaxy profile used is not in the usual units of surface brightness but has been converted to log$`(intensity)`$ and arbitrarily shifted in the Y direction to allow simple comparison of the slopes. We measure a slope of $`1.9\pm 0.1`$ for the galaxy profile using the measured intensities at all radii. It thus appears that the GC profile is flatter than the underlying stellar profile. In order to investigate the SD profile within the central regions of the galaxy, we use the HST images. This profile is shown in Fig. 15 after correction for the areas of the annular bins missing from the WFPC2 field–of–view. Each annulus contains $`25`$ GCs with the exception of the two inner bins which contain $`10`$ GCs. We find that the HST surface density follows a power-law profile exterior to $`r11`$ arcsec (i.e. the outer 6 data points). We again fitted these points with a function of the form $`\rho =\rho _0r^\alpha `$ and find $`\alpha =0.90\pm 0.15`$ and $`\rho _0=0.27\pm 0.19`$ GC arcsec<sup>-2</sup>. The slope of this fit is within the range of values measured for other ellipticals by other authors (e.g. ?) and also consistent with the slope calculated from the Keck sample. Also shown in Fig. 15 is the stellar profile of the galaxy. This has a slope of $`1.9\pm 0.1`$. The GC profile is significantly flatter than that of the galaxy, which is often the case in other systems (e.g. Grillmair, Pritchet & van den Bergh 1986; ????). Within $`11`$ arcsec there is evidence for a flattening of the profile indicating the presence of a ‘core region’. The core radius of the NGC 1700 GC distribution was also measured by ? to be $`2.7\pm 1.6`$ kpc corresponding to $`11`$ arcsec, although only 39 GCs were detected. The presence of a core region is also commonly seen in other ellipticals and ? define a relationship between GC system core radius and parent galaxy luminosity. This relation is in the sense that more luminous ellipticals have more extended GC systems. An elliptical with the luminosity of NGC 1700 (i.e. $`M_\mathrm{V}=22.47`$) is expected to have a core radius of $`3`$ kpc, corresponding to $`12`$ arcsec, although there is a large uncertainty on this value (?). Our data are thus consistent with the expectation from other ellipticals and with the previous direct measurement of the NGC 1700 GC system core radius. #### 4.4.2 Total number of globular clusters We next make an estimate of the total number, $`N_\mathrm{T}`$, of GCs possessed by NGC 1700. In order to estimate $`N_\mathrm{T}`$ we integrate under the surface density profile out to some limiting radius. We will assume a flat profile interior to our second data point ($`r=11.3`$ arcsec) with $`\mathrm{log}(\mathrm{SD})=1.52`$ (as shown in Fig. 15), although this assumption has little effect on the final calculation. Outside this radius we assume a power-law profile. The faint magnitude limit for our Keck data ($`V25`$) is much brighter than the expected peak magnitude of a standard GC luminosity function at the distance of NGC 1700. Thus the correction for incompleteness would be very large. In addition, the Keck data has a higher contamination rate than HST which is difficult to account for without a corresponding ‘blank sky’ image. We therefore integrate under the power-law fitted to the HST points only, for which we are confident there is very little contamination, and we expect a much smaller incompleteness correction. We use the Keck data to help us to define a reasonable outer limiting radius. The dominant sources of error in determining $`N_\mathrm{T}`$ are the choice of the outer limiting radius and the correction for incompleteness at faint magnitudes. Fig. 14 indicates that the surface density declines out to a radius of $`210`$ arcsec outside of which there is a large drop in surface density. This suggests that the limit of the GC distribution has been reached. We thus chose to integrate out to 210 arcsec, which corresponds to a galactocentric distance of 50 kpc, in order to calculate the total number of GCs within this radius. This limiting radius is comparable to the extents of the GC systems seen in other ellipticals (see e.g ?). The resulting number of GCs was then corrected to account for the lack of completeness in the HST sample at faint magnitudes. Based on our limiting magnitude we estimate this correction factor to lie in the range of 2–3, thus adopting a correction factor of 2.5, with an uncertainty of $`\pm 0.5`$. We estimate that we are detecting roughly $`(40\pm 8)`$ per cent of the total GC system within the limits of the HST field–of–view. We multiplied the number of GCs calculated from the SD profile by 2.5 to obtain a completeness corrected value for $`N_\mathrm{T}`$ within 50 kpc of $`N_\mathrm{T}=1320\pm 270`$. Our value is greater than that found by ? who estimated a total number of $`528\pm 48`$ GCs. However they considered only those GCs out to the radius of the HST field–of–view, whereas our Keck data indicates that the GCs are still present out to $`200`$ arcsec. #### 4.4.3 Specific frequency The globular cluster specific frequency, $`S_N`$, is a useful quantity as it can provide valuable constraints on the formation mechanisms of the GC system as well as providing information on the formation history of the parent galaxy and potentially the nature of the progenitor spirals. From $`N_\mathrm{T}`$, we calculate the total GC specific frequency $`S_N(\mathrm{total})`$ within 50 kpc from the galaxy centre, from the following relation (?) $$S_N(\mathrm{total})=N_\mathrm{T}\times 10^{0.4(M_\mathrm{V}+15)}$$ where $`M_\mathrm{V}`$ is the absolute magnitude of the stars associated with the population of GCs in question. In the case of $`S_N(\mathrm{total})`$, the required luminosity is the absolute magnitude of the galaxy, i.e. $`M_\mathrm{V}=22.47`$. We will assume a 20 per cent error in the distance to NGC 1700, corresponding to an error in $`M_\mathrm{V}`$ of 0.4 mag. This yields for $`N_\mathrm{T}=1320\pm 270`$, a total specific frequency of $`S_N(\mathrm{total})=1.4_{0.6}^{+1.0}`$. Elliptical galaxies have a wide range of GC specific frequencies. The mean value for ellipticals in the compilation of ? is $`S_N=5.1\pm 0.6`$, though there are a couple of cases of ellipticals with $`S_N`$ values less than 1.0. Although possessing large errors, our $`S_N(\mathrm{total})`$ is relatively low compared to most typical ellipticals and is more consistent with disc galaxies. It is clear from the age estimates discussed herein that NGC 1700 is a relatively young galaxy. Will the GC population eventually come to resemble those seen around typical old ellipticals ? To address this question we use the Worthey models to predict the total specific frequency once sufficient time has passed for NGC 1700 to have an age comparable to those of typical elliptical galaxies today. Old, present day ellipticals have stellar population ages of 10–15 Gyr so in order to make this comparison we compute the luminosity of a young starburst component after it has aged by 10 Gyr. As the mass of the starburst is uncertain, we consider two cases: Firstly we assume a starburst that represents 10 per cent of the galaxy by mass and is embedded in a 10 Gyr population making up the remaining 90 per cent of the galactic stellar mass. After 10 Gyr, the models indicate that the young stellar population will fade by $`\mathrm{\Delta }V=1.48`$ mag (for an $`[\mathrm{𝐹𝑒}/H]=+0.5`$ population). As this population constitutes only 10 per cent of the galaxy mass however, the global fading is only $`\mathrm{\Delta }V=0.18`$ mag. This results in a predicted $`S_N(\mathrm{total})`$ of $`1.6`$, which is still lower than that expected for most typical present day ellipticals and suggests that NGC 1700 may form a relatively ‘globular cluster poor’ elliptical galaxy. If the starburst fraction is less than $`10`$ per cent (as might be expected in an elliptical plus spiral merger), then the global fading is further reduced. Secondly, we consider a starburst that constitutes 50 per cent of the galaxy mass. The remaining 50 per cent is made up of the old stellar population. After 10 Gyr, the fading of the young stellar population results in a global fading of $`\mathrm{\Delta }V=0.81`$ mag. The $`S_N(\mathrm{total})`$ predicted after 10 Gyr is $`2.8`$. This is more consistent with elliptical galaxy specific frequencies, though still lower than the mean $`S_N`$ value calculated from ?. The merger model of GC formation by ? states that the number of new globular clusters formed in a merger will be proportional to the available gas mass, and that most ‘normal’ old ellipticals were formed from progenitors which were relatively gas–rich compared to present day spirals. It is interesting to note therefore that if one accepts the merger hypothesis and the fact that NGC 1700 is a relatively young elliptical, its GC specific frequency may be expected to be relatively low. Note, however that our ‘predicted’ values of $`S_N`$ should be considered as very approximate as there is not only an inherent difficulty in estimating our sample completeness (and hence the number of GCs), but also a large uncertainty associated with the assumptions made regarding the future evolution of the galaxy luminosity. ## 5 Conclusions We have presented results from new $`B`$, $`V`$ and $`I`$ band imaging of NGC 1700, taken using the Keck telescope and reanalysed previous HST imaging. From the morphology of the galaxy and photometry of its globular cluster system, we have derived new estimates of the galaxy age. Subtraction of an elliptical model from the Keck images revealed the presence of two symmetric tidal tail-like structures extending $`40`$ kpc to the North-West and South-East of the galaxy. The presence of tidal tails is thought to be a classic signature that a galaxy has undergone a merger event, involving the collision of two spiral galaxies during the last few Gyr. If the observed tidal features are indeed genuine tidal tails, then this would suggest that NGC 1700 has undergone a major merger event during its recent history. If they are merely plumes or other fine structure, then the situation of a disc galaxy merging with an elliptical is also possible. Based on the fraction of galaxy light contained within these tails we have estimated their age at $`3.2\pm 1.5`$ Gyr. We have also shown that NGC 1700 possesses boxy isophotes for radii of $`7.5<r<20.0`$ kpc, and that this boxiness is largely caused by the presence of the tail structures. From the HST imaging, we detected 146 globular clusters. Of these, 34 were in common with our GC candidates detected in the Keck data. We then used these common GCs to refine our Keck GC selection. After further restricting our sample in magnitude and colour we obtained a final Keck object list of 312 GCs. These show a bimodal colour distribution with peaks at $`BI=1.54`$ and 1.98. The colour of the blue population is consistent with that of Galactic GCs. Assuming that the blue population is indeed an old metal poor system and measuring the offset in magnitude and colour between the two populations, we find that the red GCs are younger and more metal rich. The ? stellar population models suggest that the red population is 2.5–5.0 Gyr old and has super-solar metallicity. The equivalent models of ? predict a significantly older age (5.0–8.0 Gyr) and a lower metallicity ($`[\mathrm{𝐹𝑒}/H]0.2`$). The bimodality in $`BI`$ is supported by a hint of bimodality seen in the $`VI`$ colours from HST data. We discuss other age estimates from the literature based on fine structure, globular cluster colours, Fundamental Plane residuals, stellar dynamics, and spectroscopic line strengths. We find that, although possessing significant errors, the various age estimates generally indicate a young age for NGC 1700 of about $`3.0\pm 1.0`$ Gyr. The fact that the ‘dynamical/structural’ and ‘star formation’ ages are similar suggests that NGC 1700 has undergone an episode of enhanced star formation triggered by a merger event, which may have created the galaxy itself. The surface density profile of GCs reveals a flatter (i.e. more extended) profile than the underlying galaxy starlight. The total number of GCs and present total specific frequency are estimated to be $`N_\mathrm{T}=1320\pm 270`$ and $`S_N(\mathrm{total})=1.4_{0.6}^{+1.0}`$ respectively. The value for $`S_N(\mathrm{total})`$, even considering the large errors, is quite low compared to the majority of typical ellipticals and is more consistent with spiral galaxies. We predict that after 10 Gyr the total $`S_N`$ will have increased to a maximum of $`2.8`$ if NGC 1700 has undergone a starburst that constitutes 50 per cent of its stellar mass. This is consistent with low $`S_N`$ ellipticals. If the starburst constitutes only 10 per cent of the total stellar mass, the predicted $`S_N`$ is somewhat lower at $`S_N1.6`$. This suggests that NGC 1700 will form a relatively ‘globular cluster poor’ elliptical galaxy once it reaches a comparable age to typical ‘old’ ellipticals. ## Acknowledgments We thank Trevor Ponman, Alejandro Terlevich and Edward Lloyd–Davies for help and useful discussions. We also thank the referee, Keith Ashman for his helpful comments. Some of the data presented herein were obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. This work was also based on observations with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract No. NAS5-26555. Part of this research was funded by NATO Collaborative Research grant CRG 971552.
warning/0004/astro-ph0004280.html	ar5iv	text	# Observations of the supernova remnant W28 at TeV energies ## 1 Introduction Supernova remnants (SNRs) have long been thought to be the dominant source of cosmic rays (CR) at energies below 100 TeV (for a review see e.g. Blandford & Eichler Blandford:1 (1987)). SNR, via the diffusive shock process, are able to accelerate electrons and hadrons and meet the energetics of the observed cosmic rays. The TeV gamma-ray flux predicted from SNR is the most accessible tracer of CR acceleration and its detection would be convincing evidence for the SNR origin of galactic CR. Models of the TeV gamma-ray emission from SNR predict distinct spectral features, according to the hadronic and/or electronic nature of the parent CR accounting for the gamma-ray flux (see Drury et al. Drury:1 (1994), Naito and Takahara Naito:1 (1994), Baring et al. Baring:1 (1999), and references therein for a summary). Ground-based surveys of SNR at gamma-ray energies (TeV to PeV) have been carried out on several promising northern hemisphere candidates (e.g. IC443, Tycho’s SNR, W51, W44, G78.2$`+`$2.1). The Whipple (Buckley et al. Buckley:1 (1998), Lessard et al. Lessard:1 (1999)), HEGRA (Hess et al. Hess:1 (1997) at TeV energies, and Prosch et al. Prosch:1 (1996) at multi-TeV energies), CAT (Goret et al. Goret:1 (1999)) and CYGNUS (Allen et al. Allen:1 (1995)) groups have reported upper limits. Recently however, the HEGRA has seen marginal evidence for TeV gamma-rays from the young SNR Cas-A, after deep observation (Pühlhofer et al. Puhlhofer:1 (1999)). In the southern hemisphere, the CANGAROO has reported the detection of TeV gamma-rays from SNR SN1006 (Tanimori et al. Tanimori:1 (1998)) and SNR RX J1713.7$``$3946 (Muraishi et al. Muraishi:1 (2000)), and if confirmed, will be strong evidence in favour of the production of cosmic rays electrons in SNR. W28 (also SNR G6.4$``$0.1 from Green Green:1 (1998)) is considered an archetypal composite (mixed or M-type) supernova remnant, characterised by a centrally filled X-ray and shell-like radio morphology (Rho & Petre Rho:1 (1998), Long et al. Long:1 (1991)). The ROSAT X-ray emission appears best explained by a thermal model (Rho et al. Rho:2 (1996)) although Tomida et al. (Tomida:1 (1998)) from the analysis of ASCA data, has suggested the presence of a weak a non-thermal component in the south west region. The limb-brightened radio emission (20, 6 & 2cm) shows a synchrotron spectrum of varying spectral index (Andrews et al. Andrews:1 (1983)). A radio point source at $`l=6.6^{},b=0.16^{}`$ (G6.6$``$0.1) is defined (Altenhoff et al. Altenhoff:1 (1978), Andrews et al. Andrews:1 (1983)), hereafter referred as A83 in this paper. A glitching radio pulsar, PSR J1801$``$23 (PSR B1758$`23`$, $`P=`$416ms, $`\dot{P}=113\times 10^{15}`$ ss<sup>-1</sup>), lies at the northern radio edge (Kaspi et al. Kaspi:1 (1993)). An upper limit to this pulsar’s characteristic age is estimated at 58 000 years, and it’s spin-down luminosity ($`\dot{E}6.2\times 10^{34}`$ erg s<sup>-1</sup>) is at the lower edge of luminosity values when compared to the known gamma-ray (EGRET & COMPTEL) pulsars. The age of W28 is estimated (Kaspi et al. Kaspi:1 (1993)) in the range 35 000 to 150 000 years, with upper and lower limits taken from the assumptions that W28 is currently in either the radiative or Sedov phases of expansion. According to Kaspi et al. Kaspi:1 (1993), the distance of PSR J1801$``$23 (9 to 16.5 kpc) derived from it’s dispersion measure (DM) appears inconsistent with that derived for the remnant. Estimates for the remnant distance are set at 1.8 kpc (Goudis Goudis:1 (1976) $`\mathrm{\Sigma }`$-D relation) and 3.3 kpc respectively (Lozinskaya Lozinskaya:1 (1981), from mean optical velocities), indicating that the pulsar/W28 association is possibly a line-of-sight coincidence. However, Frail et al. (Frail:1 (1993)) have noted the large uncertainty in using the DM as a distance estimate for this pulsar due to the high concentration of ionised material in the line of sight, and conclude there is sufficient evidence for the pulsar/remnant association. The unidentified EGRET source 3EG J1800$``$2338 (95% error circle 0.32 radius) (Hartman et al. Hartman:1 (1999)), listed as 2EG J1801$``$2312 in the second EGRET catalogue (Thompson et al. Thompson:1 (1996)), lies on the edge of the radio shell and was thought to be associated with the remnant (Esposito et al. Esposito:1 (1996), Zhang & Cheng Zhang:1 (1998)). 3EG J1800$``$2338 has a relatively hard spectral index (Hartman et al. Hartman:1 (1999)) with no apparent sign of a turnover at 1 GeV (Merck et al. Merck:1 (1996)). Lamb & Macomb (Lamb:1 (1997)) also point out that 3EG J1800$``$2338 is visible above 1 GeV at 5.4$`\sigma `$ significance, and is centred very close to the A83 radio position. The 3EG position of the EGRET source is displaced by about 0.5 relative to the 2EG position, yet still lies comfortably within the SNR radio shell, and remains a strong example of an EGRET source/SNR association (Romero et al. Romero:1 (1999)). The 3EG error circle however, now excludes PSR J1801$``$23 and the molecular clouds. W28 lies in a complex region of the galactic plane with many HII regions and dense molecular clouds (Wootten Wootten:1 (1981)) contributing to the ISM surrounding the SNR. Over forty OH (1720 MHz) maser emission sites are concentrated at the eastern and northern edges of the SNR (Claussen et al. Claussen:1 (1997)), along the SNR and molecular cloud interface. The distribution of shocked and unshocked gas in this region is also consistent with the idea of the SNR shock passing through the cloud (Arikawa et al. Arikawa:1 (1999)). OH maser emission (1720 MHz) is considered a strong indicator of collisional pumping with matter densities $`10^5`$ cm<sup>-3</sup> (Claussen et al. Claussen:1 (1997, 1999)). Enhanced levels of TeV $`\gamma `$-ray emission via the decay of neutral pions may be expected from such areas associated with the masers and molecular cloud (Aharonian et al. Aharonian:1 (1994)). Fig. 1 indicates the sites of interest in relation to the radio continuum emission (327 MHz). The presence of these interesting objects make W28 a prime southern hemisphere candidate for study at TeV gamma-ray energies. We report here on the comprehensive analysis of data taken in 1994 and 1995 with the CANGAROO 3.8 metre telescope. This work follows analysis of data taken in 1992 (Kifune et al. Kifune:1 (1993)) in which weak evidence for a gamma-ray signal was reported. At that time, only ON source data were collected, making an estimation of the background rate difficult. Mori et al. (Mori:1 (1995)) reported briefly on an analysis of 1994 data centred on PSR J1801$`23`$, in which a $`\pm 0.7^{}`$ field of view was searched. Both ON and OFF source data were collected and no evidence for TeV $`\gamma `$-ray emission was seen from various point-like sources including the pulsar and both radio and X-ray maxima. The 1995 data were centred on the radio position A83, located $`0.3^{}`$ away from PSR J1801$`23`$. A search for point-like and diffuse sources of TeV emission was carried out on the 1994 and 1995 datasets out to $`\pm 1^{}`$ from the tracking centre of each dataset, using an extended source analysis. We have used an improved set of cuts to those used in the analysis of data taken on the Vela Pulsar/Nebula (Yoshikoshi et al. Yoshikoshi:1 (1997)). These cuts were designed to minimise the loss of gamma-ray sensitivity for off-axis sources and in particular maintain reliable statistics over the search region. ## 2 The CANGAROO 3.8m telescope The CANGAROO 3.8 m imaging telescope is located near Woomera, Australia ($`137^{}47^{}`$E, $`31^{}06^{}`$S, 160m a.s.l.). The 3.8 metre diameter mirror, of focal length 3.8 metres, is used to image the Čerenkov emission from gamma-ray and cosmic-ray induced extensive air showers (EAS) onto a high resolution multi-phototube camera. The camera consisted of 224 photomultiplier tubes in 1994. In April 1995 an extra 32 tubes were added to the corners of the camera bringing the total to 256 tubes in a 16$`\times `$16 square grid arrangement. However the extra corner tubes have not been used in the analysis of the 1995 W28 dataset in order to retain consistent imaging properties with the 1994 dataset. Each camera tube is a Hamamatsu R2248 with a photocathode size 0.12$`{}_{}{}^{}\times 0.12^{}`$ on a side, and the total field of view of the camera is 2.9. An event trigger is registered when the summed output of triggered tubes exceeds a preset threshold, denoted HSUM. Images with a minimum of between 3$``$5 tubes, depending on the image’s compactness, trigger the telescope. A vertical event rate due to cosmic rays of $``$2 Hz is achieved. Monte Carlo simulations of the telescope indicate an gamma-ray energy threshold of $``$1.5 TeV at the vertical (Roberts et al. Roberts:1 (1998)), where the energy threshold is defined as that representing the half-maximum of the differential distribution of triggered energies. Tracking calibration is performed by monitoring the paths of bright (visual mag. 3–6) stars in the field of view, providing an absolute tracking accuracy of $`0.02^{}`$ (Yoshikoshi Yoshikoshi:2 (1996)). A more detailed technical description of the telescope appears in Hara et al. (Hara:1 (1993)). Data are recorded on clear moonless nights. An ON source run is generally followed by an OFF source run displaced in right ascension to provide a background run of matching zenith and azimuth angle distributions. However, since small sections of data from both observation seasons were removed due to cloud effects, a normalisation, described later, was used in estimating the statistical significance of any ON source excess. Pulse charge (ADC) and timing (TDC) information for each tube is recorded for each event. Calibration of the ADCs and TDCs is achieved by recording events triggered with a blue LED flasher before each observation. Tube signals are accepted as part of an image if they meet a number of criteria: 1. The TDC value of a tube must lie between $`\pm 37.5`$ns, referenced against the event trigger time. The event trigger time is registered when the HSUM threshold is met. 2. The ADC value must be greater than one standard deviation above the RMS noise (comprising skynoise and electronic noise) for that tube. 3. The tube must not be isolated. An isolated tube is one which is not adjacent to any other accepted tube. 4. The tube must not have an outlying relative gain value. Tubes with relative gains outside the range 1.0$`\pm `$0.3 contribute significantly to trigger differences across the camera face. This factor is particularly important when comparing regions over the entire camera face. The telescope is an altitude-azimuth type, which introduces a rotation of the camera relative to the sky about the tracking position during data collection. A ‘de-rotation’ is applied to the tube positions to account for this effect, and is necessary when considering off-axis sources. The images are parameterised according to the moment-based method of Hillas (Hillas:1 (1985)). Pre-processing steps designed to minimise the effects of electronic interference are described below. The camera is divided into groups of eight tubes which share common high voltage and other circuitry, and a special cut, box (Yoshikoshi Yoshikoshi:2 (1996)), is designed to remove images arising from electronic contamination, and are concentrated in only one or two tube boxes, This box cut, in combination with a total ADC sum (adc) cut rejecting events with fewer than 200 ADC counts, is very effective at removing such artifacts. Monte Carlo simulations show that the box cut does not reduce the power of the image cuts, and that the optimum adc cut lies at $``$200 ADC counts. Mirror degradation resulted in an event rate drop by about a factor of two from 1994 to 1995, indicating a higher energy threshold for the 1995 dataset. The results quoted in this work are normalised to a 1.5 TeV threshold for gamma-rays, using different raw triggering efficiencies for gamma-rays (Sect. 3), which take into account the increase in trigger threshold between 1994 and 1995. In addition, only events with width $`>=`$ 0.01 and with the number of triggered tubes, ntubes$`4`$ are accepted. The cuts described above are termed noise cuts. The ON–OFF statistical significance is calculated following Li & Ma (Li:1 (1983)), before and/or after application of all image cuts. $$S=\frac{ON\beta OFF}{\sqrt{ON+\beta ^2OFF}}$$ (1) and is used to assess the likelihood of a gamma-ray signal. In order to account for the mismatch of observation times between ON and OFF source data (and hence zenith angle-dependent event rates), and trigger rate differences due to subtle changes in weather conditions and/or telescope response during observation runs, a normalisation is applied to the ON–OFF statistical significance. This normalisation, $`\beta `$ is defined as the ratio of the events available for image parametrisation, i.e. after noise cuts. A final systematic check on the ON–OFF statistics after application of all cuts, performed on a run-by-run basis, is explained in Sect. 4. ## 3 Extended source analysis The analysis of extended sources has required the development of new techniques in TeV gamma-ray astronomy (Akerlof et al. Akerlof:1 (1991), Hess et al. Hess:1 (1997), Buckley et al. Buckley:1 (1998), Connaugton et al. Connaugton:1 (1998)). These methods use shape parameters such as width and length in combination with source position-dependent orientation and location cuts (e.g. alpha, asymmetry and dis). The orientation and location parameters are recalculated at every trial source position and a skymap is created. In the analysis of CANGAROO Vela Pulsar data, Yoshikoshi et al. (Yoshikoshi:1 (1997)) used a set of cuts based on alpha, length, width, concentration and dis. A continuous probability distribution was initially used to locate the position of the most significant point in the skymap. A gamma-ray flux at this point was then estimated from the ON–OFF excess obtained after using a combination of shape and location cuts with alpha$`10^{}`$. These cuts were optimised using Monte Carlo simulations and therefore are a priori decisions. We adopt the same a priori philosophy here in order to determine the significance of any result without the need to consider statistical penalties. The set of cuts described here were also used in the analysis of Centaurus A data described by Rowell et al. (Rowell:1 (1999)). The Monte Carlo simulation package, MOCCA92 (Hillas Hillas:3 (1995)), was used to generate Čerenkov images from extensive air showers (EAS). Gamma-ray primaries were sampled from a power law above 0.8 TeV with integral spectral index $``$1.6. Cosmic ray primaries, represented by a combination of proton, helium and nitrogen primaries, were sampled from a power law of spectral index $``$1.65 above 1.5 TeV. In designing an analysis suited to off-axis and extended sources, it is important to consider the off-axis sensitivity of the CANGAROO camera, particularly considering that it has a relatively small field of view and operates at a gamma-ray threshold of about 1.5 TeV. The behaviour of various image parameters as a function of gamma-ray source position has been investigated. A gamma-ray point source was placed at six positions across the camera and each position considered independently. The resulting distributions of width and miss showed little variation with source distance from the camera centre, in contrast to those of length and dis (Fig. 2). For increasing source distances from the camera centre, larger values of length and dis are possible, behaviour which is readily understood in terms of camera edge effect reduction on one side as the source approaches the camera edge. In particular, the length parameter for gamma-rays is very similar to that for cosmic-rays for sources beyond about 1.0 from the camera centre. Indeed, even for on-axis sources, the field of view of the 3.8 m camera imposes some edge effects, a consequence of operating at a relatively high gamma-ray energy threshold. Thus, for off-axis sources, a reduced gamma ray efficiency will result when using cuts derived for an on-axis source. Firstly, there is a clear case for increasing the gamma-ray efficiency while maintaining the quality factor on the grounds of improved event statistics. An increase in gamma-ray efficiency entails an increase in background acceptance, thus improving an estimate of the background and diluting any unaccounted systematic effects. It is also possible to examine the effect of the gamma-ray efficiency on the ON–OFF significance after cut application. We can express Eq. (1) in another way, incorporating the gamma-ray efficiency, after application of image cuts: $$S=\frac{N_\gamma }{\sqrt{\frac{N_\gamma }{ϵ_\gamma }+\frac{2N_b}{Q^2}}}$$ (2) where the gamma-ray signal, $`N_\gamma `$ ($`N_\gamma `$=ON–OFF) and the background, $`N_b`$ ($`N_b=`$OFF), represent those prior to image cuts. The efficiency for gamma-ray selection is given by $`ϵ_\gamma `$ and that for the background (CR), $`ϵ_b`$. The quality or Q-factor of the cuts is given by $`Q=ϵ_\gamma /\sqrt{ϵ_b}`$ . Here, we set the normalisation factor from Eq. (1), $`\beta =1`$. We can see that the significance obtained after image cuts is dependent on the quality factor and, somewhat slightly, on the gamma-ray cut efficiency. For sources with a high gamma-ray to CR flux ratio (e.g. $`>`$0.l) with $`Q`$4 (as for the CANGAROO telescope), both denominator terms are then similar, the sensitivity of $`S`$ to $`ϵ_\gamma `$ becomes more apparent. Such a situation may arise in the case of searches for bursts from AGN and/or gamma-ray bursts over short time scales. For signal to noise ratios expected of SNRs such as W28 however, we would expect only a minor improvement in $`S`$ from the above arguement. Thus, the main motivation for increasing $`ϵ_\gamma `$ here is simply to work with increased statistics. The variable cuts on length and for dis were incorporated into the cut ensemble. Along with the 3rd moment of the image, asymmetry, we use an approximation of the distance between the assumed source position and calculated source position for each image, $`D`$. This distance, expressed in units of standard deviation, is given by:- $$D=\sqrt{\left(\frac{miss}{\sigma _{miss}}\right)^2+\left(\frac{disdis_{ex}}{\sigma _{dis}}\right)^2}$$ (3) where $`\sigma _{miss},\sigma _{dis}`$ are the variances for miss and dis respectively. The variances, $`\sigma _{miss}`$ and $`\sigma _{dis}`$ represent the transverse and longitudinal errors in the most likely source position for an image. $`D`$ can be characterised as the source density function (SDF). When used in combination with shape and asymmetry cuts, a cut on $`D`$ provides a gamma-ray acceptance of $``$40% for an on-axis point source. $`D`$ is similar to the normalised cluster, or Mahalonobis distance for miss and dis (Hillas & West Hillas:2 (1991)), although here we are neglecting cross-term variances. $`D`$ can also be considered a discrete analogue of the probability distribution function used by Yoshikoshi et al. (Yoshikoshi:1 (1997)). The longitudinal error in source location is greater than the transverse error by about a factor of two. Following Akerlof et al. (Akerlof:1 (1991)), some reduction in the longitudinal error can be achieved by making use of the elongation (defined as length/width) of the image such that the expected longitudinal source distance is given by $`dis_{ex}=g(1(width/length))`$ where $`g`$ is an empirically derived constant. For the 3.8 m camera we derive a value of $`g=1.25`$ using simulations. Fig. 3 gives distributions of the parameter $`D`$ for simulated gamma-rays and cosmic rays, and a comparison to real OFF source data. Linear fits were found for the cut on length and value of $`\sigma _{dis}`$, respectively, as a function of source displacement from the camera centre. The total cut combination is listed in Table 1. These cuts (Table 1) provide a constant gamma-ray efficiency of $``$40% (and cosmic-ray efficiency) and quality factor $``$4 at the same value of $`D`$. Without the variable length and dis criteria, the gamma-ray acceptance quickly reduces to less than 30% for sources outside 0.5 from the camera centre for no significant improvement in Q-factor. Table 2 gives the simulated performance of the cut $`D`$ for various source positions. A comparison to the performance obtained by the set of cuts used in the Vela Pulsar analysis (Yoshikoshi et al. Yoshikoshi:1 (1997)) is included. These cuts were used to obtain a significant excess of gamma-ray-like events from a region displaced 0.14 from the pulsar position. The Vela Pulsar cuts are based on the same noise cuts listed in Table 1, except that fixed values of length and dis are used and a cut on alpha$`10^{}`$ is substituted for $`D2.0`$. The Vela Pulsar cuts do not provide a constant gamma-ray acceptance over the search region, and the background (cosmic ray) acceptance also decreases sharply. The resulting flux or upper limit from the search is found by dividing the ON–OFF excess, $`N`$ (or 3$`\sigma `$ upper limit thereof), by the position-dependent exposure (for an extended source, averaged over the integration region): $$F(>E)=\frac{N}{ϵ_\gamma \overline{\eta }AT}\mathrm{photons}\mathrm{cm}^2\mathrm{s}^1$$ (4) where $`\overline{\eta }`$ is the raw trigger efficiency for gamma-rays averaged over the integration region, $`ϵ_\gamma `$ is a constant gamma-ray selection efficiency for the SDF cut, $`A`$ is the area over which gamma-rays are simulated ($`A=1.96\times 10^9`$ cm<sup>2</sup>) and $`T`$ is the total observation time. We set $`ϵ_\gamma `$ to be the average of the simulated gamma-ray cut efficiencies out to $`\pm 1^{}`$, i.e. $`ϵ_\gamma `$ = 0.41. An unavoidable decline in $`\eta `$ for gamma-rays will occur as the source position is displaced further off-axis. Simulations show that $`\eta `$ decreases to less than half the on-axis value at the corners of the search. A linear fit was used to characterise $`\eta `$ at all points within $`\pm 1^{}`$ such that $`\eta _d=\eta _0(1.00.39d)`$ as a function of the source displacement from the camera centre, $`d`$. We find $`\eta _0`$(1995 data)$`=0.12`$ at the camera centre above the minimum simulated energy of 0.8 TeV for the 1995 dataset. The raw gamma-ray trigger efficiency for 1994 data was estimated by scaling the 1995 value by the ratio of observed event rates after noise cuts, giving $`\eta _0`$(1994 data)$`=0.21`$. In characterising $`ϵ_\gamma `$ and $`\overline{\eta }`$, we are assuming that the gamma-ray flux of an extended source is isotropically distributed. To calculate the flux applicable to the energy threshold of 1.5 TeV (Roberts et al. Roberts:1 (1998)) we used the gamma-ray spectral index of –1.6 adopted in the simulations. For a point-like search, the flux was taken as that from the point of interest, using a single value of $`\eta `$. The nature of our gamma-ray selection cuts naturally incorporates the gamma-ray point spread function (PSF). The cut on $`D`$ accepts events with derived source positions within an optimal radius. A search for an extended source therefore does not require any extra area to account for the PSF in addition to the source area itself. In an extended source search, at a suitably high number of assumed source positions in the region of interest, we sum the events passing all cuts, taking care not to count an event more than once. Skymaps of the statistical significance, $`S`$, of a gamma-ray signal were generated over a $`\pm 1^{}`$ area at 0.05 steps. At each grid point, representing an assumed source position, image parameters were calculated and the number of events passing the cuts of Table 1 cumulatively summed for all data. Since the resolution of the grid (0.05), is smaller than the effective acceptance area of the cuts (the cut on $`D`$ alone is more powerful than a cut on alpha$``$10), each skymap point value will not be fully independent of its neighbours. As a final check on data integrity, the distribution of $`S`$ obtained on a run-by-run basis (i.e. run-by-run skymap) was quantitatively assessed for systematic effects. The most important systematic effect to consider in this type of analysis is the consistency of the trigger threshold over the entire camera between ON and OFF source runs. Such an effect is difficult to compensate for after the data are taken. The Kolmolgorov-Smirnoff (KS) test is used to examine the likelihood that the distribution of $`S`$ of the skymap obtained on a run-by-run basis is derived from a normal distribution. Over the time scales of a single run ($``$5 hours), we do not expect significant contributions from a steady source of TeV gamma-rays of the strength expected of a SNR or plerion. A relatively strong KS probability of 4$`\sigma `$ was used to reject pairs of data that appeared severely affected by such systematics. For unmatched pairs, a well-behaved OFF or ON comparison run was used. A total of 10 hours data were rejected using the KS test, representing three ON/OFF pairs from the 1994 dataset. ## 4 Results Table 3 presents a summary of data accepted for analysis from the 1994 and 1995 observing seasons. The skymaps of significance $`S`$ (Fig. 4) for both years’ data, do not indicate any significant point-like excesses over a $`\pm 1^{}`$ search for $`\gamma `$-ray emission. Table 4 summarises the 3$`\sigma `$ upper limits for the positions of a number of interesting sites within the W28 region. It should be pointed out that our results will have a systematic error of the order $`20`$%, based on uncertainties in the trigger conditions, mirror reflectivity and spectral index adopted in the simulations (Yoshikoshi Yoshikoshi:2 (1996)). For example, in the next section we compare our results to a model producing a gamma-ray flux with spectral index of $`1.1`$, harder than the $`1.6`$ spectrum used in simulations. Such a difference however will contribute a systematic of $``$5%. The following positions were considered as potential sources; the position A83 given by Andrews et al. (Andrews:1 (1983)), the pulsar PSR J$`180123`$, the two strongest (by an order of magnitude) masers, labelled E & F by Claussen et al. (Claussen:1 (1997)), and the EGRET source 3EG J1800$``$2338. In addition, an extended region of radius 0.25 centred on a position to encompass the molecular clouds was considered. An average position was used for the masers E and F since they are separated by only 0.02. In assuming that the EGRET source was point-like, a search for the highest point significance within the 95% error circle was carried out. Since the statistical degrees of freedom for a non-a priori search for point-like emission over the skymap is $``$100, the highest ON–OFF excess within the EGRET error circle must be interpreted with a similar statistical penalty in mind. The point spread function (PSF) for a pure gamma-ray signal can be used to assess the location accuracy of the proposed source positions. Monte Carlo simulations show that the PSF for gamma-rays increases slightly from a standard deviation of 0.2 on-axis to 0.22 for sources at the skymap corners. The positions of our candidate sources (Table 4) lie within 1.0 on-axis. Taking a conservative estimate of the PSF as 0.22, a simplistic estimate of the source location error is obtained by adding in quadrature 0.22/$`\sqrt{100}`$=0.02 to the the absolute tracking precision (0.02), giving $`0.03^{}`$. We have used a value of 100 which is representative of the 3$`\sigma `$ upper limit excesses ($`N`$ in Eq. (4)) calculated here. We can therefore say that the features listed in Table 4 would be resolved by the CANGAROO 3.8m telescope if they were point sources of TeV gamma rays. From Table 4 we can also see that the positions of the highest significance within the EGRET 95% error circle for both year’s data are also not co-located, with their separation, $`0.6^{}`$, easily exceeding the estimated PSF. ## 5 Gamma Ray Production in SNR SNR shocks are able to accelerate particles to TeV energies. Gamma-rays are produced in secondary reactions between these high energy particles and ambient matter and radiation fields. The decay of $`\pi ^{}`$, produced in ion-ion collisions, is the prime hadronic source of gamma-rays in SNR. Electrons accelerated to multi-TeV energies in SNR give rise to bremsstrahlung and inverse Compton (IC) scattering $`\gamma `$-ray production processes in SNR (see e.g. Mastichiadis & de Jager Mastichiadis:1 (1996), Gaisser et al. Gaisser:1 (1998)). The cosmic microwave background (CMB) is usually considered the dominant soft photon source with contributions ($``$10%) from the infrared background as seed photons for the IC process. The detailed model of Baring et al. (Baring:1 (1999)) uses these processes collectively (including synchrotron and bremsstrahlung radiation) to account for observations from radio to TeV gamma-ray energies for the northern SNR, IC443. The models of Naito and Takahara (Naito:1 (1994)) and Drury et al. (Drury:1 (1994)) give predictions of the TeV gamma ray flux due to $`\pi ^{}`$ decay. As a first attempt at explaining the particle acceleration processes in W28, we make use of the Naito and Takahara model and compare its predictions of the $`\pi ^{}`$ channel to our results. By considering only the $`\pi ^{}`$ production channel, a lower limit on the predicted TeV gamma flux can be estimated. A proton parent spectrum with differential index $``$2.1 and exponential cutoff at 100 TeV has been assumed ($`dN/dEE^{2.1}\times \mathrm{exp}(E/100))`$, i.e. consistent with shock acceleration expected in a SNR. The predicted TeV flux will scale according to: $$F_\gamma \frac{E_{cr}n}{d^2}$$ (5) where $`E_cr`$ is the energy available for cosmic ray acceleration, $`n`$ is the particle number density of the ISM (cm<sup>-3</sup>) and $`d`$ is the distance to the SNR (kpc). $`E_{cr}`$ is some fraction, typically $``$10% of the total energy of the SNR (canonically 10<sup>51</sup>erg). In fact, the total SNR energy for W28 has been estimated by Rho et al. (Rho:2 (1996)) at 4$`\times 10^{50}`$erg from ROSAT X-ray data. The most interesting question concerns the possibly of enhanced TeV gamma-ray emission from regions of high ISM density (Aharonian et al. Aharonian:1 (1994)), where there is a greater chance for the interaction of cosmic rays from the SNR. The molecular clouds along the northeast and northern remnant boundary have been mapped in detail by Arikawa et al. (Arikawa:1 (1999)) at the CO J=1-0 and J=3-2 lines. The shocked (region undergone passage and compression by the SNR shock) component of the clouds is distributed along the SNR/clouds boundary, has a mean density of $`10^4`$ cm<sup>-3</sup>, and mass $`2\times 10^3`$M. The other 4$`\times 10^3`$M of the unshocked gas has a density of $`10^3`$cm<sup>-3</sup> and is displaced radially outward from the shocked gas by $``$1 arcmin. Clearly, any TeV gamma-ray flux from $`\pi ^{}`$ decay would be dominated by the shocked gas regions, given that the unshocked regions of the cloud would have a much lower density and lower energy available for CR production, and the mean matter density for regions excluding the molecular cloud is only $``$1.3 cm<sup>-3</sup> (Esposito et al. Esposito:1 (1996)). Arikawa et al. (Arikawa:1 (1999)) has derived the energy deposited into the shocked gas at $`3\times 10^{48}`$ergs, a value consistent with the Rho et al. (Rho:1 (1998)) estimate for the total SNR energy when considering the volume filling factor $`V`$, between the clouds and the SNR. $`V`$ is simplistically estimated at $`0.01`$, by taking the mass and density of the shocked gas (given above), assuming the cloud consists of H<sub>2</sub>, and a value of 10pc for the SNR radius. We also assume that $``$10% of the available SNR energy goes into cosmic ray production, a reasonable value for a Sedov-phase SNR, and also consistent with the measured energetics for the SNR and clouds. Thus we adopt values of 3$`\times 10^{47}`$ergs for $`E_{cr}`$, and $`10^4`$ cm<sup>-3</sup> for $`n`$ in Eq. (5). A working band on our flux prediction is obtained if we assume a range of values for $`d`$ from 1.8 to 3.3 kpc, as discussed in Sect. 1. In Fig. 5, we compare the model predictions based on the above scalings to our upper limit obtained from an extended source of radius 0.25 encompassing the clouds from 1994 data (the highest of our upper limits). The flux from 2EG J1801$``$2312 (essentially the same as 3EG J1800$``$2338) and its straight extrapolation to TeV energies is also included as any $`\pi ^{}`$ gamma-ray flux will be limited by the EGRET measurement. Our upper limit lies an order of magnitude below the straight extrapolation of the flux from EGRET (dashed line of Fig 5), and is able to place some constraint on the prediction of a $`\pi ^{}`$ gamma-ray flux from the shocked gas region. In order to accommodate our upper limit however a cutoff and/or slighter steeper parent spectrum than -2.1 appears necessary. In assuming an exponential cutoff at 100 TeV for the parent spectrum of accelerated hadrons, we assume that W28 follows the ’standard’ picture of particle acceleration in SNR, with particle energies limited by radiative losses, finite age of the shock and particle escape. When the cutoff for electrons is due only to radiative losses however, it can be expected that the hadron spectrum will continue (Reynolds & Keohane Reynolds:1 (1999)). Apart from the long-established ’knee’ at $`5\times 10^{15}`$ eV in the all-particle cosmic ray spectrum there is now direct experimental evidence pointing to a continuation of the proton/helium spectra up to at least $`800`$ TeV (Asakimori et al. Asakimori:1 (1998)), implying that strong cutoffs may not be required. Our result here is consistent however, with previous comparisons of upper limits (from other SNRs) to hadron-induced gamma-ray models (e.g. Buckley et al. Buckley:1 (1998), Allen et al. Allen:1 (1995) and Prosch et al. Prosch:1 (1996)) which seem to require some cutoff below the knee energy/and or spectra steeper than -2.1. Further constraints may arise if electronic components are considered. Particularly, electronic bremsstrahlung may dominate over the inverse Compton component due to the very high density of target matter in the clouds. The above discussion aside, the location of 3EG J1800$``$2338 by itself, makes it’s interpretation difficult in terms of simple CR-matter interaction, or as the result of a pulsar-powered process (proposed by Merck et al. Merck:1 (1996)). Many of the promising sites within W28 for gamma-ray production are now outside the 95% error circle, leaving just the filled X-ray centre, and southern/western portions of the radio shell. At best, we are in a position to rule out the interpretation of the EGRET source as resulting totally from $`\pi ^{}`$ decay gamma-rays with an unlimited parent spectrum, and to place limits on the parent spectral index/cutoff energy combination. ## 6 Conclusion A search for TeV gamma-ray emission from the southern SNR W28 was carried out by the CANGAROO over two observation seasons (1994 and 1995) using the atmospheric Čerenkov imaging technique. An analysis providing a consistent gamma-ray acceptance and quality factor for extended sources was used. A number of sites within a search region of $`\pm 1^{}`$ were considered as potential point-like and diffuse sources of TeV gamma-ray emission. No evidence was found for the emission of TeV gamma-rays at any of these sites which include those from the strongest two masers, an EGRET source, a radio pulsar (all as point sources) and a diffuse region containing the molecular clouds that appear to be interacting with the remnant. Our 3$`\sigma `$ upper limit from this diffuse region at 6.64$`\times 10^{12}`$ cm<sup>-2</sup>s<sup>-1</sup> $`>1.5`$ TeV, and the flux of the EGRET source 3EG J1800$``$2338 were compared with gamma-ray flux predictions from the model of Naito and Takahara (Naito:1 (1994)). Under this framework, our upper limit rules out a straight extrapolation of the EGRET flux to TeV energies. It also constrains somewhat the flux expected from the shocked region of gas in the molecular cloud, placing limits on the parent spectra for hadrons and/or cutoff energy. Our results suggest that the EGRET source probably does not result entirely of $`\pi ^{}`$ gamma-rays. This fact is supported by it’s location in relation to the molecular clouds. In a later paper, we will consider electronic bremsstrahlung and inverse Compton scattering and discuss the broader implications of our results in relation to the origin of galactic cosmic rays. Further data on W28 will no doubt be taken with the recently completed CANGAROO II telescope (Yoshikoshi et al. Yoshikoshi:3 (1999)), which, at the very least will provide tighter constraints on models of gamma-ray production for this interesting source. ###### Acknowledgements. This work is supported by a Grant-in-Aid in Scientific Research from the Japanese Ministry of Science, Sports and Culture, and also by the Australian Research Council. GR, JH, MR & TY acknowledge the receipt of JSPS postdoctoral fellowships. We also thank an anonymous referee for valuable comments.
warning/0004/math0004011.html	ar5iv	text	# 1 Introduction ## 1 Introduction Noncommutative space-time structures are finding increasing interest in gauge theories. Special examples have been discussed in the literature. Among them are 1) The canonical structure: $$[\widehat{x}^i,\widehat{x}^j]=i\mathrm{\Theta }^{ij},\mathrm{\Theta }^{ij}$$ 2) the Lie-algebra structure: $$[\widehat{x}^i,\widehat{x}^j]=ic_k^{ij}\widehat{x}^k,c_k^{ij}$$ 3) The quantum space structure: $$[\widehat{x}^i,\widehat{x}^j]=ic_{kl}^{ij}\widehat{x}^k\widehat{x}^l,c_{kl}^{ij}$$ We shall discuss a special example of the third case . For a more general review of the quantum space structure see . In all of the above cases we consider the associative algebra freely generated by the elements $`\widehat{x}^i`$ modulo the respective relations. This algebra of formal power series forms the algebra $`𝒜_x`$. $$𝒜_x=\frac{[[\widehat{x}^1\mathrm{}\widehat{x}^N]]}{R}$$ For a physicist this means that he is free to use the relations to reorder the elements of an arbitary power series. For the quantum space algebra we wish to exclude pathological cases such as the trivial case where there is no relation at all or where the product of any two elements is zero modulo the relations. To exclude such cases we shall demand the Poincaré-Birkhoff-Witt property for the algebra. By this we mean that the dimension of the space of homogeneous polynomials is the same as in the commutative case. For the first and second examples this will be the case, for the third example we require the Yang-Baxter equation. To formulate it we write the relations in the form: $$\widehat{x}^i\widehat{x}^j=\widehat{R}_{kl}^{ij}\widehat{x}^k\widehat{x}^l,\widehat{R}_{kl}^{ij}$$ and define $`N^3\times N^3`$ matrices $`\widehat{R}_{12j_1j_2j_3}^{i_1i_2i_3}`$ $`=`$ $`\widehat{R}_{j_1j_2}^{i_1i_2}\delta _{j_3}^{i_3}`$ $`\widehat{R}_{23j_1j_2j_3}^{i_1i_2i_3}`$ $`=`$ $`\delta _{j_1}^{i_1}\widehat{R}_{j_2j_3}^{i_2i_3}`$ The Yang-Baxter equation is: $$\widehat{R}_{12}\widehat{R}_{23}\widehat{R}_{12}=\widehat{R}_{23}\widehat{R}_{12}\widehat{R}_{23}$$ There are several known solutions of this equation. We are interested in such relations that allow a conjugation which makes the algebra a $``$-algebra. This is because we have to associate the observables like the coordinates with selfadjoint operators in a Hilbertspace. The $`\widehat{R}`$-matrices for the quantum groups $`SO_q(n)`$ allow such conjugations. The quantum space algebra is a comodule of a quantum group. We start from co-algebra relations $$\mathrm{\Delta }(x^i)=T_k^i\widehat{x}^k$$ and compute $`\mathrm{\Delta }(x^i)\mathrm{\Delta }(x^j)`$ $`=`$ $`T_k^iT_l^j\widehat{x}^k\widehat{x}^l`$ $`=`$ $`T_k^iT_l^j\widehat{R}_{mn}^{kl}\widehat{x}^m\widehat{x}^n`$ If we demand $`RTT`$ relations $$T_k^iT_l^j\widehat{R}_{mn}^{kl}=\widehat{R}_{kl}^{ij}T_m^kT_n^l$$ for the $`T`$-algebra, we find $$\mathrm{\Delta }(x^i)\mathrm{\Delta }(x^j)=\widehat{R}_{kl}^{ij}\mathrm{\Delta }(x^k)\mathrm{\Delta }(x^j)$$ There is always a solution to the $`RTT`$ relations given by $$T_j^i=\delta _j^i.$$ If this is the only solution then the bialgebra consists of the unit element only; not very interesting. If the $`\widehat{R}`$-matrix leads to the quantum group $`SO_q(n)`$, we have a more interesting case. Instead of introducing the quantum group $`SO_q(n)`$ we shall deal with the corresponding $`q`$-Lie algebra $`so_q(n)`$. The quantum space is then a module of this algebra. In this paper we discuss the 3-dimensional case in great detail. The algebra is introduced in Chapter 2. It has a peculiar property, there is a homomorphism of the algebra $`so_q(3)`$ into the algebra $`_q^3`$. This is discussed in Chapter 3. The full algebra can then be generated by a central element, the radius $`R`$, and elements of the tensor product of an $`su_q(2)`$ algebra and an $`su_q(1,1)`$ algebra. The generators of the $`su_q(2)`$ algebra are further restricted by relations that when the algebra is represented lead to a unique infinite-dimensional representation of $`su_q(2)`$. We call this algebra the $`t`$-algebra. The $`su_q(1,1)`$ algebra we call $`K`$-algebra. If we then demand that the $`so_q(3)`$ algebra corresponds to orbital angular momentum the $`K`$ algebra is restricted in the same sense as the $`t`$-algebra. This is discussed in Chapter 4. This clarifies the algebraic structure of the $`so_q(3)`$ module $`_q^3`$. To discuss physics we need representations of the algebra. The observables should be represented by (essentially) self-adjoint linear operators in a Hilbert space. This way we can use the well developped formalism of quantum mechanics and its interpretation scheme. In Chapter 5 we discuss the representations of the algebra. We find that they are characterized by one real parameter $`z_0`$. In all these representations we obtain a discrete spectrum for the coordinate $`X^3`$, which along with $`R`$ and $`T_{orb}^3`$, the third component of the orbital angular momentum, form a complete commuting set of observables. The scale of the spacing of the eigenvalues of $`X^3`$ is determined by the constant $`z_0`$, the eigenvalues are exponentially spaced. This we call a $`q`$-lattice. We are not surprised that noncommuting variables lead to a discretization (latticization) of space . In Chapter 6 we construct the transformation that leads to a basis where $`\stackrel{}{T}_{orb}^2`$ is diagonal. The corresponding transformation function turn out to be the $`q`$-deformed associated Legendre functions. They are defined in Appendix D in terms of the big $`q`$-Jacobi polynomials. They satisfy a difference equation, a recursion equation and have orthogonaltity properties - in complete analogy to the usual associated Legendre functions. From the self-adjointness property of $`X^3`$ we derive a completeness relation as well, this is done in Appendix E. Appendices A, B and C are devoted to the representation of the $`su_q(2)`$ and $`su_q(1,1)`$ algebras and their comultiplication. ## 2 The Algebra of the Euclidean Quantum Space $`_𝒒^\mathrm{𝟑}`$ This algebra has been discussed in Ref , we use the same notation here: $`_q^3`$: $`X^3X^+q^2X^+X^3`$ $`=`$ $`0`$ $`X^3X^{}q^2X^{}X^3`$ $`=`$ $`0`$ (2.1) $`X^{}X^+X^+X^{}`$ $`=`$ $`\lambda X^3X^3,\lambda =qq^1,q.`$ We shall assume $`q>1`$ in this paper. This non-commutative structure is our model for a non-commutative space. We can impose conjugation properties that are compatible with the relations (2.1) justifying the ‘$``$’ in $`_q^3`$: $$\overline{X^+}=qX^{},\overline{X^3}=X^3.$$ (2.2) The quantum space $`_q^3`$ has a co-module structure under the action of the quantum group $`SO_q(3)`$ and a module structure under the corresponding $`q`$-Lie algebra. $`su_q(2)`$: $`q^1T^+T^{}qT^{}T^+`$ $`=`$ $`T^3`$ $`q^2T^3T^+q^2T^+T^3`$ $`=`$ $`(q+q^1)T^+`$ (2.3) $`q^2T^{}T^3q^2T^3T^{}`$ $`=`$ $`(q+q^1)T^{}`$ The conjugation properties justifying the ‘$`u`$’ in $`su_q(2)`$ are: $$\overline{T^+}=\frac{1}{q^2}T^{},\overline{T^3}=T^3.$$ (2.4) The module structure that was found in Ref. is: $`T^3X^3`$ $`=`$ $`X^3T^3`$ $`T^3X^+`$ $`=`$ $`q^4X^+T^3+q^1(1+q^2)X^+`$ (2.5) $`T^3X^{}`$ $`=`$ $`q^4X^{}T^3q(1+q^2)X^{}`$ $`T^+X^3`$ $`=`$ $`X^3T^++q^2\sqrt{1+q^2}X^+`$ $`T^+X^+`$ $`=`$ $`q^2X^+T^+`$ (2.6) $`T^+X^{}`$ $`=`$ $`q^2X^{}T^++q^1\sqrt{1+q^2}X^3`$ $`T^{}X^3`$ $`=`$ $`X^3T^{}+q\sqrt{1+q^2}X^{}`$ $`T^{}X^+`$ $`=`$ $`q^2X^+T^{}+\sqrt{1+q^2}X^3`$ (2.7) $`T^{}X^{}`$ $`=`$ $`q^2X^{}T^{}`$ In the limit $`q=1`$ we obtain from relations (2.1)–(2.7) the commutative $`^3`$ with the Lie algebra $`so(3)`$ acting on it. As a consequence of the above relations it follows that there is a central hermitean element, the $`q`$-deformed radius: $`R^2`$ $`=`$ $`X^3X^3qX^+X^{}{\displaystyle \frac{1}{q}}X^{}X^+=q^2\overline{X^3}X^3+(1+q^2)\overline{X^+}X^+,`$ (2.8) $`\overline{R^2}`$ $`=`$ $`R^2.`$ ‘Central’ means that $`R^2`$ commutes with all the elements $`X`$ and $`T`$. There is a well-known Casimir operator for the $`su_q(2)`$ algebra: $$\stackrel{}{T}^2=\frac{q^2}{\lambda ^2}\tau ^{\frac{1}{2}}+\frac{1}{\lambda ^2}\tau ^{\frac{1}{2}}+\tau ^{\frac{1}{2}}T^+T^{}\frac{1+q^2}{\lambda ^2}.$$ (2.9) We have introduced the group-like element $$\tau =1\lambda T^3$$ (2.10) and the elements $`\tau ^{\frac{1}{2}}`$ and $`\tau ^{\frac{1}{2}}`$ as an extension of the algebra. We shall extend the algebra by the element $`R=(R^2)^{\frac{1}{2}}`$ and $`R^1=(R^2)^{\frac{1}{2}}`$ as well. The $`\tau X`$ and $`\tau T`$ commutation relations can be obtained from the $`T^3X`$ and $`T^3T`$ relations and vice versa. They are $`\tau X^3`$ $`=`$ $`X^3\tau `$ $`\tau X^+`$ $`=`$ $`q^4X^+\tau `$ (2.11) $`\tau X^{}`$ $`=`$ $`q^4X^{}\tau `$ and $`\tau T^3`$ $`=`$ $`T^3\tau `$ $`\tau T^+`$ $`=`$ $`q^4T^+\tau `$ (2.12) $`\tau T^{}`$ $`=`$ $`q^4T^{}\tau .`$ The definition of the orbital angular momentum as it was given in Ref. can be best formulated in terms of the elements $`L^+`$ $`=`$ $`{\displaystyle \frac{1}{q^2\sqrt{1+q^2}}}\tau ^{\frac{1}{2}}T^+`$ $`L^{}`$ $`=`$ $`{\displaystyle \frac{1}{q^3\sqrt{1+q^2}}}\tau ^{\frac{1}{2}}T^{}`$ (2.13) $`L^3`$ $`=`$ $`{\displaystyle \frac{1}{q^2(1q^2)}}\left(\tau ^{\frac{1}{2}}1{\displaystyle \frac{\lambda ^2}{1+q^2}}\stackrel{}{T}^2\right)`$ As the $`q`$-generalization of the fact that orbital angular momentum is orthogonal to the coordinate vector we impose the constraint $$LX=L^3X^3qL^+X^{}\frac{1}{q}L^{}X^+=0.$$ (2.14) We shall see that this defines orbital angular momentum uniquely. ## 3 The $`𝒕`$ Algebra: The algebra introduced in the previous chapter allows a homomorphism of the $`T`$ algebra into the $`X`$ algebra. This was first seen in Ref.. We find this homomorphism by interpreting Eqns. (2.5), (2.6) and (2.7) as inhomogeneous equations which can be solved for $`T`$ in terms of $`X`$. We first construct a particular solution $`t`$ and exhibit the homomorphism $`T^+`$ $``$ $`t^+={\displaystyle \frac{1}{\lambda q^3}}\sqrt{1+q^2}X^+(X^3)^1`$ $`T^{}`$ $``$ $`t^{}={\displaystyle \frac{q^2}{\lambda }}\sqrt{1+q^2}X^{}(X^3)^1`$ (3.1) $`T^3`$ $``$ $`t^3={\displaystyle \frac{1}{\lambda }}\left(1+R^2(X^3)^2\right).`$ Here we extend the algebra by the inverse of $`X^3`$. To establish the homomorphism we have to use (2.1) to show that the $`t`$ elements satisfy (2.3). Furthermore the relations (2.5) to (2.7) are fullfilled by the $`t`$ elements. It is due to (2.2) that they satisfy (2.4) as well. There are additional relations for the $`t`$ elements that follow from (2.1). They are $$\tau _t=1\lambda t^3=R^2(X^3)^2$$ (3.2) and $`t^+t^{}`$ $`=`$ $`{\displaystyle \frac{1}{\lambda ^2}}\left(1+q^2\tau _t\right)`$ (3.3) $`t^{}t^+`$ $`=`$ $`{\displaystyle \frac{1}{\lambda ^2}}\left(1+{\displaystyle \frac{1}{q^2}}\tau _t\right).`$ It follows that the Casimir operator for the $`t`$ algebra takes a definite value and that in the notation of Appendix A, where $`\overline{m}_t`$ and $`d_t`$ are defined, $$\stackrel{}{T}^2=\frac{1+q^2}{\lambda ^2},\overline{m}_t=0,d_t=\frac{q^2}{\lambda }.$$ (3.4) This value of the Casimir operator and the sign of $`\tau _t`$, which is negative, show that the $`t`$ algebra cannot be represented by the well-known finite dimensional representations of the $`T`$ algebra . In Appendix A we shall show that there are infinite-dimensional representations of the $`T`$ algebra among which there is one satisfying (3.2), (3.3) and (3.4). The representation is uniquely determined by these conditions, we present it here: $`t^3\|m_t`$ $`=`$ $`{\displaystyle \frac{1}{\lambda }}\left(1+q^2q^{4m_t}\right)\|m_t`$ $`t^+\|m_t`$ $`=`$ $`{\displaystyle \frac{1}{\lambda q}}\sqrt{q^{4m_t}1}\|m_t+1`$ (3.5) $`t^{}\|m_t`$ $`=`$ $`{\displaystyle \frac{q}{\lambda }}\sqrt{q^{4(m_t1)}1}\|m_t1`$ $`m_t0.`$ From (3.5) it follows that $$t^+\|0=0.$$ (3.6) There is no state with positive $`m_t`$. Eqns. (3.1) allow us to express the elements $`XR^1`$ in terms of the $`t`$ elements: $`X^3R^1`$ $`=`$ $`\pm (\tau _t)^{\frac{1}{2}}`$ $`X^+R^1`$ $`=`$ $`{\displaystyle \frac{\lambda q^3}{\sqrt{1+q^2}}}t^+(\tau _t)^{\frac{1}{2}}`$ (3.7) $`X^{}R^1`$ $`=`$ $`\pm {\displaystyle \frac{\lambda }{q^2\sqrt{1+q^2}}}t^{}(\tau _t)^{\frac{1}{2}}`$ The two different signs are the signs of $`X^3R^1=\sqrt{(X^3)^2R^2}`$. These elements can be viewed as homogeneous coordinates in the $`_q^3`$ space. The representations of these elements are now obtained from (3.5): $`X^3R^1\|m_t`$ $`=`$ $`\pm q^{2m_t1}\|m_t`$ $`X^+R^1\|m_t`$ $`=`$ $`{\displaystyle \frac{q}{\sqrt{1+q^2}}}\sqrt{1q^{4m_t}}\|m_t+1`$ (3.8) $`X^{}R^1\|m_t`$ $`=`$ $`\pm {\displaystyle \frac{1}{\sqrt{1+q^2}}}\sqrt{1q^{4(m_t1)}}\|m_t1`$ The different signs in (3.8) lead to inequivalent irreducible representations of the $`X`$ algebra. ## 4 The $`𝑲`$ Algebra We continue to consider the Eqns. (2.5), (2.6) and (2.7) as inhomogeneous equations that should be solved for the $`T`$’s. We have found one particular solution (3.1) and now move to the homogeneous part. This we do by the Ansatz: $`T^\pm `$ $`=`$ $`\mathrm{\Delta }^\pm +t^\pm `$ (4.1) $`T^3`$ $`=`$ $`\mathrm{\Delta }^3+t^3`$ Eqns (2.5), (2.6) and (2.7) become homogeneous equations for the $`\mathrm{\Delta }`$’s. $`X^\pm \mathrm{\Delta }^+`$ $`=`$ $`q^{\pm 2}\mathrm{\Delta }^+X^\pm `$ (4.2) $`X^3\mathrm{\Delta }^+`$ $`=`$ $`\mathrm{\Delta }^+X^3`$ $`X^\pm \mathrm{\Delta }^{}`$ $`=`$ $`q^{\pm 2}\mathrm{\Delta }^{}X^\pm `$ (4.3) $`X^3\mathrm{\Delta }^{}`$ $`=`$ $`\mathrm{\Delta }^{}X^3`$ $`X^\pm \mathrm{\Delta }^3`$ $`=`$ $`q^{\pm 4}\mathrm{\Delta }^3X^\pm `$ (4.4) $`X^3\mathrm{\Delta }^3`$ $`=`$ $`\mathrm{\Delta }^3X^3`$ These equations suggest the further Ansatz: $`K^\pm `$ $`=`$ $`\pm (\tau _t)^{\frac{1}{2}}\mathrm{\Delta }^\pm ,`$ (4.5) $`K^3`$ $`=`$ $`(\tau _t)^1\mathrm{\Delta }^3.`$ The element $`\tau _t`$ satisfies the relation (2) and as a consequence all the $`K`$s commute with all the $`X`$s and therefore with all the $`t`$’s as well $`K^AX^B`$ $`=`$ $`X^BK^A`$ (4.6) $`K^At^B`$ $`=`$ $`t^BK^A`$ Now we turn to (2.3) and compute the $`KK`$ relations: $`q^1K^+K^{}qK^{}K^+`$ $`=`$ $`K^3`$ $`q^2K^3K^+q^2K^+K^3`$ $`=`$ $`(q+q^1)K^+`$ (4.7) $`q^2K^3K^{}+q^2K^{}K^3`$ $`=`$ $`(q+q^1)K^{}`$ This is exactly the same algebra as (2.3). Any realization of the $`K`$-algebra will lead to a realization of the $`T`$-algebra: $`T^\pm `$ $`=`$ $`t^\pm \pm (\tau _t)^{\frac{1}{2}}K^\pm `$ (4.8) $`T^3`$ $`=`$ $`t^3+\tau _tK^3`$ This is a relation which is familiar from the comultiplication of two representations of the algebra (2.3): $`\mathrm{\Delta }_\beta (T^3)`$ $`=`$ $`T^31+\tau T^3`$ (4.9) $`\mathrm{\Delta }_\beta (T^\pm )`$ $`=`$ $`T^\pm 1\pm \sqrt{\tau }T^\pm `$ This comultiplication will be discussed in Appendix C. It is adjusted to representations where the first factor has negative eigenvalues of $`\tau `$. We emphasize that the representations of the $`t`$ algebra in (4.8) are restricted by the relations (3.3) whereas for the $`K`$ algebra any representation would do as long as we are not considering any conjugation properties. If we now demand the conjugation property (2.4) for the $`T`$ algebra we find for the $`K`$ algebra: $$\overline{K^3}=K^3,\overline{K^+}=\frac{1}{q^2}K^{}.$$ (4.10) Note the sign. The $`K`$ algebra belongs to the $`SU_q(1,1)`$ quantum group. If we now use the condition (2.14) for orbital angular momentum we will specify the $`K`$ algebra representation uniquely as well. It needs some computation to express the $`L`$ algebra (2.13) in terms of the $`t`$ and $`K`$ algebras. $`L^+`$ $`=`$ $`{\displaystyle \frac{1}{q^2\sqrt{1+q^2}}}\left\{(\tau _t)^{\frac{1}{2}}t^+(\tau _k)^{\frac{1}{2}}+1(\tau _k)^{\frac{1}{2}}K^+\right\}`$ $`L^{}`$ $`=`$ $`{\displaystyle \frac{1}{q^3\sqrt{1+q^2}}}\left\{(\tau _t)^{\frac{1}{2}}t^{}(\tau _k)^{\frac{1}{2}}1(\tau _k)^{\frac{1}{2}}K^{}\right\}`$ (4.11) $`L^3`$ $`=`$ $`{\displaystyle \frac{q^21}{q^4(q^2+1)}}\{{\displaystyle \frac{q^2}{\lambda ^2}}(\tau _t)^{\frac{1}{2}}(\tau _k)^{\frac{1}{2}}+(\tau _t)^{\frac{1}{2}}(\tau _k)^{\frac{1}{2}}(K^+K^{}+{\displaystyle \frac{q^2}{\lambda ^2}})`$ $`{\displaystyle \frac{1+q^2}{\lambda ^2}}(\tau _t)^{\frac{1}{2}}(\tau _k)^{\frac{1}{2}}+t^{}(\tau _k)^{\frac{1}{2}}K^+q^2t^+(\tau _k)^{\frac{1}{2}}K^{}\}`$ This already shows that we should restrict the representations such that $`\tau _k`$ has negative eigenvalues. There is an additional reason for it. We shall see in Appendix C that the coproduct (4.8) only leads to representations with positive eigenvalues of $`\tau `$ if $`\tau _k`$ has negative eigenvalues. Only in this case the representations of $`T`$ can be decomposed into finite-dimensional ones. We are here adding this as an additional assumption - not knowing if it is really necessary. With this assumption it will follow from the comultiplication rule of Appendix C that we have to choose $$d_k=\frac{1}{\lambda q^2}.$$ (4.12) Now we are ready to evaluate (2.14). This relation will be true if and only if: $$\underset{¯}{m}_k=1.$$ (4.13) This is in the notation of Appendix B. For the Casimir operator we find $$\stackrel{}{T}_k^2=\frac{1+q^2}{\lambda ^2}.$$ (4.14) This uniquely determines the $`K`$ algebra representation. The generators of the orbital angular momentum will be denoted by $`T_{orb}`$. We find the result: $`T_{orb}^3`$ $`=`$ $`t^31+\tau _tK^3`$ (4.15) $`T_{orb}^\pm `$ $`=`$ $`t^\pm 1\pm \sqrt{\tau _t}K^\pm `$ where the $`t`$ and $`K`$ representations are determined by (3.4) and (4.12), (4.13). We can add spin to orbital angular momentum: $`T^3`$ $`=`$ $`T_{orb}^31+\tau _{orb}S^3`$ (4.16) $`T^\pm `$ $`=`$ $`T_{orb}^\pm 1+\sqrt{\tau _{orb}}S^\pm `$ The spin operators $`S`$ can be in any finite-dimensional representations of the $`T`$ algebra. ## 5 Representations of the $`𝑻_{𝒐𝒓𝒃}`$ Algebra: The representation of the $`K`$ algebra that enters orbital angular momentum is characterized by (4.12), (4.13) and (4.14): $$d_k=\frac{1}{\lambda q^2},\underset{¯}{m}_k=1,\stackrel{}{T}_k^2=\frac{1+q^2}{\lambda ^2}.$$ (5.1) It is an infinite-dimensional representation with $`m_k`$ ranging from $`0`$ to $`\mathrm{}`$. $`K^3\|m_k`$ $`=`$ $`{\displaystyle \frac{1}{\lambda }}\left(1+{\displaystyle \frac{1}{q^2}}q^{4m_k}\right)\|m_k`$ $`K^+\|m_k`$ $`=`$ $`{\displaystyle \frac{1}{q\lambda }}\sqrt{1q^{4(m_k+1)}}\|m_k+1`$ (5.2) $`K^{}\|m_k`$ $`=`$ $`{\displaystyle \frac{q}{\lambda }}\sqrt{1q^{4m_k}}\|m_k1`$ $`K^{}\|0`$ $`=`$ $`0,m_k0`$ The representation $`T_{orb}`$ of orbital angular momentum is the tensor product of this representation and the $`t`$ representation given in (3.5). The eigenstates of $`T_{orb}^3`$ are characterized by the two numbers $`m_t`$ and $`m_k`$. $`T_{orb}^3\|m_t,m_k`$ $`=`$ $`{\displaystyle \frac{1}{\lambda }}\left(1q^{4(m_t+m_k)}\right)\|m_t,m_k`$ (5.3) $`T_{orb}^+\|m_t,m_k`$ $`=`$ $`{\displaystyle \frac{1}{\lambda q}}\sqrt{q^{4m_t}1}\|m_t+1,m_k`$ $`+{\displaystyle \frac{1}{\lambda }}q^{2m_t}\sqrt{1q^{4(m_k+1)}}\|m_t,m_k+1`$ $`T_{orb}^{}\|m_t,m_k`$ $`=`$ $`{\displaystyle \frac{q}{\lambda }}\sqrt{q^{4(m_t1)}1}\|m_t1,m_k`$ $`+{\displaystyle \frac{q^2}{\lambda }}q^{2m_t}\sqrt{1q^{4m_k}}\|m_t,m_k1`$ In this representation $`\stackrel{}{K}^2`$ and $`\stackrel{}{t}^2`$ are diagonal. The transformation to the basis where $`\stackrel{}{T}_{orb}^2`$ is diagonal will be constructed in the next chapter. The value of $`T_{orb}^3`$ in (5.3) shows that we have found finite-dimensional representations of $`T_{orb}`$. We obtain from (5.3) $$\tau _{orb}\|m_t,m_k=q^{4(m_t+m_k)}\|m_t,m_k.$$ (5.4) The representation of the $`X`$ algebra can be obtained from (3.8). The element $`R`$ is central, it will be diagonal in the $`m_t,m_k`$ basis. We denote the eigenvalue of $`R^2`$ by $$R^2\|m_t,m_k,M=q^{4M+2}z_0^2\|m_t,m_k,M,$$ (5.5) where $`z_0`$ is an arbitrary parameter characterizing the radius. Then we obtain from (3.8) the representation of $`X^3`$: $$X^3\|m_t,m_k,M=q^{2(m_t+M)}z_0\|m_t,m_k,M.$$ (5.6) We have absorbed the sign in (3.8) yielding inequivalent representations in the sign of $`z_0`$ which is not determined by (5.5). This and (5.4) suggest that we should introduce a notation characterizing the eigenvalue of $`X^3`$ by a quantum number as well as the eigenvalue of $`T_{orb}^3`$. $$\nu =m_t+M,m=m_t+m_k$$ (5.7) In this notation we obtain the representation which was also found in : $`X^3\|M,\nu ,m`$ $`=`$ $`q^{2\nu }z_0\|M,\nu ,m`$ $`R^2\|M,\nu ,m`$ $`=`$ $`q^{4M+2}z_0^2\|M,\nu ,m`$ $`T_{orb}^3\|M,\nu ,m`$ $`=`$ $`{\displaystyle \frac{1}{\lambda }}\left(1q^{4m}\right)\|M,\nu ,m`$ $`X^+\|M,\nu ,m`$ $`=`$ $`{\displaystyle \frac{q^2z_0}{\sqrt{1+q^2}}}\sqrt{q^{4M}q^{4\nu }}\|M,\nu +1,m+1`$ (5.8) $`X^{}\|M,\nu ,m`$ $`=`$ $`{\displaystyle \frac{qz_0}{\sqrt{1+q^2}}}\sqrt{q^{4M}q^{4(\nu 1)}}\|M,\nu 1,m1`$ $`T_{orb}^+\|M,\nu ,m`$ $`=`$ $`{\displaystyle \frac{1}{q^21}}\sqrt{q^{4(M\nu )}1}\|M,\nu +1,m+1`$ $`+{\displaystyle \frac{1}{\lambda }}\sqrt{q^{4(M\nu )}q^{4(m+1)}}\|M,\nu ,m+1`$ $`T_{orb}^{}\|M,\nu ,m`$ $`=`$ $`{\displaystyle \frac{q^2}{q^21}}\sqrt{q^{4(M\nu +1)}1}\|M,\nu 1,m1`$ $`+{\displaystyle \frac{q^2}{\lambda }}\sqrt{q^{4(M\nu )}q^{4m}}\|M,\nu ,m1`$ $`\nu M,m\nu M`$ ## 6 Reduction of the representation of $`𝑻_{𝒐𝒓𝒃}`$ The above representation (5.3) of $`T_{orb}`$ is a tensor product of two representations, with $`\stackrel{}{t}^2`$ and $`\stackrel{}{K}^2`$ diagonal. We proceed with its decomposition into a sum of irreducible representations characterized by the eigenvalues of $`\stackrel{}{T}_{orb}^2`$. From the Appendix A we know that for $`d=\lambda d_td_k=\lambda ^1`$ the eigenvalues of $`\stackrel{}{T}^2`$ are $`q[l][l+1]`$. Therefore we start with an Ansatz of the form $`\|l,m`$ $`=`$ $`{\displaystyle \underset{m_k,m_t}{}}c_{l,m}^{m_k,m_t}\|m_t,m_k,m_k0,m_t0`$ (6.1) $`\stackrel{}{T}_{orb}^2\|l,m`$ $`=`$ $`q[l][l+1]\|l,m.`$ According to (5.7) $`m=m_t+m_k`$, so that we have $$c_{l,m}^{m_k,m_t}=c_{l,m}^{m_t}\delta _{m,m_t+m_k}.$$ (6.2) From the definition (2.9) of $`\stackrel{}{T}^2`$ and the Equations (5.3) we obtain a recursion relation for the coefficients $`c_{l,m}^{m_t}`$. $`\left(q^{2l+2}+q^{2l}(q^2+1)q^{2(m+1)4m_t}\right)c_{l,m}^{m_t}=`$ $`q^{2m+1}(\sqrt{(q^{4m_t}1)(q^{4m_t}q^{4m})}c_{l,m}^{m_t+1}`$ $`+\sqrt{(q^{44m_t}1)(q^{44m_t}q^{4m})}c_{l,m}^{m_t1})`$ A comparison with the $`q`$-difference Equation (D) for the functions $`\stackrel{~}{P}_m^l`$ defined in (D.17) and (D.25) shows that (6) is solved by $$c_{l,m}^{m_t}=\{\begin{array}{cc}q^{m_tm1}\sqrt{1q^2}\stackrel{~}{P}_l^m(\pm q^{2(m_t1)2m})\hfill & \text{ for }m0\hfill \\ q^{m_t1}\sqrt{1q^2}\stackrel{~}{P}_l^{\|m\|}(\pm q^{2(m_t1)})\hfill & \text{ for }m<0\hfill \end{array}$$ (6.4) Note that $`P_l^m`$ is defined for $`m0`$ only. The orthogonality condition (D.28) for the functions $`\stackrel{~}{P}_l^m`$ suggests to start with the direct sum of two representations of the form (5.3), such that both signs of the argument of $`\stackrel{~}{P}_l^m`$ appear. $`\|l,m`$ $`=`$ $`{\displaystyle \underset{\sigma =\pm 1}{}}{\displaystyle \underset{m_t}{}}c_{l,m}^{m_t,\sigma }\|m_t,m_k,\sigma `$ (6.5) $`c_{l,m}^{m_t,\sigma }`$ $`=`$ $`\{\begin{array}{cc}\sqrt{1q^2}q^{m_t1m}\stackrel{~}{P}_l^m(\sigma q^{2(m_tm1})\hfill & \text{ for }m0\hfill \\ & \\ \sqrt{1q^2}q^{m_t1}\stackrel{~}{P}_l^{\|m\|}(\sigma q^{2(m_t1)})\hfill & \text{ for }m<0\hfill \end{array}`$ (6.9) We know that $`m_k0`$, $`m_t0`$ and $`m=m_t+m_k`$, thus $`m_t`$ is restricted by $`m_t0`$ and $`mm_t`$. The last condition comes into effect for negative values of $`m`$. Note that if $`m_t`$ takes its largest allowed value, the coefficient of $`c_{l,m}^{m_t+1}`$ in (6) vanishes. We are free to choose this $`c`$ to be zero. For $`m0`$ it then follows from (6) that $`c_{l,m}^{m_t}=0`$ for $`m_t>0`$ and for $`m<0`$ the same is true for $`m_t>m`$. The values of $`l`$ are restricted by the condition $`\|m\|l`$, as seen from (D.18). This is obviously consistent with the recursion formula (6). We have chosen the normalization in (6.9) in such a way that according to (D.28) the eigenfunctions of $`\stackrel{}{T}_{orb}^2`$ are orthonormal: $$\begin{array}{c}q^1\lambda \underset{\sigma =\pm 1}{}\underset{m_t=\mathrm{}}{\overset{\mathrm{min}\{0,m\}}{}}q^{2(m_t1)m\|m\|}\stackrel{~}{P}_l^{\|m\|}(\sigma q^{2(m_t1)m\|m\|})\stackrel{~}{P}_l^{}^{\|m\|}(\sigma q^{2(m_t1)m\|m\|})\hfill \\ \\ =\delta _{l,l^{}}\hfill \end{array}$$ (6.10) To see this for $`m<0`$ it is enough to shift the summation variable $`m_tm_t+m`$. We now assume that the two representations with $`\sigma =+1`$ and $`\sigma =1`$ also lead to a different sign of $`z_0`$ in (5.6). $$X^3\|m_t,m_k,M,\sigma =q^{2(m_t+M)}\sigma \|z_0\|\|m_t,m_k,M,\sigma $$ (6.11) Then it follows from (E.7) that the functions $`\stackrel{~}{P}_l^m`$ satisfy the following completeness relation $$\begin{array}{c}q^1\lambda \underset{l=0}{\overset{\mathrm{}}{}}q^{m_t+m_t^{}2m\|m\|}\stackrel{~}{P}_l^{\|m\|}(\sigma q^{2(m_t1)m\|m\|})\stackrel{~}{P}_l^{\|m\|}(\sigma ^{}q^{2(m_t^{}1)m\|m\|})\hfill \\ \\ =\delta _{\sigma ,\sigma ^{}}\delta _{m_t,m_t^{}}\hfill \end{array}$$ (6.12) This construction shows that for fixed $`m`$, $`\stackrel{}{T}_{orb}^2`$ is a selfadjoint operator in the basis $`\|m_t,m_k,\sigma `$, $`\sigma =\pm 1`$, and that the transformation from the basis $`\|m_t,m_k,\sigma `$ to the basis $`\|l,m`$ is an isometry. ## Appendices ## A Representations of the $`𝑻`$ Algebra When constructing representations of the $`T`$ algebra, we are aiming at representations where $`T^3`$ is selfadjoint (or essentially selfadjoint). This allows us to assume $`T^3`$ to be diagonal: $$T^3\|m=f(m)\|m.$$ (A.1) The eigenvalue of $`T^3`$ is $`f(m)`$, $`m`$ is a labelling of the eigenstates. The second equation of (2.3) shows that $`T^+\|m`$ is again an eigenstate of $`T^3`$, we choose the labelling such that this state is labelled by $`m+1`$: $$T^+\|m=c_m\|m+1.$$ (A.2) The relation (2.3) leads to a recursion formula for $`f(m)`$: $$f(m+1)=\frac{1}{q^4}f(m)+\frac{1}{q^2}(q+q^1).$$ (A.3) This recursion formula has the solution $$f(m)=\frac{1}{\lambda }dq^{4m}.$$ (A.4) From $`\overline{T^3}=T^3`$ follows that $`dq^{4m}`$ has to be real. We take $`d`$ and $`m`$ to be real. For the operator $`\tau `$ of (2.10) follows $$\tau \|m=\lambda dq^{4m}\|m.$$ (A.5) From the conjugation properties of $`T^+`$ it follows that $$T^{}\|m=q^2c_{m1}^{}\|m1.$$ (A.6) The third equation of (2.3) is the conjugate of the second one. The first equation of (2.3) amounts to a recursion formula for $`c_m^{}c_m`$ $$qc_{m1}^{}c_{m1}q^3c_m^{}c_m=f(m).$$ (A.7) This recursion formula can be solved: $$c_m^{}c_m=\frac{1}{\lambda }\left\{\frac{1}{q^2\lambda }+\alpha \lambda q^{2m}\frac{d}{q^4}q^{4m}\right\}.$$ (A.8) The real parameter $`\alpha `$ is not determined by (A.7). We see that $`c_m^{}c_m`$ becomes negative for $`m\mathrm{}`$. This is not allowed. There has to be a largest value of $`m`$, say $`\overline{m}`$, such that $$c_{\overline{m}}^{}c_{\overline{m}}=0.$$ (A.9) Then it follows from (A.2) that $`T^+`$ does not lead to a state with a larger value than $`\overline{m}`$. To analyze this situation we introduce the function: $`x`$ $`=`$ $`q^{2m}`$ (A.10) $`h(x)`$ $`=`$ $`\left\{{\displaystyle \frac{1}{q(q^21)}}+\alpha \lambda x{\displaystyle \frac{d}{q^4}}x^2\right\}`$ The function $`h(x)`$ is negative for $`x=0`$, the sign of $`h(x)`$ for $`x\mathrm{}`$ depends on the sign of $`d`$. In any case $`h(x)`$ has to have a zero for positive $`x`$ to represent $`c^{}c`$. We have to demand $$x_1=q^{2\overline{m}},h(x_1)=0.$$ (A.11) The parameter $`\alpha `$ can now be expressed in terms of $`\overline{m}`$. $$\alpha =\frac{1}{\lambda }\left\{\frac{1}{\lambda }q^{2(\overline{m}1)}+dq^{2(\overline{m}+2)}\right\}.$$ (A.12) If $`\alpha `$ takes this value $`h(x)`$ has the two zeros: $$x_1=q^{2\overline{m}},x_2=\frac{1}{\lambda d}q^{2(\overline{m}+1)}.$$ (A.13) We obtain for $`c^{}c`$: $$c_m^{}c_m=\frac{d}{\lambda q^4}\left(q^{2m}q^{2\overline{m}}\right)\left(q^{2m}\frac{1}{\lambda d}q^{2(\overline{m}+1)}\right).$$ (A.14) The representation is characterized by two parameters, $`\overline{m}`$ and $`d`$. We use this parameter in the explicit form of the matrix elements: $`T^3\|m`$ $`=`$ $`({\displaystyle \frac{1}{\lambda }}dq^{4m})\|m`$ $`T^+\|m`$ $`=`$ $`{\displaystyle \frac{1}{q^2}}\sqrt{{\displaystyle \frac{d}{\lambda }}(q^{2m}q^{2\overline{m}})\left({\displaystyle \frac{1}{\lambda d}}q^{2(\overline{m}+1)}q^{2m}\right)}\|m+1`$ (A.15) $`T^{}\|m`$ $`=`$ $`\sqrt{{\displaystyle \frac{d}{\lambda }}(q^{2(m1)}q^{2\overline{m}})\left({\displaystyle \frac{1}{\lambda d}}q^{2(\overline{m}+1)}q^{2(m1)}\right)}\|m1`$ $`\tau \|m`$ $`=`$ $`d\lambda q^{4m}\|m`$ $`\stackrel{}{T}^2`$ $`=`$ $`{\displaystyle \frac{1}{\lambda ^2\sqrt{\lambda d}}}\left(q^{2(\overline{m}+1)}+\lambda dq^{2\overline{m}}\right){\displaystyle \frac{1+q^2}{\lambda ^2}}`$ Let us now have a closer look at the condition $`c^{}c0`$. For this purpose we discuss the three cases $`d>0`$, $`d=0`$ and $`d<0`$ separately. $`d>0`$ There has to be a smallest value of m, say $`\underset{¯}{m}`$, such that $`\|\underset{¯}{m}0`$ and $`T^{}\|\underset{¯}{m}=0`$, therefore $$c_{\underset{¯}{m}1}^{}c_{\underset{¯}{m}1}=0.$$ (A.16) From (A.14) follows for $`\overline{m}0`$ $$d=\frac{1}{\lambda },\underset{¯}{m}=\overline{m}.$$ (A.17) The number of states between $`\overline{m}`$ and $`\underset{¯}{m}`$ has to be integer: $$2\overline{m}+1n.$$ (A.18) This shows that $`\overline{m}`$ has to be integer or half integer and we found the $`2l+1`$-dimensional representation ($`\overline{m}=l`$) of $`so_q(3)`$. $`d=0`$ In this case $`h(x)`$ is a linear function: $$h(x)=\frac{1}{q(q^21)}+\alpha \lambda x.$$ (A.19) Now $`\alpha `$ has to be positive for $`h`$ to have a zero for positive $`x`$. From (A.12) follows $$\alpha =\frac{1}{q^2\lambda ^2}q^{2\overline{m}}.$$ (A.20) The representation can be obtained from (A). The parameter $`\overline{m}`$ that characterizes the representation can take any real value. The representation is infinite-dimensional, however, $`\tau `$ is not invertible. $`d<0`$ This is the situation that arises for the $`t`$ algebra, as can be seen from (3.2). In this case $`x_2`$ is negative. We only have a largest value of $`m`$. The representation is infinite-dimensional and $`\overline{m}`$ is not restricted. The matrix elements are obtained from (A). We write them such as to exhibit the positive square roots: $`T^3\|m`$ $`=`$ $`({\displaystyle \frac{1}{\lambda }}dq^{4m})\|m`$ $`T^+\|m`$ $`=`$ $`{\displaystyle \frac{1}{q^2}}\sqrt{{\displaystyle \frac{d}{\lambda }}}\sqrt{(q^{2m}q^{2\overline{m}})\left(q^{2m}{\displaystyle \frac{1}{\lambda d}}q^{2(\overline{m}+1)}\right)}\|m+1`$ (A.21) $`T^{}\|m`$ $`=`$ $`\sqrt{{\displaystyle \frac{d}{\lambda }}}\sqrt{(q^{2(m1)}q^{2\overline{m}})\left(q^{2(m1)}{\displaystyle \frac{1}{\lambda d}}q^{2(\overline{m}+1)}\right)}\|m1`$ $`\tau \|m`$ $`=`$ $`d\lambda q^{4m}\|m`$ $`\tau `$ has negative eigenvalues only, thus $`\tau ^{\frac{1}{2}}`$ and $`\stackrel{}{T}^2`$ will not be real. ## B Representations of the $`𝑲`$ Algebra The algebraic relations of the $`K`$ algebra are the same as the relations of the $`T`$ algebra, they are different only as a $``$ algebra: $`\overline{K^3}=K^3,`$ $`\overline{K^+}={\displaystyle \frac{1}{q^2}}K^{}`$ (B.1) $`\overline{T^3}=T^3,`$ $`\overline{T^+}={\displaystyle \frac{1}{q^2}}T^{}`$ This makes the $`K`$ algebra a $`su_q(1,1)`$ algebra. All the results that depend only on the algebraic relations are the same as for the $`T`$ algebra. $`\text{(}\text{A.1}\text{)}:`$ $`K^3\|m`$ $`=\varphi (m)\|m`$ (B.2) $`\text{(}\text{A.2}\text{)}:`$ $`K^+\|m`$ $`=\gamma _m\|m`$ (B.3) $`\text{(}\text{A.4}\text{)}:`$ $`\varphi (m)`$ $`={\displaystyle \frac{1}{\lambda }}d_kq^{4m}`$ (B.4) We again take $`d_k`$ and $`m`$ real. $$\text{(}\text{A.5}\text{)}:\tau _k\|m=\lambda d_kq^{4m}\|m$$ (B.5) For $`K^{}`$ there is a change in sign due to (B.1): $`\text{(}\text{A.6}\text{)}:`$ $`K^{}\|m`$ $`=q^2\gamma _{m1}^{}\|m1`$ (B.6) $`\text{(}\text{A.8}\text{)}:`$ $`\gamma _m^{}\gamma _m`$ $`={\displaystyle \frac{1}{\lambda }}\left\{{\displaystyle \frac{1}{q^2\lambda }}+\alpha \lambda q^{2m}{\displaystyle \frac{d_k}{q^4}}q^{4m}\right\}`$ (B.7) Now $`\gamma _m^{}\gamma _m`$ becomes positive for $`m\mathrm{}`$, we do not have to cut off the spectrum at a largest value of $`m`$. We shall see that all the representations are infinite-dimensional. We introduce the function $`\kappa (x)`$ analogous to $`h(x)`$ in (A.10): $$\kappa (x)=\left\{\frac{1}{q^2\lambda ^2}\alpha x+\frac{d_k}{\lambda q^4}x^2\right\},x=q^{2m}.$$ (B.8) The representations of the $`K`$ algebra are: $`K^3\|m`$ $`=`$ $`({\displaystyle \frac{1}{\lambda }}d_kq^{4m})\|m`$ $`K^+\|m`$ $`=`$ $`\sqrt{\kappa (q^{2m})}\|m+1`$ (B.9) $`K^{}\|m`$ $`=`$ $`q^2\sqrt{\kappa (q^{2(m1)})}\|m1`$ They are characterized by $`\alpha `$ and $`d_k`$ and restricted by the condition $`\kappa (q^{2m})0`$. To discuss this condition we determine the zeros of $`\kappa (x)`$ $`\kappa (x_{1,2})`$ $`=`$ $`0`$ (B.10) $`x_{1,2}`$ $`=`$ $`{\displaystyle \frac{\lambda }{2d_kq^4}}\left\{\alpha \pm \sqrt{\alpha ^24d_kq^6\lambda ^3}\right\}`$ We discuss the cases $`d_k>0`$, $`d_k=0`$ and $`d_k<0`$ separately and start with $`d_k>0`$: In this case $`\kappa (x)`$ has no positive zero for $`\alpha <2q^3\sqrt{d_k\lambda ^3}=\alpha _0`$. The range of $`m`$ is not restricted, it can be of the form $`m_0+n`$, $`n_0`$. If $`\alpha \alpha _0`$ we will have two positive zeros and $`\kappa (x)`$ can be written in the form $$\kappa (x)=\frac{1}{q^2\lambda ^2x_1x_2}(xq^{2\overline{m}})(xq^{2\underset{¯}{m}}).$$ (B.11) The values of the zeros $`x_1`$, $`x_2`$ determine the parameters $`\alpha `$ and $`d_k`$ and therefore the representation. $`d_k`$ $`=`$ $`{\displaystyle \frac{q^2}{\lambda }}q^{2(\overline{m}+\underset{¯}{m})}`$ (B.12) $`\alpha `$ $`=`$ $`{\displaystyle \frac{1}{q^2\lambda ^2}}(q^{2\underset{¯}{m}}+q^{2\overline{m}})`$ There are now two inequivalent representations. We find that in one representation the allowed values of $`m`$ are $$m\overline{m},m=\overline{m},\overline{m}1,\overline{m}2,\mathrm{}$$ (B.13) For the other representation we find $$m>\underset{¯}{m},m=\underset{¯}{m}+1,\underset{¯}{m}+2,\mathrm{}$$ (B.14) Now we consider $`d_k=0`$: The function $`\kappa `$ becomes linear. It is positive at $`x=0`$ and, depending on $`\alpha `$, stays positive or becomes negative for $`x\mathrm{}`$. If $`\alpha <0`$ there is no restriction in the range of $`m`$, $`m=m_0+n,`$ $`n`$. If $`\alpha >0`$ there is a lowest eigenvalue of $`m`$, we are at the situation of (B.14). Finally we consider $`d_k<0`$: In this case $`\kappa (x)`$ is positive for $`x0`$ and negative for $`x\mathrm{}`$. There is one zero for $`x>0`$. This can also be seen from (B.10) because the square root will now be larger than $`\alpha `$. The relevant zero of $`\kappa (x)`$ is: $`x_1`$ $`=`$ $`{\displaystyle \frac{\lambda }{2\|d_k\|q^4}}\left\{\alpha +\sqrt{\alpha ^2+4\|d_k\|q^6\lambda ^3}\right\}`$ $`=`$ $`q^{2\underset{¯}{m}_k}`$ Now all values of $`\alpha `$ are allowed. The range of $`m`$ will be as in (B.14). For orbital angular momentum we encounter the representation $$d_k=\frac{1}{q^2\lambda },x_1=q^2.$$ (B.16) This leads to $`\alpha =0`$ and $$\kappa (q^{2m})=\frac{1}{\lambda ^2q^2}\left(1q^{4(m+1)}\right).$$ (B.17) The respective representation is shown in (5.2). ## C Comultiplication The standard comultiplication rule for the algebra (2.3) is: $`\mathrm{\Delta }(T^3)`$ $`=`$ $`T^31+\tau T^3`$ (C.1) $`\mathrm{\Delta }(T^\pm )`$ $`=`$ $`T^\pm 1+\tau ^{\frac{1}{2}}T^\pm `$ As a consequence, $`\tau `$ is group-like: $$\mathrm{\Delta }(\tau )=\tau \tau .$$ (C.2) The algebra (2.3) is the same for the $`T`$ algebra and the $`K`$ algebra, they are distinguished by their conjugation properties (B.1). As long as $`\tau ^{\frac{1}{2}}`$ is hermitean, (C.1) will respect the conjugation properties and we have a comultiplication within the $`T`$ algebra or the $`K`$ algebra respectively. From (C.2) follows that $`\mathrm{\Delta }(\tau ^{\frac{1}{2}})`$ will be hermitean if $`\tau ^{\frac{1}{2}}`$ is. $$\mathrm{\Delta }(\tau ^{\frac{1}{2}})=\tau ^{\frac{1}{2}}\tau ^{\frac{1}{2}}$$ (C.3) If $`\tau ^{\frac{1}{2}}`$ is not hermitean $`\mathrm{\Delta }(T)`$ will have no definite conjugation properties even if $`T`$ has. We now turn to the product of representations as it follows from the comultiplication rule (C.1). If we have two representations of the algebra (2.3) we obtain a new one by the rule $`\mathrm{\Delta }(T^3)`$ $`=`$ $`T_1^31+\tau _1T_2^3`$ (C.4) $`\mathrm{\Delta }(T^\pm )`$ $`=`$ $`T_1^\pm 1+\tau _1^{\frac{1}{2}}T_2^\pm `$ From the discussion above follows that we can multiply two representations of the $`T`$ algebra ($`K`$ algebra) to obtain a representation of the $`T`$ algebra ($`K`$ algebra) as long as $`\tau _1^{\frac{1}{2}}`$ is hermitean. From now on we shall drop the indices $`1`$ and $`2`$ again, first and second representations will be defined by the position in the product (C.4). That the $`\tau ^{\frac{1}{2}}`$ of the first representation is hermitean means $`d_1>0`$. We shall discuss this situation first. $`d_1>0`$: The product of two representations of the $`T`$ algebra ($`K`$ algebra) will be a $`T`$ algebra ($`K`$ algebra). From (C.2) follows $$d=\lambda d_1d_2.$$ (C.5) If $`d_2`$ is negative $`d`$ will be negative as well. For the $`T`$ algebra $`d`$ positive restricts $`d`$ to be $`d=\frac{1}{\lambda }`$. This characterizes the finite-dimensional representations. From (C.5) follows that the product of two finite-dimensional representations is finite-dimensional as expected but also that the product of a finite-dimensional representation $`(d_1=\frac{1}{\lambda })`$ with an infinite-dimensional representation $`(d_2<0)`$ leads to $`d<0`$ and cannot be reduced to finite-dimensional representations. For the $`K`$ algebra all representations are infinite-dimensional. We now turn to the case that $`d_1`$ is negative, $`\tau _1^{\frac{1}{2}}`$ will be anti-hermitean. $`d_1<0`$: In this case the product of two representations will in general not have well-defined conjugation properties. We can, however, start from a modified comultiplication rule: $`\mathrm{\Delta }_\beta (T^3)`$ $`=`$ $`T^31+\tau T^3`$ (C.6) $`\mathrm{\Delta }_\beta (T^\pm )`$ $`=`$ $`T^\pm 1\pm (\tau )^{\frac{1}{2}}T^\pm `$ If $`(\tau )^{\frac{1}{2}}`$ is hermitean this rule allows us to multiply a representation of the $`T(K)`$ algebra by a representation of the $`K(T)`$ algebra to obtain a $`T(K)`$ algebra. $`T^31+\tau K^3`$ (C.7) $`T^\pm 1\pm (\tau )^{\frac{1}{2}}K^\pm `$ will be a representation of the $`T`$ algebra whereas $`K^31+\tau _kT^3`$ (C.8) $`K^\pm 1\pm (\tau _k)^{\frac{1}{2}}T^\pm `$ will be a representation of the $`K`$ algebra. For the comultiplication (C.6) $`\tau `$ will be group-like as well and it follows again that $$d=\lambda d_1d_2.$$ (C.9) But now $`d_1`$ is negative. Of special interest is the case that $`d_1`$ and $`d_2`$ are both negative, then $`d`$ is positive. If we multiply $`T\times K`$ to obtain a $`T`$ algebra then we know that $`d=\frac{1}{\lambda }`$ and, as a consequence $$d_1d_2=\frac{1}{\lambda ^2}$$ (C.10) to obtain a representation with well-defined conjugation properties. This is exactly the case for the construction of the $`T_{orb}`$ algebra in the main part of this paper. ## D The big $`𝒒`$-Jacobi polynomials In this appendix we recall some basics about $`q`$-special functions , , , in particular the big $`q`$-Jacobi polynomials. First, we introduce some useful notation. The expressions $$[a]=\frac{q^aq^a}{qq^1}\stackrel{q1}{}a,[a]!=\underset{k=1}{\overset{a}{}}[k]\stackrel{q1}{}a!$$ (D.1) are known as symmetric $`q`$-numbers and symmetric $`q`$-factorials respectively. The corresponding $`q`$-binomial coefficient is $$\left[\begin{array}{c}n\hfill \\ k\hfill \end{array}\right]=\{\begin{array}{cc}\frac{[n]!}{[k]![nk]!}\hfill & \text{ for }nk,\hfill \\ & \\ 0\hfill & \text{ for }n<k\text{ or }n,k<0.\hfill \end{array}$$ (D.2) Of course $`\left[\begin{array}{c}n\hfill \\ k\hfill \end{array}\right]\stackrel{q1}{}\left(\begin{array}{c}n\hfill \\ k\hfill \end{array}\right)`$. There are also “unsymmetric” counterparts of these objects: the basic $`q`$-number $$\frac{1q^a}{1q}\stackrel{q1}{}a$$ (D.3) and the $`q`$-shifted factorial (Pochammer-symbol) $$(a;q)_k=\underset{n=0}{\overset{k1}{}}(1aq^n),(a_1,a_2,\mathrm{},a_i;q)_k=\underset{m=1}{\overset{i}{}}(a_m;q)_k.$$ (D.4) The Jackson integral of a function $`f(x)`$ is defined for $`q>1`$ by $$_0^ad_{q^1}xf(x)=(1q^1)\underset{\nu =0}{\overset{\mathrm{}}{}}aq^\nu f(aq^\nu )$$ (D.5) With the help of the $`q`$-shifted factorials, the basic hypergeometric function can be introduced $`{}_{r}{}^{}\varphi _{s}^{}\left(\begin{array}{c}a_1,\mathrm{},a_r\hfill \\ b_1,\mathrm{},b_s\hfill \end{array}\|q^1;x\right)=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(a_1,\mathrm{},a_r;q^1)_k}{(b_1,\mathrm{},b_s;q^1)_k}}(1)^{(1+sr)k}q^{\frac{1}{2}(1+sr)k(k1)}{\displaystyle \frac{x^k}{(q^1;q^1)_k}}`$ This series plays in the theory of $`q`$-special functions a role analogous to that of the hypergeometric series in the theory of usual special functions. We have considered a base $`q^1`$ here, because in this way $`{}_{r}{}^{}\varphi _{s}^{}`$ is well-defined for $`q>1`$, which is the case we are interested in here. The big $`q`$-Jacobi polynomials , are constructed in terms of the basic hypergeometric series as $$P_l(x;a,b,c;q^1)={}_{3}{}^{}\varphi _{2}^{}\left(\begin{array}{c}q^l,abq^{(l+1)},x\hfill \\ aq^1,cq^1\hfill \end{array}\|q^1;q^1\right).$$ (D.9) For the applications we consider in this paper, we are interested in the case $`P_l^m(x)P_{lm}(x;q^{2m},q^{2m},q^{2m};q^2),m0`$ (D.17) $`={\displaystyle \underset{k=0}{\overset{lm}{}}}(1)^k{\displaystyle \frac{q^{k(m+1)}(x;q^2)_k}{(q^{2(m+1)};q^2)_k}}\left[\begin{array}{c}lm\hfill \\ k\hfill \end{array}\right]\left[\begin{array}{c}l+m+k\hfill \\ k\hfill \end{array}\right]\left[\begin{array}{c}m+k\hfill \\ k\hfill \end{array}\right]^1`$ Notice that the $`P_l^m`$ are polynomials of the order $`lm`$ in $`x`$. Due to the factor $`\left[\begin{array}{c}lm\hfill \\ k\hfill \end{array}\right]`$, which vanishes for $`k>lm`$ according to the definition (D.2). The sum in (D) actually becomes finite. Due to the same factor the polynomials $`P_l^m`$ vanish, if the condition $`ml`$ is not satisfied: $$P_l^m(x)=0\text{for }l<m.$$ (D.18) The further condition $`m0`$ is necessary for the polynomials $`P_m^l`$ to be well-defined, due to the factor $`(q^{2(m1)};q^2)_k`$ in the denominator of the basic hypergeometric function which otherwise vanishes for negative $`m`$. Some of the first big $`q`$-Jacobi polynomials are $$\begin{array}{ccc}P_0^0(x)=1,\hfill & P_1^0(x)=x,\hfill & \\ & & \\ P_2^0(x)=\frac{1}{q[2]}([3]x^2q^2),\hfill & P_3^0(x)=\frac{x}{q^5[2]}([5]q^2x^2[3]),\hfill & \\ & & \\ P_1^1(x)=1,\hfill & P_2^1(x)=x,\hfill & \\ & & \\ P_3^1(x)=\frac{1}{q^5[4]}(q^4[5]x^21).\hfill & & \end{array}$$ (D.19) From Ref. and we learn that the polynomials $`P_l^m(x)`$ satisfy a recurrence relation $$xq^m[2l+1]P_l^m(x)=q^l[l+m+1]P_{l+1}^m(x)+q^{l1}[lm]P_{l1}^m(x)$$ (D.20) a $`q`$-difference equation $`\left(q^{12m}(q^{2l+1}+q^{2l1})x^2q^{4(m+1)}(q^2+1)\right)P_l^m(x)=`$ $`q^{2(2m+1)}(x^21)P_l^m(xq^2)+(x^2q^{4(m+1)})P_l^m(xq^2)`$ and the orthonormality condition $$_{q^{2(m+1)}}^{q^{2(m+1)}}d_{q^2}xw_l^m(x)w_l^{}^m(x)P_l^m(x)P_l^{}^m(x)=\delta _{l,l^{}}.$$ (D.22) Here, the weight function $`w_l^m`$ is defined by $`w_l^m(x)`$ $``$ $`\sqrt{{\displaystyle \frac{(q^{4(m+1)};q^4)_{\mathrm{}}}{(q^4,q^{4(m+1)2};q^4)_{\mathrm{}}(q^2;q^2)_{\mathrm{}}}}}\sqrt{{\displaystyle \frac{[2m+1]}{2[2l+1]}}}`$ $`\times q^{\frac{1}{2}(l^2+l+2lm3m^2+m+3)}\sqrt{(x^2q^{4m};q^4)_m}\sqrt{{\displaystyle \frac{(q^2;q^2)_{lm}}{(q^{2(2m+1)};q^2)_{lm}}}}`$ Actually, as it is done e.g. in , the big $`q`$-Jacobi polynomials can be alternatively defined as those polynomials in $`x`$ which are orthonormal with respect to the Jackson integral with the weight function $`w_l^m(x)`$ in (D). The weight function has the following scaling properties $`w_l^m(xq^2)`$ $`=`$ $`w_l^m(x)\sqrt{{\displaystyle \frac{(1x^2)}{(1x^2q^{4m})}}},`$ (D.24) $`w_{l1}^m(x)`$ $`=`$ $`w_l^m(x)q^l\sqrt{{\displaystyle \frac{[l+m][2l+1]}{[lm][2l1]}}}.`$ It is useful for the purposes of this paper to absorb the weight function in the definition of the polynomials themselves and to introduce the further notation $$\stackrel{~}{P}_l^m(x)=w_l^m(x)P_l^m(x).$$ (D.25) With the help of (D.24) it turns out that (D.20) and (D) become respectively $$xq^{m+1}\stackrel{~}{P}_l^m(x)=\sqrt{\frac{[lm+1][l+m+1]}{[2l+1][2l+3]}}\stackrel{~}{P}_{l+1}^m(x)+\sqrt{\frac{[l+m][lm]}{[2l+1][2l1]}}\stackrel{~}{P}_{l1}^m(x)$$ (D.26) and $`\left((q^{2l+1}+q^{2l1})q^1x^2(q^2+1)q^{2(m+2)}\right)\stackrel{~}{P}_l^m(x)=`$ $`q^{2(m+1)}\sqrt{(x^21)(x^2q^{4m}1)}\stackrel{~}{P}_l^m(xq^2)`$ $`+\sqrt{(x^2q^{4(m+1)+1})(x^2q^4)}\stackrel{~}{P}_l^m(xq^2)`$ By using (D.22) and the definition of the Jackson integral (D.5) we obtain the following orthonormality condition for the functions $`\stackrel{~}{P}_l^m(x)`$ $$(1q^2)\underset{\sigma =\pm 1}{}\underset{n=\mathrm{}}{\overset{0}{}}q^{2(nm1)}\stackrel{~}{P}_l^m(\sigma q^{2(nm1)})\stackrel{~}{P}_l^{}^m(\sigma q^{2(nm1)})=\delta _{l,l^{}}.$$ (D.28) Moreover, the functions $`\stackrel{~}{P}_l^m(x)`$ have the property that they transform under a parity transformation like $$\stackrel{~}{P}_l^m(x)=(1)^{lm}\stackrel{~}{P}_l^m(x).$$ (D.29) In the particular case $`m=0`$ the big $`q`$-Jacobi polynomials become the big $`q`$-Legendre polynomials, which in the limit $`q1`$ yield the usual Legendre polynomials. In the same limit from the polynomials $`P_l^m(x)`$ we recover the Jacobi polynomials with the normalization $`P_l^m(1)=1`$. ## E Diagonalization of $`𝑿^\mathrm{𝟑}`$ In this appendix we study the transformation which is inverse to the transformation (6.5), (6.9) constructed in Section 5. We show how the big $`q`$-Jacobi polynomials can be used to diagonalize $`X^3`$ in the basis where $`\stackrel{}{T}_{orb}^2`$, $`T_{orb}^3`$, $`R^2`$ are diagonal. The representation where $`\stackrel{}{T}_{orb}^2`$, $`T_{orb}^3`$, $`R^2`$ are diagonal can be found in , , $`\stackrel{}{T}_{orb}^2\|M,l,m`$ $`=`$ $`q[l][l+1]\|M,l,m`$ $`X^3\|M,l,m`$ $`=`$ $`r_0q^{2M+m}\{\sqrt{{\displaystyle \frac{[l+m+1][lm+1]}{[2l+1][2l+3]}}}\|M,l+1,m`$ (E.1) $`+\sqrt{{\displaystyle \frac{[l+m][lm]}{[2l+1][2l1]}}}\|M,l1,m\}`$ $`X^+\|M,l,m`$ $`=`$ $`r_0q^{2M+m}\{q^l\sqrt{{\displaystyle \frac{[l+m+1][l+m+2]}{[2][2l+1][2l+3]}}}\|M,l+1,m+1`$ $`q^{l+1}\sqrt{{\displaystyle \frac{[lm][lm1]}{[2][2l+1][2l1]}}}\|M,l1,m+1\}`$ $`X^{}\|M,l,m`$ $`=`$ $`r_0q^{2M+m}\{q^l\sqrt{{\displaystyle \frac{[lm+1][lm+2]}{[2][2l+1][2l+3]}}}\|M,l+1,m1`$ $`q^{l1}\sqrt{{\displaystyle \frac{[l+m][l+m1]}{[2][2l+1][2l1]}}}\|M,l1,m1\}`$ where $$0l<\mathrm{},lml.$$ (E.2) We make the following Ansatz for an eigenfunction of $`X^3`$ $$X^3\underset{M,l,m}{}d_{M,l,m}\|M,l,m=z\underset{M,l,m}{}d_{M,l,m}\|M,l,m,$$ (E.3) with $`z`$ the corresponding eigenvalue. By using (E) we obtain a recursion relation for the coefficients $`d_{M,l,m}`$ $`zd_{M,l,m}`$ $`=`$ $`r_0{\displaystyle \frac{q^{2M+m}}{\sqrt{[2l+1]}}}\{\sqrt{{\displaystyle \frac{[lm+1][l+m+1]}{[2l+3]}}}d_{M,l+1,m}`$ $`+\sqrt{{\displaystyle \frac{[l+m][lm]}{[2l1]}}}d_{M,l1,m}\}.`$ A comparison with the recursion relation (D.26) for the functions $`\stackrel{~}{P}_l^m`$ defined in (D.25) in terms of the Jacobi polynomials shows that a solution of (E) is $$d_{M,l,m}^{\nu ,\sigma }=\{\begin{array}{cc}\sqrt{1q^2}q^{\nu M1m}\stackrel{~}{P}_l^m(\sigma q^{2(\nu M1m)})\hfill & \text{ for }m0\hfill \\ \sqrt{1q^2}q^{\nu M1}\stackrel{~}{P}_l^{\|m\|}(\sigma q^{2(\nu M1)})\hfill & \text{ for }m<0\hfill \end{array}$$ (E.5) where $`z=\sigma r_0q^{1+2\nu }`$, $`\sigma =\pm 1`$. By comparing with the form of the eigenvalues of $`X^3`$ (5.8) we see that we have to restrict $$\nu ,M,\nu M,m\nu M.$$ (E.6) Notice that the argument of the functions would correspond to $`x=\mathrm{cos}\theta =\frac{z}{r}`$ classically, apart from the $`q`$-factor $`q^{(m+\|m\|)}1`$ for $`q1`$. $`X^3`$ is a self-adjoint operator in this representation. This was shown in Ref. . Now, the set of eigenfunctions of a self-adjoint operator is complete, therefore we expect a completeness relation to hold for the eigenfunctions of $`X^3`$. In fact, (D.28) can be interpreted in this way. As the sum (D.28) contains two sums, one where the argument of $`P_m^l`$ is positive and one where it is negative, we obtain a representation where the eigenvalues of $`X^3`$ can have both signs, so that we automatically find the direct sum of two representations of the type (5.8). The normalization of the coefficients $`d_{M,l,m}`$ in (E.5) has been chosen in such a way as to yield exactly (D.28). As the eigenfunctions of a selfadjoint operator corresponding to different eigenvalues are orthogonal, since the normalization constant is already fixed by (D.28), we argue that the following relation holds $$(1q^2)\underset{l=0}{\overset{\mathrm{}}{}}q^{\nu +\nu ^{}2}\stackrel{~}{P}_l^{\|m\|}(\sigma q^{2(\nu 1)})\stackrel{~}{P}_l^{\|m\|}(\sigma ^{}q^{2(\nu ^{}1)})=\delta _{\nu ,\nu ^{}}\delta _{\sigma ,\sigma ^{}}$$ (E.7) where $`\sigma ,\sigma ^{}=\pm 1`$ are the signs of the argument of the functions and $`\nu \mathrm{min}\{m,0\}`$. This is an interesting result for itself about the Jacobi polynomials.
warning/0004/cond-mat0004057.html	ar5iv	text	# Spin 1/2 Magnetic Impurity in a 2D Magnetic System Close to Quantum Critical Point. ## I Introduction There is no need to explain the importance of the Kondo problem for condense matter physics. It is relevant to the magnetic impurities in metals, heavy fermion compounds , tunneling phenomena in quantum dots , correlated lattice fermion systems , and many other physical systems. The problem of magnetic impurity in a two-dimensional (2D) insulating system with long range antiferromagnetic order has also attracted great interest. This includes an impurity spin with an on-site and sublattice symmetric coupling, as well as an isolated ferromagnetic bond . Because of Adler’s relation for the impurity-spin-wave interaction these systems have no nontrivial infrared dynamics at zero temperature in spite of the gapless spectrum of Goldstone spin waves . This makes the impurity problem for the insulating state much simpler than the Kondo one. However it has been realized recently that in the case when the 2D magnetic insulating system is close to the quantum critical point the infrared dynamics of the impurity is highly nontrivial and to a large extend is similar to that for the Kondo problem. The quantum critical impurity problem has been very recently addressed by Vojta, Buragohain, and Sachdev . They demonstrated that there are nontrivial critical indexes for the impurity Green’s function and for magnetic susceptibilities. In the present work we consider the same quantities calculating more accurately the critical indexes and in some cases the prefactors. In addition we consider specific heat and interaction of the impurities . ## II The Hamiltonian To be specific we consider a magnetic impurity in two coupled Heisenberg planes. The two coupled Heisenberg planes is an ideal example of a 2D critical system which can be described by O(3) nonlinear $`\sigma `$-model. The two planes model is very well studied at zero temperature both numerically and analytically . Finite temperature properties are also well understood . Hamiltonian of the system with impurity is of the form $`H=H_2+H_{imp},`$ (1) $`H_2=J{\displaystyle \underset{<ij>}{}}\left(𝐒_i^{(1)}𝐒_j^{(1)}+𝐒_i^{(2)}𝐒_j^{(2)}\right)+J_{}{\displaystyle \underset{i}{}}𝐒_i^{(1)}𝐒_i^{(2)},`$ (2) $`H_{imp}=j𝐒_0^{(1)}𝐬.`$ (3) Here $`𝐒_𝐢^{(𝐧)}`$ is spin 1/2 on the square lattice. the index $`i`$ numerates cites, and the index $`n`$ numerates planes. $`J`$ is an antiferromagnetic coupling in the plane, and $`J_{}`$ is an antiferromagnetic coupling between the planes. Spin of the impurity, $`s=1/2`$, is coupled to one of the planes. It is known that in this system at $`J_{}=2.525\pm 0.002J`$ there is a quantum phase transition from quantum disordered state to the Neel state. In the present work we consider only the quantum disordered phase including the critical point. As we have already pointed out an interesting regime arises only close to the critical point and therefore we concentrate on the vicinity of this point. Note that in terms of the non-linear $`O(3)`$ $`\sigma `$-model an effective Lagrangian of the system is $`L=_\mu \stackrel{}{\phi }_\mu \stackrel{}{\phi }m^2\stackrel{}{\phi }^2+\gamma \stackrel{}{\phi }\stackrel{}{s}`$, $`\stackrel{}{\phi }^2=a^2`$, where the parameters can be expressed in terms of parameters of the original Hamiltonian (1). In our considerations we will use only the original Hamiltonian (1). Using bond operator representation $$𝐒^{(1,2)}=\frac{1}{2}\left(\pm 𝐭\pm 𝐭^{}i𝐭^{}\times 𝐭\right),$$ (4) the two plane Hamiltonian $`H_2`$ from (1) can be rewritten in terms of the operators $`𝐭_𝐢=(t_{i,x},t_{i,y},t_{i,z})`$, and then diagonalized by a combination of the Fourier and Bogoliubov transformations with account of the hard core constraint, see Refs. : $`𝐭_i={\displaystyle \underset{𝐪}{}}e^{i\mathrm{𝐪𝐫}_𝐢}(u_𝐪t_𝐪+v_𝐪𝐭_𝐪^{}),`$ (5) $`H_2{\displaystyle \underset{𝐪}{}}\omega _𝐪𝐭_𝐪^{}𝐭_𝐪,`$ (6) where $`u_𝐪,v_𝐪=\sqrt{\frac{A_𝐪}{2\omega _𝐪}\pm \frac{1}{2}}`$ are Bogoliubov parameters. The operator $`𝐭_𝐪^{}`$ creates quasiparticle of the system. This quasiparticle has spin 1 and we call it spin wave or magnon. Near the critical point the spin-wave excitation energy is $$\omega _𝐪\sqrt{\mathrm{\Delta }^2+c^2q^2},$$ (7) where $`\mathrm{\Delta }J`$ is the spin-wave gap and $`c1.9J`$ is the spin-wave velocity. The function $`A_𝐪`$ in the vicinity of the critical point is q-independent: $`A_𝐪A2.4J`$. In the bond operator representation the impurity Hamiltonian is of the form $$H_{imp}=\frac{j}{2}𝐬\left(𝐭_\mathrm{𝟎}+𝐭_{\mathrm{𝟎}}^{}{}_{}{}^{}i𝐭_{\mathrm{𝟎}}^{}{}_{}{}^{}\times 𝐭_\mathrm{𝟎}\right)\frac{j}{2}𝐬\left(𝐭_\mathrm{𝟎}+𝐭_{\mathrm{𝟎}}^{}{}_{}{}^{}\right).$$ (8) Here we have dropped the term $`𝐭_{\mathrm{𝟎}}^{}{}_{}{}^{}\times 𝐭_\mathrm{𝟎}`$. The matter is that we are interested in nontrivial long-range dynamics, but one can prove that all diagrams generated by the $`𝐭_{\mathrm{𝟎}}^{}{}_{}{}^{}\times 𝐭_\mathrm{𝟎}`$ term are infrared convergent, and therefore its contribution is less important. An alternative way is to redefine $`H_{imp}`$ as $`H_{imp}\frac{j}{2}𝐬\left(𝐒_{\mathrm{𝟎}}^{}{}_{}{}^{(1)}𝐒_{\mathrm{𝟎}}^{}{}_{}{}^{(2)}\right)`$, then the $`𝐭_{\mathrm{𝟎}}^{}{}_{}{}^{}\times 𝐭_\mathrm{𝟎}`$ term is canceled out exactly. Note that for an integer impurity spin the $`𝐭_{\mathrm{𝟎}}^{}{}_{}{}^{}\times 𝐭_\mathrm{𝟎}`$ term could be much more important because it can give a bound state of the spin wave with the impurity, and hence can lead to the full screening of the impurity. Using (5) we rewrite $`H_{imp}`$ in terms of quasiparticle operators $`𝐭_𝐪`$. $$H_{imp}=\frac{j}{4}\underset{𝐪}{}(u_𝐪+v_𝐪)\stackrel{}{t}_𝐪^{}\stackrel{}{\sigma }+h.c.\frac{j\sqrt{A}}{2\sqrt{2}}\underset{𝐪}{}\frac{1}{\sqrt{\omega _𝐪}}\stackrel{}{t}_𝐪^{}\stackrel{}{\sigma }+h.c.,$$ (9) where $`\stackrel{}{\sigma }`$ is the impurity Pauli matrix. Hereafter we set $`J=1`$, so all energies are measured in units of $`J`$. ## III The impurity Green’s function at zero temperature Let us calculate the impurity self energy $`\mathrm{\Sigma }`$ shown in Fig. 1. In single loop approximation, see Fig.1a, using eq. (9) we find $$\mathrm{\Sigma }^{(1)}(ϵ)=\frac{3Aj^2}{8}\frac{1}{\omega _𝐪(ϵ\omega _𝐪)}\frac{d^2q}{(2\pi )^2}.$$ (10) Close to the critical point $`\omega _𝐪=\sqrt{c^2q^2+\mathrm{\Delta }^2}`$, and hence simple integration gives for $`ϵ=0`$ (i.e. at the position of the quasiparticle pole) $$\mathrm{\Sigma }^{(1)}(0)=\alpha ^2\mathrm{ln}\frac{\mathrm{\Lambda }}{\mathrm{\Delta }}.$$ (11) Here $$\alpha ^2=\frac{3Aj^2}{16\pi c^2}0.04j^2,$$ (12) is dimensionless coupling constant, and $`\mathrm{\Lambda }2J`$ is the ultraviolet cutoff. Similar calculation for the second order self energy shown in Fig.1b gives $$\mathrm{\Sigma }^{(2)}(0)=\frac{\alpha ^4}{\mathrm{\Delta }}\mathrm{ln}\frac{\mathrm{\Lambda }}{\mathrm{\Delta }}.$$ (13) As soon as $`\alpha ^2/\mathrm{\Delta }1`$ the second order self energy is small compared to the first order one, $`\mathrm{\Sigma }^{(2)}\mathrm{\Sigma }^{(1)}`$, and hence the perturbation theory is justified. However at $`\alpha ^2/\mathrm{\Delta }>1`$ the expansion does not converge and hence one has to sum all orders of perturbation theory. Exactly at the critical point, $`\mathrm{\Delta }=0`$, the expansion diverges at arbitrary small $`\alpha `$. The problem under consideration has a small parameter which is independent of the interaction. This is 1/N, where N=3 is number of components of the spin-wave excitation (O(N) $`\sigma `$-model). In the leading in N approximation only the rainbow diagrams contribute to the impurity self energy. Summation of these diagrams leads to the usual noncrossing approximation (= self consistent Born approximation) for the impurity Green’s function $$G(ϵ)=\frac{1}{ϵ\mathrm{\Sigma }(ϵ)},$$ (14) $$\mathrm{\Sigma }(ϵ)=\alpha ^2G(ϵ\omega )𝑑\omega .$$ (15) It is convenient to rescale the variables: $`ϵϵ/\alpha ^2`$, $`\omega \omega /\alpha ^2`$, $`G\alpha ^2G`$, $`\mathrm{\Sigma }\mathrm{\Sigma }/\alpha ^2`$. In the new variables $`\alpha `$ disappears from the eqs. (14,15). Dependence on $`\alpha `$ remains only in the limits of $`\omega `$-integration: $`\mathrm{\Sigma }=_{\mathrm{\Delta }_\alpha }^{\mathrm{\Lambda }_\alpha }G𝑑\omega `$, where $`\mathrm{\Delta }_\alpha =\mathrm{\Delta }/\alpha ^2`$ and $`\mathrm{\Lambda }_\alpha =\mathrm{\Lambda }/\alpha ^2`$. Consider first the critical point, i.e. $`\mathrm{\Delta }_\alpha =0.`$ In this case eqs. (14,15) can be solved analytically. The answer is $$G(ϵ)=\frac{1}{\sqrt{2(ϵ_0ϵ)}},$$ (16) where $`ϵ_0`$ is the impurity binding energy. Eq. (16) is valid if $`\|ϵ_0ϵ\|1`$. For illustration we present in Fig. 2 spectral function, $`\frac{1}{\pi }\text{Im}G(ϵ)`$, obtained by direct numerical solution of eqs. (14,15) at $`\mathrm{\Delta }_\alpha =0`$ and $`\mathrm{\Lambda }_\alpha =1,2,5`$. Agreement with analytical solution (16) is perfect. At the large $`\mathrm{\Lambda }_\alpha =\mathrm{\Lambda }/\alpha ^2`$ the impurity binding energy $`ϵ_02.5`$. In the original variables it means that $`ϵ_02.5\alpha ^2`$, when $`\alpha ^21`$. As one shall expect the Green’s function (16) has no quasiparticle pole. However the pole appears away from the critical point when the spin-wave gap is nonzero . For illustration we present in Fig. 3 the spectral functions obtained by numerical solution of Eqs. (14,15) for $`\mathrm{\Lambda }_\alpha =2`$, $`\mathrm{\Delta }_\alpha =0`$, and $`\mathrm{\Lambda }_\alpha =2`$, $`\mathrm{\Delta }_\alpha =0.05`$. It is clear that the quasiparticle peak absorbs spectral weight of the incoherent Green’s function (16) from the area $`ϵ_0ϵ\mathrm{\Delta }_\alpha `$. Therefore the quasiparticle residue $`Z\sqrt{\mathrm{\Delta }_\alpha }`$ (c.f. with Ref. ). Slightly more detail analysis of the Eqs. (14,15) shows that $$Z0.8\sqrt{\mathrm{\Delta }_\alpha }.$$ (17) We have considered above the leading in N approximation. Let us estimate now $`1/N`$ correction which is due to the single loop contribution to the impurity-magnon vertex function shown in Fig.4. Taking into account the algebraic relation for the Pauli matrices, $`\sigma _\mu \sigma _\nu \sigma _\mu =\sigma _\nu `$, we find the vertex function given by Fig. 4 $$\mathrm{\Gamma }(ϵ,\lambda )=\mathrm{\Gamma }_{bare}\left(1\frac{1}{N}G(ϵ\omega )G(ϵ\omega \lambda )𝑑\omega \right).$$ (18) Here $`ϵ`$ and $`\lambda `$ are energies of incoming impurity and spin wave correspondingly, and $`\mathrm{\Gamma }_{bare}`$ is the bare vertex given by eq. (9). Taking $`\lambda ϵϵ_0`$ and using Green’s function (16) we find after integration in (18) $$\mathrm{\Gamma }_\lambda =\mathrm{\Gamma }_{bare}\left(1\frac{1}{2N}\mathrm{ln}\frac{\mathrm{\Lambda }_\alpha }{\lambda }\right).$$ (19) This is the the first term of $`1/N`$ expansion, and keeping in mind scaling behavior we find $$\mathrm{\Gamma }_\lambda \lambda ^y,y\frac{1}{2N}=\frac{1}{6}.$$ (20) The impurity self energy is given by the diagram presented in Fig. 5. It differs from (15) by the vertexes: the bare vertexes are replaced by the “exact” ones given by eq. (20). Effectively this introduces an additional factor $`\omega ^{2y}`$ in the integrand in eq. (15). Assuming power behavior of the Green’s function $$G(ϵ)\frac{1}{(ϵ_0ϵ)^x}$$ (21) and performing integration we find $$\mathrm{\Sigma }(ϵ)(ϵ_0ϵ)^{2yx+1}+const.$$ (22) Substitution of this self energy into Dyson equation (14) gives the following condition of self consistency $$(ϵ_0ϵ)^x(ϵ_0ϵ)^{x2y1}.$$ (23) Therefore the Green’s function critical index with account of the leading $`1/N`$ correction is $$x=\frac{1}{2}+y=\frac{1}{2}(1+\frac{1}{N})0.67$$ (24) Away from the critical point the quasiparticle pole appears in the Greens function. It absorbs spectral weight of the incoherent Green’s function (21) from the area $`ϵ_0ϵ\mathrm{\Delta }_\alpha `$. Therefore the quasiparticle residue scales as $$Z\mathrm{\Delta }_\alpha ^z,z=1x=\frac{1}{2}y=\frac{1}{2}(1\frac{1}{N})0.33.$$ (25) ## IV Magnetic moments of the impurity, Susceptibilities Following Vojta, Buragohain, and Sachdev we consider two different types of the magnetic interaction. The first one is an interaction when the magnetic field $`h`$ interacts only with the impurity $$\stackrel{~}{H}_M^{(I)}=2\mathrm{𝐬𝐡}.$$ (26) The second case is homogeneous magnetic field $$H_M^{(II)}=2\mathrm{𝐬𝐡}\underset{i,n}{}2𝐒_𝐢^{(𝐧)}𝐡.$$ (27) It is clear that in the first case the renormalized magnetic moment is proportional to the quasiparticle residue $`Z`$, and hence, according to (25) it scales as $$\mu ^{(I)}\mathrm{\Delta }^z.$$ (28) If we consider the system exactly at the critical point, but at finite temperature, then the effective spin-wave gap is equal to the temperature, see e.g. Refs. : $$\mathrm{\Delta }\mathrm{\Delta }_T0.962T.$$ (29) Together with (28) this gives the following dependence of the impurity magnetic susceptibility on temperature. $$\chi _{imp}^{(I)}=\frac{\mu ^2}{T}\frac{1}{T^{2y}}\frac{1}{T^{1/N}}=\frac{1}{T^{0.33}}.$$ (30) For homogeneous magnetic field the impurity magnetic moment is not renormalized because the interaction (27) is proportional to the total spin which is conserved. The magnetic moment is certainly redistributed over the volume of size $`r1/\mathrm{\Delta }_T`$, but the value is conserved. For this reason the impurity susceptibility is given by usual the Curie law. $$\chi _{imp}^{(II)}=\frac{1}{T}.$$ (31) This conclusion is not quite trivial since the magnon cloud size $`r\mathrm{}`$ at $`T0`$. Besides eq. (31) does not agree with a conclusion from the paper . Therefore in the next section we also present a diagrammatic prove of our statement. It is rather technical section and a reader who is satisfied by the general arguments can skip it. ## V Corrections to the impurity magnetic moment In the present section only homogeneous magnetic field (27) is considered, therefore to simplify notations we omit the superscript (II). We restrict our consideration by single loop corrections, and follow the way used in Ref. for calculation of the spin-wave magnetic moment. First we discuss a zero temperature case. Single loop self energy is given by eq. (11). This as a simple perturbation theory in $`\alpha `$, but one can also consider this as a contribution to renormalization group equations with running coupling constant $`\alpha `$ and running infrared cutoff $`\mathrm{\Delta }`$. Anyway, correction to the quasiparticle residue due to (11) is following $$\delta Z=\frac{\mathrm{\Sigma }^{(1)}}{\mathrm{\Delta }}=\frac{\alpha ^2}{\mathrm{\Delta }}.$$ (32) The impurity magnetic moment is renormalized according to the diagrams shown in Fig. 6. First contribution comes from the correction to the quasiparticle residue: $`\delta \mu _{6a}=\delta Z`$. Straightforward calculation gives $`\delta \mu _{6b}=\alpha ^2/(3\mathrm{\Delta })`$ and $`\delta \mu _{6c}=4\alpha ^2/(3\mathrm{\Delta })`$. Altogether this gives $$\delta \mu _6=\delta \mu _{6a}+\delta \mu _{6b}+\delta \mu _{6c}=0.$$ (33) So as one shall expect there is no renormalization of the magnetic moment. At finite temperature the relation (33) remains valid since the only thing we have to do is to replace $`\mathrm{\Delta }\mathrm{\Delta }_T`$. However at finite temperature there are also additional diagrams which are due to the heat bath of the excited magnons . First of all these are the two contributions to the self energy shown in Fig. 7 Both contributions are proportional to the magnon mean occupation number $$n_𝐪=\frac{1}{e^{\omega _𝐪/T}1},$$ (34) but they are of the opposite sign and exactly cancel each other. So they do not influence the position of the quasiparticle pole. However these diagrams contribute equally to the quasiparticle residue $$\delta Z_T=2\alpha ^2\frac{n_𝐪d\omega _𝐪}{\omega _𝐪^2}=\frac{2\alpha ^2}{\mathrm{\Delta }_T}_1^{\mathrm{}}\frac{dx}{x^2(e^x1)}.$$ (35) This gives a correction to the inpurity magnetic moment. Other thermally induced corrections to the magnetic moment are given by diagrams presented in Fig. 8. Straightforward calculation gives $`\delta \mu _{8a}=\delta \mu _{8b}=\frac{\alpha ^2}{3\mathrm{\Delta }_T}_1^{\mathrm{}}𝑑x/[x^2(e^x1)]`$ , and $`\delta \mu _{8c}=\delta \mu _{8d}=\frac{4\alpha ^2}{3\mathrm{\Delta }_T}_1^{\mathrm{}}𝑑x/[x^2(e^x1)]`$. Total thermally induced correction to the impurity magnetic moment is equal to $$\delta \mu _T=\delta Z_T+\delta \mu _{8a}+\delta \mu _{8b}+\delta \mu _{8c}+\delta \mu _{8d}=0.$$ (36) Together with eq. (33) this proves that the impurity magnetic moment is not renormalized. Thus, in spite of the cloud of virtual and thermal magnons, effectively the impurity in the external magnetic field can be described as a spin 1/2 system with unrenormalized magnetic moment. Hence we immediately come to eq. (31). We would like to stress that the arguments related to the magnetic moment guarantee only singular in $`T`$ part of the susceptibility, therefore instead of (31) it is more correct to write $`\chi _{imp}^{(II)}=1/T+const`$. ## VI Specific heat related to the impurity Binding energy of the impurity has been calculated in section III. Exactly at the critical point, at zero temperature, and at small coupling constant, $`\alpha ^21`$, the binding energy is $`ϵ_02.5\alpha ^2`$. A finite spin-wave gap $`\mathrm{\Delta }`$ pushes the position of the quasiparticle pole up, see Fig. 3. For a small gap the dependence of $`ϵ_0`$ on the gap is linear, and the coefficient of the proportionality is approximately 2: $$ϵ_0\alpha ^2(2.5+2\mathrm{\Delta }).$$ (37) The origin of $`\mathrm{\Delta }`$ in this equation is not important: whether $`\mathrm{\Delta }`$ is nonzero because the system is away from the critical point, or the system is at the critical point, but $`\mathrm{\Delta }`$ is nonzero because of temperature. In the later case we can use eq. (29) for the gap and hence the impurity specific heat is $$C_{imp}=\frac{dϵ_0}{dT}2\alpha ^2.$$ (38) Note that this is highly unusual result because $`C_{imp}0`$ at $`T0`$. For comparison: the bulk specific heat of a 2D antiferromagnet at the quantum critical point is quadratic in temperature, $`C_{bulk}T^2`$, see e.g. Refs. . From (38) one concludes that the impurity entropy is $`S_{imp}=C_{imp}𝑑T/T\mathrm{ln}T`$. Strictly speaking this is nonsense because it gives $`S_{imp}(T=0)=\mathrm{}`$. This probably indicates that there is a small nonzero critical index $`\xi `$ in the specific heat dependence: $`C_{imp}T^\xi `$. Equation (37) has been derived in the noncrossing approximation. One can check that single loop vertex corrections considered at the end of the section III do not change this equation. It probably means that nonzero $`\xi `$ can appear only due to the two loop corrections, $`\xi 1/N^2`$. There is no doubt that this is a very interesting problem which deserves further analytical analysis and which can be also studied in numerical simulations. ## VII Static interaction between two distant impurities Interaction between two impurities is given by diagram shown in Fig. 9. Since the impurities are localized, we have to integrate over all possible momenta transfer q. Away from the critical point there are two distinct regimes: $`c/r<\mathrm{\Delta }`$ and $`c/r>\mathrm{\Delta }`$, where $`r`$ is distance between the impurities and $`c1.9J`$ is the spin-wave velocity. In the first regime the interaction drops down exponentially with distance and it is not an interesting case. We consider only $`c/r>\mathrm{\Delta }`$ case which is also relevant to the critical point. We also assume that the coupling constant is strong enough: $`\alpha ^2/\mathrm{\Delta }>1`$. The impurity-magnon interaction is given by the Hamiltonian (9). According to the consideration in section III it is renormalized as $`H_{imp}Z_q\mathrm{\Gamma }(q)`$, where $`Z_q`$ is quasiparticle residue (17), and $`\mathrm{\Gamma }(q)`$ is the vertex (20). Therefore the interaction corresponding to the diagram Fig.9 is of the form $$V(r)=\frac{Z_q^2\mathrm{\Gamma }_q^2e^{i\mathrm{𝐪𝐫}}}{\omega _𝐪}\frac{d^2q}{(2\pi )^2}=\stackrel{}{\sigma }_1\stackrel{}{\sigma }_1\frac{Aj^2}{4}\frac{\left[Z_q\mathrm{\Gamma }_q/\mathrm{\Gamma }_{bare}\right]^2e^{i\mathrm{𝐪𝐫}}}{\omega _𝐪^2}\frac{d^2q}{(2\pi )^2},$$ (39) where $`\stackrel{}{\sigma }_1`$ and $`\stackrel{}{\sigma }_2`$ are the impurities Pauli matrixes. One power of $`\omega _𝐪`$ in the denominator appears because of the spin-wave propagator and another power is due to the bare interaction (9). With account of (12), (17), and (20) this gives $$V(r)\stackrel{}{\sigma }_1\stackrel{}{\sigma }_1\alpha ^{24z4y}q^{2(z+y1)}e^{i\mathrm{𝐪𝐫}}d^2q\frac{\stackrel{}{\sigma }_1\stackrel{}{\sigma }_1}{r}.$$ (40) This is a long range ferromagnetic interaction. An interesting fact is that the interaction is independent of the bare coupling constant $`\alpha `$. Another interesting fact is that there is no renormalization of the power in the $`1/r`$ dependence. However we stress that these facts have been proven only in the one loop approximation (one loop above the leading noncrossing approximation). ## VIII conclusion The dynamics of the magnetic impurity in a 2D antiferromagnet close to the critical point is highly nontrivial. To some extend it is similar to the dynamics in the Kondo problem. In the present paper we have considered spin 1/2 impurity. We have calculated indexes for critical behavior of the Green’s function, vertexes, magnetic moments and magnetic susceptibilities. We have also considered the impurity specific heat and long range interaction between two impurities which is due to the spin-wave exchange. ###### Acknowledgements. I am grateful to A.V. Chubukov, K. Le Hur, M. Yu. Kuchiev, M. Troyer and especially to A. Sandvik for very helpful discussions. The main part of this work has been done during my stay at ITP UCSB and ITP ETH Zurich. It was supported by NSF Grant PHY94-07194 and by the Swiss National Fund.
warning/0004/hep-ph0004209.html	ar5iv	text	# References Parapositronium (p-Ps) and orthopositronium (o-Ps) are bound states of $`e^+e^{}`$ whose lifetimes for 2$`\gamma `$ and 3$`\gamma `$ channels, respectively, have been measured with high precision: $`\mathrm{\Gamma }(\text{p-Ps}\gamma \gamma )`$ $`=`$ $`7090.9(1.7)\mu s^1\text{[1]}`$ (1) $`\mathrm{\Gamma }(\text{o-Ps}\gamma \gamma \gamma )`$ $`=`$ $`\{\begin{array}{c}7.0514(14)\mu s^1\text{[2]},\\ 7.0482(16)\mu s^1\text{[3]},\\ 7.0398(29)\mu s^1\text{[4]}.\end{array}`$ (5) The corresponding theoretical predictions which include perturbative QED corrections to a non-relativistic treatment of the bound state wave function, have been computed also with high accuracy : $`\mathrm{\Gamma }(\text{p-Ps}\gamma \gamma )`$ $`=`$ $`{\displaystyle \frac{\alpha ^5m}{2}}\left(1(5{\displaystyle \frac{\pi ^2}{4}}){\displaystyle \frac{\alpha }{\pi }}+2\alpha ^2\mathrm{ln}{\displaystyle \frac{1}{\alpha }}+1.75(30)\left({\displaystyle \frac{\alpha }{\pi }}\right)^2{\displaystyle \frac{3\alpha ^3}{2\pi }}\mathrm{ln}^2{\displaystyle \frac{1}{\alpha }}\right)`$ (6) $`=`$ $`7989.50(2)\mu s^1`$ $`\mathrm{\Gamma }(\text{o-Ps}\gamma \gamma \gamma )`$ $`=`$ $`{\displaystyle \frac{2(\pi ^29)\alpha ^6m}{9\pi }}\left(110.28661(1){\displaystyle \frac{\alpha }{\pi }}{\displaystyle \frac{\alpha ^2}{3}}\mathrm{ln}{\displaystyle \frac{1}{\alpha }}+B_0\left({\displaystyle \frac{\alpha }{\pi }}\right)^2{\displaystyle \frac{3\alpha ^3}{2\pi }}\mathrm{ln}^2{\displaystyle \frac{1}{\alpha }}\right)`$ (7) $`=`$ $`(7.0382+B_0\mathrm{\hspace{0.25em}0.39}\times 10^4)\mu s^1,`$ where $`m`$ denotes de mass of the electron and $`\alpha `$ is the fine structure constant. As it can be observed, predictions for parapositronium are in very good agreement with experiment, while the experimental results for orthopositronium reported in Refs. largely disagree with the theoretical expectation. Recent theoretical efforts have focused on a more complete evaluation of the non-logarithmic $`O(\alpha ^2)`$ perturbative corrections for orthopositronium, with the result $`B_0=44.52(26)`$ or, equivalently, $$\mathrm{\Gamma }(\text{o-Ps}\gamma \gamma \gamma )=7.039934(10)\mu s^1.$$ (8) This result renders even closer the theoretical prediction to the experimental measurement of Ref. . Although these achievements of perturbative QED for bound states are impressive, little is known about the effects of non-perturbative corrections. In the present paper we implement a mechanism, imposed by analyticity and gauge invariance, which allows to introduce binding energy (BE) effects in the positronium decay amplitudes. These new corrections affect the $`\gamma \gamma `$ decay rate of parapositronium at the order $`\alpha ^4`$, while corrections of order $`\alpha ^2`$ are induced in the orthopositronium decay rate and in the single photon spectrum of p-Dm$`e^+e^{}\gamma `$ and o-Ps$`\gamma \gamma \gamma `$ decays (p-Dm (o-Dm) denotes the J=0(1) bound state of the dimuonium system $`\mu ^+\mu ^{}`$). Moreover, the single photon energy spectrum in the last two decays previously mentioned is softened with respect to results obtained when BE effects are neglected. Our implementation of these non-perturbative corrections relies on well grounded basis of QED such as the requirement of a correct behavior of the decay amplitude in terms of the photon energies (analyticity) and on electromagnetic gauge invariance. We briefly discuss at the end, the possible implications for analogous decays of quarkonium states. Details of our calculations will be given elsewhere . Let us first try to explain how problems related to analyticity of the positronium decay amplitudes appears in the usual approaches. Current calculations of the positronium rates assume a factorization approximation of the dynamics contained in the decay amplitude. The amplitudes for positronium decays into a given final state $`X`$ is approximated as the product of the amplitude for the bound state annihilation (described by $`\mathrm{\Phi }_0`$, the $`e^+e^{}`$ wavefunction at the origin), times the annihilation amplitude of free leptons into the final state $`X`$. Under this assumption, the positronium decay rates can be written as : $$\mathrm{\Gamma }(\text{Ps}(^{2J+1}S)X)=\frac{1}{2J+1}\|\mathrm{\Phi }_0\|^2(v_{rel}\sigma (e^+e^{}X))_{v_{rel}0},$$ (9) where $`\sigma (e^+e^{}X)`$ denotes the total cross section for $`e^+e^{}X`$ and $`v_{rel}`$ is the relative velocity of the constituents in their center of mass frame. Since the leptons in the intermediate stage are taken on their mass-shell, the annihilation amplitude for $`e^+e^{}n\gamma `$ is automatically gauge-invariant. However, in this approximation the positronium decay amplitudes do not exhibit the expected analytical behavior. Indeed, for definiteness let us consider the 3$`\gamma `$ decay mode of orthopositronium where the photon energy is kinematically allowed to vanish. The behavior of the $`e^+e^{}3\gamma `$ amplitude in terms of the photon momenta is driven in the soft-photon limit by the intermediate electron propagators as follows: $$\frac{i}{\mathit{}\mathit{}_im}=\frac{i(\mathit{}\mathit{}_i+m)}{2l_ik},$$ (10) where $`k`$ is the four-momentum of the electron that emits a photon of four-momentum $`l_i`$. Note that in the soft-photon limit $`l_i0`$, the positronium decay amplitude seems to diverge as $`l_i^1`$ . Actually, this is only apparent since selection rules cancels these infrared divergencies in the static limit ($`k=(m,0,0,0))`$, and the amplitude starts indeed at order $`l_i^0`$. This is in contradiction with the fact that in the soft-photon limit the amplitude must vanish since o-Ps $`\gamma \gamma \gamma `$ involve only neutral external bosons. We can try to cure this bad analytical behavior by realizing that electrons in the intermediate stage are always off their mass-shell due to BE effects. Actually, the $`e^+e^{}`$ bound states involve two mass scales: the mass $`M`$ of the bound state and the mass $`m`$ of the constituent electrons. These masses differ by terms of order $`\alpha ^2`$ and are related through the binding energy $`E_{bind}`$, $$E_{bind}M2m=\frac{1}{4}m\alpha ^2.$$ (11) Thus, the relevant momentum scale for intermediate electrons is determined by $`M/2`$. In this case, the lepton propagator involved in the $`e^+e^{}3\gamma `$ decay amplitude becomes ($`k^2=M^2/4`$): $$\frac{i}{\mathit{}\mathit{}_im}=\frac{i(\mathit{}\mathit{}_i+m)}{2l_ik\gamma ^2},$$ (12) where $`\gamma ^2m^2M^2/4M^2\alpha ^2/16`$, takes into account the bound state nature of the $`e^+e^{}`$ pair in the rescattering process. The presence of $`\gamma ^2`$ in the denominator would provide to the decay amplitude a better analytical behavior. Unfortunately, the gauge invariance of the $`e^+e^{}\gamma \gamma \gamma `$ amplitude is spoiled due to the presence of the two mass scales $`M`$ and $`m`$. Thus, a more sophisticated procedure is required to restore analyticity without destroying gauge invariance. A natural way to incorporate BE effects is by considering a model with a loop of virtual leptons for the decays of positronium at lowest order. For definiteness we consider the $`n\gamma `$ decay of the positronium state $`B(^{2J+1}S)`$: $`J=0(1)`$ being the spin of the p-Ps(p-Dm) state, $`n=2`$ for p-Ps or p-Dm, and $`n=3`$ for o-Ps or o-Dm decays (the $`e^+e^{}\gamma `$ mode of paradimuonium can be reached from p-Dm$`\gamma \gamma ^{}`$ where $`\gamma ^{}`$ is a virtual photon). In this model we need to introduce a coupling $`F_B\mathrm{\Gamma }`$ to describe the $`Be^+e^{}`$ vertex, where $`\mathrm{\Gamma }=\gamma _5(\mathit{\eta ̸})`$ for para(ortho)-positronium decay, $`\eta _\alpha `$ is the polarization four-vector of the $`B(^3S)`$ state and $`F_B`$ is a form factor that describes the structure of the vertex. The evaluation of the $`B(^{2J+1}S)X`$ decay amplitude follows standard rules (see Fig. 1). Applying the Feynman rules to Fig. 1 the decay amplitude is given by (contraction with photon polarization vectors $`ϵ_\mu (l_1)\mathrm{}`$ must be understood): $`^{\mu \nu \mathrm{}}(Bn\gamma )`$ $`=`$ $`{\displaystyle \frac{d^4q}{(2\pi )^4}F_B\mathrm{Tr}\left\{\mathrm{\Gamma }\frac{i(\mathit{}\frac{\mathit{}}{2}+m)}{(q\frac{P}{2})^2m^2}\mathrm{\Gamma }^{\mu \nu \mathrm{}}\frac{i(\mathit{}+\frac{\mathit{}}{2}+m)}{(q+\frac{P}{2})^2m^2}\right\}},`$ (13) where $`\mathrm{\Gamma }^{\mu \nu \mathrm{}}`$ is a properly symmetrized amplitude for annihilation of virtual lepton pairs into $`n\gamma `$. It is important to emphasize that this amplitude is gauge-invariant if $`F_B`$ is constant. To clearly illustrate how this loop model reproduces, in the limit of zero BE, the lowest order amplitudes of the on-shell approximation, let us consider the following ansatz for $`F_B`$: $$F_B=iC\mathrm{\Phi }_0\frac{8\pi \gamma }{(𝐪^2+\gamma ^2)^2}(𝐪^2+\gamma ^2)$$ (14) where $`\gamma `$ contains the BE, $`\mathrm{\Phi }_0=\sqrt{\alpha ^3m^3/8\pi }`$ is the ground state wavefunction of $`e^+e^{}`$ at the origin, and $`C=2/\sqrt{M}`$ is a normalization constant. Using well known representations of Dirac-delta functions we can check that in the limit of zero BE ($`\gamma 0,mM/20`$) we obtain: $$\frac{F_B}{(q\frac{P}{2})^2m^2}\frac{1}{(q+\frac{P}{2})^2m^2}\frac{1}{2}C\mathrm{\Phi }_0(2\pi )^4\delta ^{(4)}(q)\frac{\sqrt{𝐪^2+m^2}+\frac{M}{2}}{q_0^2(\sqrt{𝐪^2+m^2}+\frac{M}{2})^2}.$$ (15) Thus, upon (trivial) integration of Eq. (10) one gets: $$^{\mu \nu \mathrm{}}(Bn\gamma )=\frac{C\mathrm{\Phi }_0}{8M}\mathrm{Tr}\left\{\mathrm{\Gamma }\left(\mathit{}M\right)\overline{\mathrm{\Gamma }}^{\mu \nu \mathrm{}}\left(\mathit{}+M\right)\right\},$$ (16) where $`\overline{\mathrm{\Gamma }}^{\mu \nu \mathrm{}}`$ denotes the reduced vertex evaluated at $`q=0`$. Note that the factors $`(\mathit{}\pm M)`$ play the role of projector operators external to the action of photon vertices. Setting in the rest frame of positronium and using $`2m=M`$ (and the condition $`P_\alpha \eta ^\alpha =0`$ for orthopositronium) we arrive at the well known results of the factorization approximation, namely: $$^{\mu \nu \mathrm{}}(Bn\gamma )=C\mathrm{\Phi }_0\frac{M}{2\sqrt{2}}\mathrm{Tr}\left\{\frac{1+\gamma _0}{\sqrt{2}}\mathrm{\Gamma }\overline{\mathrm{\Gamma }}^{\mu \nu \mathrm{}}\right\}.$$ (17) Let us return to the issues concerning gauge-invariance in the context of the present model when BE is not neglected. Gauge invariance requires that the amplitude in Eq. (10) (contracted with photon polarizations) vanishes when $`ϵ_\alpha (l_i)l_{i\alpha }`$ for any external photon. The amplitude for p-Ps$`\gamma \gamma `$ is always gauge-invariant, no matter the specific form of $`F_B`$. This follows from the fact that Lorentz covariance for the axial-$`\gamma \gamma `$ vertex implies that the amplitude should corresponds to an effective operator $`F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }`$, which is automatically gauge-invariant. This is not the case for o-Ps$`\gamma \gamma \gamma `$ decays. In this case the amplitude of Eq. (10) is gauge invariant provided $`F_B`$ remains the same under the shifts $`qq+l_i`$ of the integration variable for all photons. Since this is achieved only if $`F_B`$ is constant, it means that other contributions should be added to the orthopositronium decay amplitude in this loop model in order to compensate for this lack of gauge invariance. In view of these difficulties, we propose and ansatz for positronium decay amplitudes that fulfills gauge invariance and analyticity simultaneously. Our recipe contains three steps: $`(a)`$ evaluate the expression $`\overline{\mathrm{\Gamma }}^{\mu \nu \mathrm{}}`$ in Eq. (13) for constituents masses $`M/2`$ (this modified expression will be denoted by $`\overline{\mathrm{\Gamma }}_{M/2}^{\mu \nu \mathrm{}}`$); this will ensure gauge invariance, $`(b)`$ introduce BE effects by multiplying the amplitude $`\overline{}(B(^{2J+1}S)n\gamma )`$ (i.e. the amplitude of Eq. (13) where we replace $`m`$ by $`M/2`$ in the argument of the Trace operator) by a factor: $`𝒜_1`$ for p-Ps (p-Dm) states and $`_{i=1}^3𝒜_i)`$ for o-Ps (o-Dm), where we have defined $`𝒜_iP.l_i/(P.l_i+\gamma ^2)`$ (note that $`𝒜_1=𝒜_2`$ for p-Ps$`2\gamma `$). This will warrant the correct analytical properties of the amplitude in the soft-photon limit and, $`(c)`$ obtain the decay rate by integration over the physical phase-space determined by the mass $`M`$. Let us first note that the proposed ansatz is fulfilled automatically for the p-Ps$`\gamma \gamma `$ decay. The reduced vertex in Eq. (13) is given by (remember $`q=0`$ and $`P^2=M^2`$): $`\overline{\mathrm{\Gamma }}^{\mu \nu }`$ $`=`$ $`(ie)^2{\displaystyle \frac{\gamma ^\mu \left(\frac{\mathit{}}{2}+\mathit{}_1+m\right)\gamma ^\nu +\gamma ^\nu \left(\frac{\mathit{}}{2}\mathit{}_1+m\right)\gamma ^\mu }{\left(\frac{P}{2}l_1\right)^2m^2}}`$ (18) $`=`$ $`(ie)^2{\displaystyle \frac{\gamma ^\mu \left(\frac{\mathit{}}{2}+\mathit{}_1+m\right)\gamma ^\nu +\gamma ^\nu \left(\frac{\mathit{}}{2}\mathit{}_1+m\right)\gamma ^\mu }{\left(\frac{P}{2}l_1\right)^2\frac{M^2}{4}}}\left(𝒜_1\right).`$ In the second line of Eq. (15) we can replace $`mM/2`$ in the numerator because terms proportional to $`m`$ in $`\overline{\mathrm{\Gamma }}^{\mu \nu }`$ cancel when performing the trace in Eq. (13). Therefore, gauge invariance is preserved independently of $`m`$ and we have: $$^{\mu \nu }(\text{p-Ps}\gamma \gamma )=\overline{}^{\mu \nu }(\text{p-Ps}\gamma \gamma )=\frac{C\mathrm{\Phi }_0}{8M}\mathrm{Tr}\left\{\gamma _5(\mathit{}M)\overline{\mathrm{\Gamma }}_{M/2}^{\mu \nu }(\mathit{}+M)\right\}𝒜_1.$$ (19) This amplitude satisfies analyticity, as required. Secondly, the reduced vertex for orthopositronium decay $`\overline{\mathrm{\Gamma }}^{\mu \nu \rho }`$ can be worked in the following way. In order to accomplish gauge invariance, we are forced to replace $`m`$ by $`M/2`$ for the mass of the constituents and simultaneously add a third analytical factor $`A_i`$ to each of the six amplitudes contributing to $`\overline{\mathrm{\Gamma }}^{\mu \nu \rho }`$. This gives rise to the final gauge-invariant and analytical amplitude for orthopositronium decay: $$\overline{}^{\mu \nu \rho }=\frac{C\mathrm{\Phi }_0}{8M}\mathrm{Tr}\left\{\mathit{\eta ̸}(\mathit{}M)\overline{\mathrm{\Gamma }}_{M/2}^{\mu \nu \rho }(\mathit{}+M)\right\}\underset{i=1}{\overset{3}{}}𝒜_i,$$ (20) where $`\eta _\alpha `$ represents the polarization four-vector of orthopositronium. In the soft photon limit ($`l_i0`$), $`\overline{}^{\mu \nu \rho }`$ vanishes as required. Up to now we have considered the effects of nonzero BE corrections in the dynamics of positronium decays as expressed in the decay amplitude. It is clear that the physical phase space for these decays is determined by the masses of external particles, in particular the initial available energy $`M`$. Thus, we will study the effects of these non-perturbative corrections in observables associated to positronium decays, which contain the BE effects in the dynamics (amplitude) and the kinematics (phase-space). 2$`\gamma `$ decay of parapositronium. The decay amplitude for p-Ps$`\gamma \gamma `$ obtained from Eq. (16) can be expressed as: $$\overline{}^{\mu \nu }(B(^1S)\gamma \gamma )=2Ce^2\mathrm{\Phi }_0\frac{\epsilon ^{\mu \nu \alpha \beta }l_{1\alpha }P_\beta }{Pl_1}𝒜_1.$$ (21) The corresponding rate for this two-body decay where photons fly apart with energy $`M/2`$ in the rest frame of p-Ps, is given by ($`2P.l_1=M^2`$): $`\mathrm{\Gamma }(\text{p-Ps}\gamma \gamma )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\alpha ^5m\times {\displaystyle \frac{4m^2}{M^2}}\left[{\displaystyle \frac{1}{1+2\gamma ^2/M^2}}\right]^2`$ (22) $``$ $`{\displaystyle \frac{1}{2}}\alpha ^5m\left(1{\displaystyle \frac{\alpha ^4}{64}}\right),`$ where we have used the definition of $`\gamma ^2`$ as given after Eq. (9) and we have neglected corrections of $`O(\alpha ^6)`$. Thus, (non-perturbative) BE corrections for $`2\gamma `$ decays of parapositronium are negligible small and are beyond present experimental precision. para-dimuonium (p-Dm) decay into $`e^+e^{}\gamma `$ The decay mechanism for this process is similar to the previous one, but where one virtual photon converts into a $`e^+e^{}`$ pair. Therefore the corresponding amplitude is: $$\overline{}^\mu (\text{p-Dm}e^+e^{}\gamma )=2Ce^3\mathrm{\Phi }_0\frac{\epsilon ^{\mu \nu \alpha \beta }l_{1\alpha }P_\beta }{Pl_1}𝒜_1\frac{\left\{\overline{u}(p)\gamma _\nu v(p^{})\right\}}{r^2},$$ (23) where $`l_1`$ and $`rp+p^{}=Pl_1`$ denote, respectively, the four-momenta of the photon and the $`e^+e^{}`$ pair. The single photon spectrum in this case is given by ($`a4m_e^2/M^2,x=2E_\gamma /M`$): $$\frac{d\mathrm{\Gamma }(\text{p-Dm}e^+e^{}\gamma )}{dx}=\frac{16\alpha ^3\mathrm{\Phi }_0^2}{3M^2}\sqrt{1\frac{a}{1x}}\left[a+2(1x)\right]\frac{x^3}{(1x)^2}\left(x+\frac{2\gamma ^2}{M^2}\right)^2.$$ (24) Observe that this spectrum falls as $`x^3`$ when $`x0`$ due to non-zero BE effects, instead of the usual behavior (proportional to $`x`$) expected when these effects are neglected. It is interesting to note that when $`x2\gamma ^2/M^2`$, Eq. (21) behaves as the corresponding photon energy spectrum in $`\pi ^0e^+e^{}\gamma `$ decay, for a point-like pion vertex. Thus, BE corrections affect the shape of the spectrum or, conversely, actually probes the structure of the bound state. Indeed, Eq. (21) can be seen as a result that extrapolates the spectrum between the $`e^+e^{}\gamma `$ decay of a pseudoscalar point particle ($`\pi ^0`$ case) and the standard bound state calculations (paradimuonium decay without BE corrections). A closed (but long) analytic expression for the decay rate can be obtained from integration of Eq. (21). A useful approximation that takes into account leading non-vanishing corrections of $`O(\alpha ^2)`$ is: $$\mathrm{\Gamma }(\text{p-Dm}e^+e^{}\gamma )=\frac{\alpha ^6m}{6\pi }\left[F_0(1a)^{3/2}\frac{\alpha ^2}{2}+O(\alpha ^4)\right],$$ (25) where the function $`F_0(4/3)\sqrt{1a}(a4)+2\mathrm{ln}[(1+\sqrt{1a})/(1\sqrt{1a})]`$ fixes the lowest order rate. 3$`\gamma `$ decays of orthopositronium The squared unpolarized amplitude for o-Ps$`\gamma \gamma \gamma `$ obtained from Eq. (17), in terms of dimensionless photon energy variables $`x_i=2E_{\gamma i}/M`$ ($`x_1+x_2+x_3=2`$) is given by: $$\underset{pols}{}\|(3\gamma )\|^2\left[\left(\frac{1x_1}{x_2x_3}\right)^2+\left(\frac{1x_2}{x_1x_3}\right)^2+\left(\frac{1x_3}{x_1x_2}\right)^2\right]\underset{i=1}{\overset{3}{}}\left(\frac{x_i}{x_i+\frac{2\gamma ^2}{M^2}}\right)^2.$$ (26) As in the previous case, the single photon spectrum will be softened by BE corrections due to the last factor in the squared amplitude. The phase space integration of Eq. (23) can be performed in analytic form and expressed in terms of dilogarithmic functions. However, it is more illustrative to express the decay rate in terms of an expansion in powers of $`\alpha ^2`$ (or $`\gamma ^2`$). If we express the decay rate in terms of the mass $`m`$ of constituents, we obtain: $$\mathrm{\Gamma }(\text{o-Ps}3\gamma )\alpha ^6m\frac{2(\pi ^29)}{9\pi }\left[1\frac{5}{4}\alpha ^2\right].$$ (27) Thus, BE corrections affects the decay rate of orthopositronium at order $`\alpha ^2`$ which are indeed relevant when confronted to accuracy of present experiments. A comparison of Eqs. (4) and (24) indicates that the net effects of BE corrections is to resize the coefficient appearing in nonlogarithmic $`O(\alpha ^2)`$ corrections of orthopositronium, namely: $$B_0B_0\frac{5\pi ^2}{4}44.52(26)12.34.$$ (28) i.e. a non-negligible 28% correction at order $`\alpha ^2`$. Before concluding, let us address some comments on possible implications of BE corrections in analogous decays of quarkonia. Notice that BE in quarkonia is not simply related to the masses of the constituent quarks and to the coupling $`\alpha _s`$ as in Eq. (8) due to confinement. If we assume, however, that BE corrections do affect quarkonium decays in a similar form as in positronium, we can address some apparent conflicts in some of their inclusive hadronic and radiative decays. First, the photon spectrum measured in $`J/\psi gg\gamma `$ decays seems to be softer than predicted by perturbative QCD which indicate possible large non-perturbative effects. This is precisely the effect induced by BE corrections in the single photon spectrum of o-Ps$`\gamma \gamma \gamma `$. Second, different inclusive rates in quarkonium decays as discussed in will, in general, get decreased by BE corrections. This may increase the relatively low values of the strong coupling constants extracted from ratios of experimental quarkonium rates . Finally, these effects will also manifest in the “14 % rule” observed in the ratio $`BR(\psi (2S)X)/BR(J/\psi X)`$ for single photon mediated decays , because BE are different for these two radial excitations of charmonium. In this paper we have computed binding energy corrections to positronium decays in a specific ansatz where these non-perturbative effects are introduced as a necessity to account for analyticity and gauge invariance of the corresponding QED amplitudes. The BE corrections (of order $`\alpha ^2`$) to orthopositronium decay are indeed relevant in view of recent efforts to achieve a precise comparison of theory and experiment. When extended to the quarkonium sector, these BE corrections may contribute to solve apparent discrepancies observed between experimental data and perturbative QCD calculations of some inclusive rates of quarkonia. Acknowledgements: C.S. acknowledges financial support from FNRS (Belgium). G.L.C. was partially supported by Conacyt (México) under contract No. 32429.
warning/0004/hep-ex0004027.html	ar5iv	text	# Measurement of 𝐽/𝜓 and 𝜓⁢(2⁢𝑆) polarization in 𝑝⁢𝑝̄ collisions at √𝑠=1.8 TeV ## Abstract We have measured the polarization of $`J/\psi `$ and $`\psi (2S)`$ mesons produced in $`p\overline{p}`$ collisions at $`\sqrt{s}=1.8`$ TeV, using data collected at CDF during 1992-95. The polarization of promptly produced $`J/\psi `$ \[$`\psi (2S)`$\] mesons is isolated from those produced in $`B`$-hadron decay, and measured over the kinematic range $`4[5.5]<P_T<20\mathrm{G}eV/c`$ and $`\|y\|<0.6`$. For PT > 12 > subscript𝑃𝑇12P_{T}\mathrel{\vbox{\kern 0.0pt\hbox{$>$} \kern 0.0pt\hbox{$\sim$} }}12 $`\mathrm{G}eV/c`$ we do not observe significant polarization in the prompt component. T. Affolder,<sup>21</sup> H. Akimoto,<sup>43</sup> A. Akopian,<sup>36</sup> M. G. Albrow,<sup>10</sup> P. Amaral,<sup>7</sup> S. R. Amendolia,<sup>32</sup> D. Amidei,<sup>24</sup> K. Anikeev,<sup>22</sup> J. Antos,<sup>1</sup> G. Apollinari,<sup>10</sup> T. Arisawa,<sup>43</sup> T. Asakawa,<sup>41</sup> W. Ashmanskas,<sup>7</sup> M. Atac,<sup>10</sup> F. Azfar,<sup>29</sup> P. Azzi-Bacchetta,<sup>30</sup> N. Bacchetta,<sup>30</sup> M. W. Bailey,<sup>26</sup> S. Bailey,<sup>14</sup> P. de Barbaro,<sup>35</sup> A. Barbaro-Galtieri,<sup>21</sup> V. E. Barnes,<sup>34</sup> B. A. Barnett,<sup>17</sup> M. Barone,<sup>12</sup> G. Bauer,<sup>22</sup> F. Bedeschi,<sup>32</sup> S. Belforte,<sup>40</sup> G. Bellettini,<sup>32</sup> J. Bellinger,<sup>44</sup> D. Benjamin,<sup>9</sup> J. Bensinger,<sup>4</sup> A. Beretvas,<sup>10</sup> J. P. Berge,<sup>10</sup> J. Berryhill,<sup>7</sup> B. Bevensee,<sup>31</sup> A. Bhatti,<sup>36</sup> M. Binkley,<sup>10</sup> D. Bisello,<sup>30</sup> R. E. Blair,<sup>2</sup> C. Blocker,<sup>4</sup> K. Bloom,<sup>24</sup> B. Blumenfeld,<sup>17</sup> S. R. Blusk,<sup>35</sup> A. Bocci,<sup>32</sup> A. Bodek,<sup>35</sup> W. Bokhari,<sup>31</sup> G. Bolla,<sup>34</sup> Y. Bonushkin,<sup>5</sup> D. 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Clark,<sup>13</sup> A. Connolly,<sup>21</sup> J. Conway,<sup>37</sup> J. Cooper,<sup>10</sup> M. Cordelli,<sup>12</sup> J. Cranshaw,<sup>39</sup> D. Cronin-Hennessy,<sup>9</sup> R. Cropp,<sup>23</sup> R. Culbertson,<sup>7</sup> D. Dagenhart,<sup>42</sup> F. DeJongh,<sup>10</sup> S. Dell’Agnello,<sup>12</sup> M. Dell’Orso,<sup>32</sup> R. Demina,<sup>10</sup> L. Demortier,<sup>36</sup> M. Deninno,<sup>3</sup> P. F. Derwent,<sup>10</sup> T. Devlin,<sup>37</sup> J. R. Dittmann,<sup>10</sup> S. Donati,<sup>32</sup> J. Done,<sup>38</sup> T. Dorigo,<sup>14</sup> N. Eddy,<sup>16</sup> K. Einsweiler,<sup>21</sup> J. E. Elias,<sup>10</sup> E. Engels, Jr.,<sup>33</sup> W. Erdmann,<sup>10</sup> D. Errede,<sup>16</sup> S. Errede,<sup>16</sup> Q. Fan,<sup>35</sup> R. G. Feild,<sup>45</sup> C. Ferretti,<sup>32</sup> R. D. Field,<sup>11</sup> I. Fiori,<sup>3</sup> B. Flaugher,<sup>10</sup> G. W. Foster,<sup>10</sup> M. Franklin,<sup>14</sup> J. Freeman,<sup>10</sup> J. Friedman,<sup>22</sup> Y. Fukui,<sup>20</sup> I. Furic,<sup>22</sup> S. Galeotti,<sup>32</sup> M. Gallinaro,<sup>36</sup> T. Gao,<sup>31</sup> M. Garcia-Sciveres,<sup>21</sup> A. F. Garfinkel,<sup>34</sup> P. Gatti,<sup>30</sup> C. Gay,<sup>45</sup> S. Geer,<sup>10</sup> D. W. Gerdes,<sup>24</sup> P. Giannetti,<sup>32</sup> P. Giromini,<sup>12</sup> V. Glagolev,<sup>8</sup> M. Gold,<sup>26</sup> J. Goldstein,<sup>10</sup> A. Gordon,<sup>14</sup> A. T. Goshaw,<sup>9</sup> Y. Gotra,<sup>33</sup> K. Goulianos,<sup>36</sup> C. Green,<sup>34</sup> L. Groer,<sup>37</sup> C. Grosso-Pilcher,<sup>7</sup> M. Guenther,<sup>34</sup> G. Guillian,<sup>24</sup> J. Guimaraes da Costa,<sup>14</sup> R. S. Guo,<sup>1</sup> R. M. Haas,<sup>11</sup> C. Haber,<sup>21</sup> E. Hafen,<sup>22</sup> S. R. Hahn,<sup>10</sup> C. Hall,<sup>14</sup> T. Handa,<sup>15</sup> R. Handler,<sup>44</sup> W. Hao,<sup>39</sup> F. Happacher,<sup>12</sup> K. Hara,<sup>41</sup> A. D. Hardman,<sup>34</sup> R. M. Harris,<sup>10</sup> F. 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Zucchelli<sup>3</sup> (CDF Collaboration) <sup>1</sup> Institute of Physics, Academia Sinica, Taipei, Taiwan 11529, Republic of China <sup>2</sup> Argonne National Laboratory, Argonne, Illinois 60439 <sup>3</sup> Istituto Nazionale di Fisica Nucleare, University of Bologna, I-40127 Bologna, Italy <sup>4</sup> Brandeis University, Waltham, Massachusetts 02254 <sup>5</sup> University of California at Los Angeles, Los Angeles, California 90024 <sup>6</sup> Instituto de Fisica de Cantabria, CSIC-University of Cantabria, 39005 Santander, Spain <sup>7</sup> Enrico Fermi Institute, University of Chicago, Chicago, Illinois 60637 <sup>8</sup> Joint Institute for Nuclear Research, RU-141980 Dubna, Russia <sup>9</sup> Duke University, Durham, North Carolina 27708 <sup>10</sup> Fermi National Accelerator Laboratory, Batavia, Illinois 60510 <sup>11</sup> University of Florida, Gainesville, Florida 32611 <sup>12</sup> Laboratori Nazionali di Frascati, Istituto Nazionale di Fisica Nucleare, I-00044 Frascati, Italy <sup>13</sup> University of Geneva, CH-1211 Geneva 4, Switzerland <sup>14</sup> Harvard University, Cambridge, Massachusetts 02138 <sup>15</sup> Hiroshima University, Higashi-Hiroshima 724, Japan <sup>16</sup> University of Illinois, Urbana, Illinois 61801 <sup>17</sup> The Johns Hopkins University, Baltimore, Maryland 21218 <sup>18</sup> Institut für Experimentelle Kernphysik, Universität Karlsruhe, 76128 Karlsruhe, Germany <sup>19</sup> Korean Hadron Collider Laboratory: Kyungpook National University, Taegu 702-701; Seoul National University, Seoul 151-742; and SungKyunKwan University, Suwon 440-746; Korea <sup>20</sup> High Energy Accelerator Research Organization (KEK), Tsukuba, Ibaraki 305, Japan <sup>21</sup> Ernest Orlando Lawrence Berkeley National Laboratory, Berkeley, California 94720 <sup>22</sup> Massachusetts Institute of Technology, Cambridge, Massachusetts 02139 <sup>23</sup> Institute of Particle Physics: McGill University, Montreal H3A 2T8; and University of Toronto, Toronto M5S 1A7; Canada <sup>24</sup> University of Michigan, Ann Arbor, Michigan 48109 <sup>25</sup> Michigan State University, East Lansing, Michigan 48824 <sup>26</sup> University of New Mexico, Albuquerque, New Mexico 87131 <sup>27</sup> The Ohio State University, Columbus, Ohio 43210 <sup>28</sup> Osaka City University, Osaka 588, Japan <sup>29</sup> University of Oxford, Oxford OX1 3RH, United Kingdom <sup>30</sup> Universita di Padova, Istituto Nazionale di Fisica Nucleare, Sezione di Padova, I-35131 Padova, Italy <sup>31</sup> University of Pennsylvania, Philadelphia, Pennsylvania 19104 <sup>32</sup> Istituto Nazionale di Fisica Nucleare, University and Scuola Normale Superiore of Pisa, I-56100 Pisa, Italy <sup>33</sup> University of Pittsburgh, Pittsburgh, Pennsylvania 15260 <sup>34</sup> Purdue University, West Lafayette, Indiana 47907 <sup>35</sup> University of Rochester, Rochester, New York 14627 <sup>36</sup> Rockefeller University, New York, New York 10021 <sup>37</sup> Rutgers University, Piscataway, New Jersey 08855 <sup>38</sup> Texas A&M University, College Station, Texas 77843 <sup>39</sup> Texas Tech University, Lubbock, Texas 79409 <sup>40</sup> Istituto Nazionale di Fisica Nucleare, University of Trieste/ Udine, Italy <sup>41</sup> University of Tsukuba, Tsukuba, Ibaraki 305, Japan <sup>42</sup> Tufts University, Medford, Massachusetts 02155 <sup>43</sup> Waseda University, Tokyo 169, Japan <sup>44</sup> University of Wisconsin, Madison, Wisconsin 53706 <sup>45</sup> Yale University, New Haven, Connecticut 06520 The production of heavy quarkonia states, $`c\overline{c}`$ and $`b\overline{b}`$, provides a useful system for the study of Quantum Chromodynamics (QCD), as it involves both perturbative and nonperturbative energy scales. In $`p\overline{p}`$ collisions, charmonium production occurs through three mechanisms: direct production, the decay of heavier charmonia, and the decay of $`b`$-flavored hadrons. The first two mechanisms are collectively known as “prompt” because they are observed to occur at the $`p\overline{p}`$ interaction point. The Collider Detector at Fermilab (CDF) collaboration previously reported results on the production of $`J/\psi `$ and $`\psi (2S)`$ mesons . The measured cross sections for direct production were on the order of 50 times larger than predicted by the Color Singlet Model (CSM) . However, calculations based on the Nonrelativistic QCD (NRQCD) factorization formalism are able to account for the observed cross sections by including color octet production mechanisms. This leads to the prediction that directly produced $`\psi `$ mesons will be increasingly transversely polarized at high $`P_T`$ . (In this Letter we use $`\psi `$ to denote either $`J/\psi `$ or $`\psi (2S)`$ mesons.) This is because the production of $`\psi `$ mesons with $`P_TM_\psi `$ is dominated by gluon fragmentation. It is predicted that the gluon’s transverse polarization is preserved as the $`c\overline{c}`$ pair evolves into a bound state $`\psi `$ meson. On the other hand, the Color Evaporation Model predicts an absence of polarization . In this Letter, we report on measurements of the polarization of promptly produced $`\psi `$ mesons at CDF. Our analysis also yields as a byproduct the effective polarization of the $`\psi `$ mesons produced in $`B`$-hadron decays. CDF is a multi-purpose detector designed to study high energy $`p\overline{p}`$ collisions produced by the Tevatron . The CDF coordinate system is defined with the $`z`$ axis along the proton beam direction. The polar angle $`\theta `$ is defined relative to the $`z`$ axis, $`r`$ is the perpendicular radius from this axis, and $`\varphi `$ is the azimuthal angle. Pseudorapidity is defined as $`\eta \mathrm{ln}[\mathrm{tan}(\theta /2)]`$. Three charged-particle tracking detectors immersed in a 1.4 T solenoidal magnetic field surround the beamline. This tracking system is contained within a calorimeter, while drift chambers outside the calorimeter identify muon candidates. The innermost tracking device is a four-layer silicon microstrip detector (SVX) located at radii between 2.9 and 7.9 cm from the beam axis. The SVX is surrounded by a set of time projection chambers (VTX) extending out to a radius of 22 cm. An 84 layer cylindrical drift chamber (CTC) measures the particle trajectories in the region between 30 and 132 cm from the beam. This tracking system has high efficiency for detecting charged particles with momentum transverse to the beam $`P_T>0.4`$ GeV$`/c`$ and \|η\| < 1.1 < 𝜂1.1\|\eta\|\mathrel{\vbox{\kern 0.0pt\hbox{$<$} \kern 0.0pt\hbox{$\sim$} }}1.1. Together, the CTC and SVX measure charged particle transverse momenta with a precision of $`\sigma _{P_T}/P_T=0.0070.001P_T`$ (with $`P_T`$ in GeV$`/c`$). The impact parameter resolution is $`\sigma _d=(13+40/P_T)\mu `$m for tracks with SVX and CTC information. The central muon detection system consists of four layers of planar drift chambers separated from the interaction point by five interaction lengths of material. This system detects muons with PT > 1.4 > subscript𝑃𝑇1.4P_{T}\mathrel{\vbox{\kern 0.0pt\hbox{$>$} \kern 0.0pt\hbox{$\sim$} }}1.4 GeV$`/c`$ and \|η\| < 0.6 < 𝜂0.6\|\eta\|\mathrel{\vbox{\kern 0.0pt\hbox{$<$} \kern 0.0pt\hbox{$\sim$} }}0.6. Dimuon candidates used in this analysis are collected using a three-level $`\mu ^+\mu ^{}`$ trigger. The first level trigger requires that two candidates be observed in the muon chambers. For each muon candidate the first level trigger efficiency rises from $``$40% at $`P_T=1.5`$ GeV$`/c`$ to $``$93% for muons with $`P_T>3.0`$ GeV$`/c`$. The second level trigger requires one or more charged particle tracks in the CTC, reconstructed using the central fast track processor (CFT). The CFT performs a partial reconstruction of all charged tracks with $`P_T`$ above $``$2 GeV$`/c`$. Muon candidates found by the first level trigger are required to match a CFT track within 15 degrees in azimuth. The third level trigger performs three-dimensional track reconstruction and accepts dimuon masses in a broad window around the $`J/\psi `$ and $`\psi (2S)`$ masses. The data used in this study correspond to an integrated luminosity of 110 pb<sup>-1</sup> and were collected between 1992 and 1995. Following the online data collection, additional requirements are made offline to identify the signals and to reduce the backgrounds. To identify muon candidates and reduce the rate from sources such as $`\pi /K`$ meson decay-in-flight, we require that each track observed in the muon chambers be associated with a matching CTC track. These matches are required to pass a maximum $`\chi ^2`$ cut of 9 and 12 (for 1 degree of freedom) in the $`\varphi `$ and $`z`$ views respectively. Also, we require $`P_T`$ greater than about 2 GeV$`/c`$ for each muon candidate. This requirement ensures that the muon trigger and reconstruction efficiencies are well understood, to avoid biases in the decay angular distributions of the charmonia states studied below. The measurement of the polarization of $`\psi `$ mesons is made by analyzing their decays to $`\mu ^+\mu ^{}`$ in the helicity basis, in which the spin quantization axis lies along the $`\psi `$ direction in the $`p\overline{p}`$ center-of-mass (lab) frame. We define $`\theta ^{}`$ as the angle between the $`\mu ^+`$ direction in the $`\psi `$ rest frame and the $`\psi `$ direction in the lab frame. The normalized angular distribution $`I(\mathrm{cos}\theta ^{})`$ is given by $$I(\mathrm{cos}\theta ^{})=\frac{3}{2(\alpha +3)}(1+\alpha \mathrm{cos}^2\theta ^{})$$ (1) Unpolarized $`\psi `$ mesons have $`\alpha =0`$ whereas $`\alpha =+1`$ or $`1`$ correspond to fully transverse or longitudinal polarizations respectively. Experimentally, the acceptance is severely reduced as $`\|\mathrm{cos}\theta ^{}\|`$ approaches 1, due to the $`P_T`$ cuts on the muons. Our method for determining $`\alpha `$ is to fit the observed distributions of $`\mathrm{cos}\theta ^{}`$ to distributions derived from simulated $`\psi \mu ^+\mu ^{}`$ decays. The Monte Carlo simulation accounts for the geometric and kinematic acceptance of the detector as well as the reconstruction efficiency as a function of $`\mathrm{cos}\theta ^{}`$. In order to extract the polarization parameter $`\alpha `$ for promptly produced $`\psi `$ mesons, we separate the prompt component from the $`B`$-decay component using the proper decay length of each event. For $`\psi `$ candidates with one or both muons reconstructed in the SVX (the SVX sample), we define a vector pointing from the $`p\overline{p}`$ collision point to the $`\psi `$ decay vertex. The transverse decay length $`L_{xy}`$ is then defined as the projection of this vector onto the $`\psi `$ transverse momentum. The proper decay length $`ct`$ is related to the transverse decay length by $`ct=(M_\psi L_{xy})/(F_{corr}^\psi P_T^\psi )`$, where $`M_\psi `$ is the $`\psi `$ mass. Here $`F_{corr}^\psi `$ is a correction factor obtained from Monte Carlo studies , which accounts for the fact that we are using the $`\psi `$ $`P_T`$ instead of the $`B`$-hadron $`P_T`$. Prompt events have $`ct`$ consistent with zero whereas $`B`$-decays have an exponential $`ct`$ distribution; the detector resolution smears the $`ct`$ distribution. We fit the $`ct`$ distribution to obtain the relative fractions of prompt and $`B`$-decay production. Details of this fitting procedure are given in . The measured fraction of $`J/\psi `$ mesons which come from $`B`$-hadron decay increases from $`(13.0\pm 0.3)\%`$ at $`P_T^{J/\psi }=4\mathrm{G}eV/c`$ to $`(40\pm 2)\%`$ at 20 $`\mathrm{G}eV/c`$. For $`\psi (2S)`$ mesons, an increase from $`(21\pm 2)\%`$ to $`(35\pm 4)\%`$ is seen in the range from 5.5 to 20 $`\mathrm{G}eV/c`$. The proper decay length measurement allows us to divide the data into two samples: a short-lived sample dominated by prompt production, and a long-lived sample dominated by $`B`$-decays. The short-lived sample is defined by $`0.1ct0.013[0.01]`$ cm, and the long-lived sample by $`0.013[0.01]ct0.3`$ cm, for the $`J/\psi `$ $`[\psi (2S)]`$ analysis respectively. The boundary between the two $`ct`$ regions has been optimized separately for the $`J/\psi `$ and $`\psi (2S)`$ samples, to maximize the purity of prompt decays in the short-lived sample and $`B`$-decays in the long-lived sample. Depending on $`P_T^\psi `$, the prompt fraction in the short-lived sample ranges from 85%\[86%\] to 96%\[95%\], and the $`B`$-decay fraction in the long-lived sample ranges from 83%\[86%\] to 98%\[91%\], for $`J/\psi [\psi (2S)]`$ respectively. The $`J/\psi `$ polarization is measured in seven $`P_T`$ bins, covering a range of $`420\mathrm{G}eV/c`$. Using a 3 standard deviation mass window around the $`J/\psi `$ peak, our data sample consists of 180,000 signal $`J/\psi `$ events, with a signal to background ratio of about 13. The $`J/\psi `$ sample is divided into three subsamples: the short-lived and long-lived SVX samples described above, and a third sample (the CTC sample) in which neither muon has SVX information and no $`ct`$ measurement is made. In each $`P_T^{J/\psi }`$ bin, we measure the prompt polarization ($`\alpha _P`$) and the effective polarization of $`J/\psi `$ mesons from $`B`$-hadron decays ($`\alpha _B`$). (We refer to $`\alpha _B`$ as “effective” because $`\theta ^{}`$ is defined using the lab frame, not the $`B`$-hadron rest frame — in effect this dilutes any polarization from the $`B`$-decay toward zero.) We find that it is not feasible to make separate polarization measurements for direct $`J/\psi `$ production and for production from $`\chi _c`$ and $`\psi (2S)`$ decays. The latter sources account for 36$`\pm `$6% of the prompt component, with only a small $`P_T^{J/\psi }`$ dependence . The $`J/\psi `$ polarization is measured by fitting $`\mathrm{cos}\theta ^{}`$ distributions in data to a set of Monte Carlo templates . The templates are generated by processing simulated samples of $`J/\psi \mu ^+\mu ^{}`$ decays with a detector and trigger simulation. The polarization is obtained using a $`\chi ^2`$ fit of the data to a weighted sum of transversely polarized and longitudinally polarized templates. The fitted weights yield the polarization. Two transverse/longitudinal template pairs are generated, using measured prompt and $`B`$-decay $`P_T^{J/\psi }`$ spectra . The $`\mathrm{cos}\theta ^{}`$ distribution of background events is modeled in the fit using sidebands around the $`J/\psi `$ mass peak. The fit is performed simultaneously on the SVX short-lived, SVX long-lived and CTC samples, with two fit parameters: $`\alpha _P`$ and $`\alpha _B`$. To account for the mixture of prompt and $`B`$-decay components in each sample, the relative fractions of prompt and $`B`$-decay templates in each are fixed in the fit using the results of the lifetime fit. The B-decay fraction in the CTC sample is assumed to be the same as in the SVX sample, because the two samples differ primarily in the $`z`$ position of the primary vertex. Within each $`P_T^{J/\psi }`$ bin, a small correction is applied to the $`P_T^{J/\psi }`$ distributions of the Monte Carlo samples so that they match with those in the data. As an example, the fit in the $`P_T`$ range 12-15 $`\mathrm{G}eV/c`$ is shown in Fig. 1. Three sources of systematic uncertainty are evaluated: the trigger efficiency, the fitted prompt and $`B`$-decay fractions, and the $`P_T^{J/\psi }`$ spectra used in making the Monte Carlo templates. Except in the lowest $`P_T`$ bins, the systematic uncertainties are much smaller than the statistical uncertainties. Our fit results are listed in Table I, and $`\alpha _P`$ is compared with a theoretical NRQCD prediction in Fig. 2. The measurement of the $`\psi (2S)`$ polarization is made in three $`P_T`$ bins covering $`5.520.0\mathrm{G}eV/c`$. Both muons are required to be reconstructed in the SVX. The resulting dimuon mass distribution is fitted with a Gaussian signal and a linear background. We find a total of $`1855\pm 65`$ signal $`\psi (2S)`$ events, with a signal to background ratio of about 1 in a $`3`$ standard deviation mass window around the $`\psi (2S)`$ mass. As discussed above, the sample in each $`P_T`$ bin is further divided into two subsamples based on the $`ct`$ distribution. Because the statistics are lower than in the $`J/\psi `$ case, we use 10 bins in $`\|\mathrm{cos}\theta ^{}\|`$. The number of signal events in each $`\|\mathrm{cos}\theta ^{}\|`$ bin is obtained by fitting its mass distribution. The resulting $`\|\mathrm{cos}\theta ^{}\|`$ distributions in the two $`ct`$ subsamples are fitted simultaneously to the predicted number of events to extract the $`\psi (2S)`$ polarizations for prompt and $`B`$-decay production. The number of predicted events in each $`\|\mathrm{cos}\theta ^{}\|`$ bin is derived by weighting the normalized angular distribution $`I(\mathrm{cos}\theta ^{})`$ with the detector acceptance . We use the measured prompt and $`B`$-decay $`P_T^{\psi (2S)}`$ distributions to calculate the acceptance. As in the $`J/\psi `$ case, there is a small correlation between the measured $`P_T^{\psi (2S)}`$ distributions and the polarization. A correction is applied iteratively in the fits to account for this dependence. Figure 3 shows the observed angular distributions with their polarization fits for the short-lived sample in the three $`P_T^{\psi (2S)}`$ bins. Three sources of systematic uncertainty are considered: the uncertainty in the event yield from the mass fits in the $`\|\mathrm{cos}\theta ^{}\|`$ bins, the uncertainty due to the error on the fitted prompt and $`B`$-decay fractions, and the uncertainty on the $`\|\mathrm{cos}\theta ^{}\|`$ acceptance from the Monte Carlo modeling of the $`P_T^{\psi (2S)}`$ distributions. The uncertainty due to the trigger efficiency is negligible in the $`P_T^{\psi (2S)}`$ range used. The systematic uncertainties are much smaller than the statistical uncertainties. The fitted values of $`\alpha _P`$ and $`\alpha _B`$ as a function of $`P_T^{\psi (2S)}`$ are listed in Table II, and $`\alpha _P`$ is shown in Fig. 2 with the NRQCD predictions . In conclusion, we have measured the polarization of $`J/\psi `$ and $`\psi (2S)`$ mesons produced in 1.8 TeV $`p\overline{p}`$ collisions. The polarization from $`B`$-decays is generally consistent with zero, as expected. In both the $`J/\psi `$ and $`\psi (2S)`$ cases, we do not observe increasing prompt transverse polarization at PT > 12GeV/c > subscript𝑃𝑇12G𝑒𝑉𝑐P_{T}\mathrel{\vbox{\kern 0.0pt\hbox{$>$} \kern 0.0pt\hbox{$\sim$} }}12\ {\mathrm{G}eV}\!/c. Our measurements are limited by statistics, especially for the $`\psi (2S)`$, but they appear to indicate that no large transverse prompt polarization is present at high $`P_T`$, in disagreement with NRQCD factorization predictions. We thank the Fermilab staff and the technical staffs of the participating institutions for their vital contributions. This work was supported by the U.S. Department of Energy and National Science Foundation; the Italian Istituto Nazionale di Fisica Nucleare; the Ministry of Education, Science, Sports and Culture of Japan; the Natural Sciences and Engineering Research Council of Canada; the National Science Council of the Republic of China; the Swiss National Science Foundation; the A. P. Sloan Foundation; and the Bundesministerium für Bildung und Forschung, Germany.
warning/0004/cond-mat0004361.html	ar5iv	text	# Local functional models of critical correlations in thin-films ## Abstract Recent work on local functional theories of critical inhomogeneous fluids and Ising-like magnets has shown them to be a potentially exact, or near exact, description of universal finite-size effects associated with the excess free-energy and scaling of one-point functions in critical thin films. This approach is extended to predict the two-point correlation function $`G`$ in critical thin-films with symmetric surface fields in arbitrary dimension $`d`$. In $`d=2`$ we show there is exact agreement with the predictions of conformal invariance for the complete spectrum of correlation lengths $`\xi ^{(n)}`$ as well as the detailed position dependence of the asymptotic decay of $`G`$. In $`d=3`$ and $`d4`$ we present new numerical predictions for the universal finite-size correlation length and scaling functions determining the structure of $`G`$ across the thin-film. Highly accurate analytical closed form expressions for these universal properties are derived in arbitrary dimension. Finite-size scaling theory predicts that the free-energy, order-parameter and correlation length of a critical thin-film Ising-like magnet (or of a fluid between parallel-plates at the bulk critical point) are characterised by universal amplitudes and scaling functions that depend only on dimensionality $`d`$ and the qualitative nature of of the boundary conditions . These have been the topic of widespread theoretical interest with approaches including sophisticated $`\epsilon `$-expansion treatments as $`d4`$ and in particular studies that invoke conformal invariance in $`d=2`$ . In a recent article Boran and Upton (BU) have revisited the problem of critical thin-films using a local-functional model similar to one introduced earlier by Fisher and Upton . Specifically BU consider critical films of thickness $`L`$ and infinite area $`A`$ with both like sign $`(++)`$ and opposite sign $`(+)`$ symmetry-breaking surface fields and obtain explicit expressions for both the universal Casimir amplitude $`A_{ab}`$ associated with the finite-size free-energy and universal scaling function $`\mathrm{\Psi }_{ab}(x/L)`$ for the equilibrium magnetisation profile $`m(x)`$, for arbitrary dimension $`d`$. For the $`(++)`$ geometry in particular where there is less ambiguity in the choice of functional, the predictions of BU are in astonishing agreement with previous results obtained from the $`\epsilon `$-expansion and conformal invariance. For example in $`d=2`$ the local functional result for $`\mathrm{\Psi }_{++}(x/L)`$ is exact whilst the numerical value obtained for $`A_{++}`$ is in near perfect agreement with the conformal result $`A_{++}=\pi /48`$ . Whilst the local functional theory is phenomenological in nature, the agreement with microscopic theory is compelling and the model offers the exciting possibility of predicting (near) exact results for universal amplitudes and scaling functions in $`d=3`$. In the present work we extend the local functional theory to allow calculation of the two-point connected correlation function $`G(𝐫_\mathrm{𝟏},𝐫_\mathrm{𝟐})=m(𝐫_\mathrm{𝟏})m(𝐫_\mathrm{𝟐})m(𝐫_\mathrm{𝟏})m(𝐫_\mathrm{𝟐})`$ for the $`(++)`$ thin-film. First we discuss the $`d=2`$ Ising strip and compare the theory with exact results obtained from conformal invariance. This is a stringent test of the accuracy of the local functional theory and our analysis shows an extraordinary level of agreement which we find even more striking than the successes obtained by BU for $`A_{++}`$ and $`\mathrm{\Psi }_{++}(x/L)`$ mentioned above. In particular we derive the precise spectrum of universal correlation lengths as well as the exact and highly non-trivial position dependence of the leading-order and next-to-leading order asymptotic decay of $`G`$. We are also able to obtain the exact Casimir amplitude $`A_{++}`$, improving on the estimate of BU. Secondly we present detailed predictions for the universal finite-size correlation length and scaling function $`\mathrm{\Xi }_d(x/L)`$ characterising the asymptotic decay of $`G`$ for the $`d=3`$ and $`d4`$ thin-films. These features of $`G`$ have not been discussed in detail before and we are able to derive accurate analytic expressions for them valid in all dimensions. To begin we recall some pertinent universal finite-size scaling properties of critical thin-films and the exact predictions for $`G`$ in $`d=2`$ arising from conformal invariance. Consider a $`d`$ dimensional Ising-like thin film with surface fields $`h_1`$, $`h_2`$ acting on the spins in the surface planes $`x=0`$ and $`x=L`$. These fields mimic the preferential local adsorption of the confining walls on a fluid in a parallel-plate geometry. We define the reduced excess free-energy $`f^\times (L)(FALf_b)/Ak_BT`$ by subtracting from the total free-energy $`F`$ the appropriate contribution from the bulk free-energy density $`f_b`$. By further subtracting from $`f^\times (L)`$ the semi-infinite (reduced) surface tensions of the two-independent surfaces we obtain the desired finite-size contribution to the free-energy, $`\mathrm{\Delta }f^\times (l)`$. Finite size scaling theory predicts that exactly at the bulk critical point, for asymptotically large widths $`L`$ (and dimension $`d<4`$), this can be written $`\mathrm{\Delta }f^\times (L)=A_{ab}/L^{d1}+\mathrm{}`$ where the Casimir amplitude $`A_{ab}`$ depends only on whether the surface fields have like $`(++)`$ or opposite $`(+)`$ sign (or are zero). Hereafter we restrict attention to the $`(++)`$ geometry where the results of the local functional theory are most striking. A second universal quantity is the normalised scaling function $`\mathrm{\Psi }_{++}(x/L)`$ describing the magnetisation profile $`m(x)`$ across the thin-film. For a semi-infinite system the asymptotic large distance decay of the critical magnetisation profile is $`m(x)Cx^{\beta /\nu }`$ where $`C`$ is a non-universal amplitude and $`\beta `$, $`\nu `$ are standard bulk critical exponents for the magnetisation and correlation length respectively. In the $`(++)`$ critical thin-film finite-size scaling predicts that for large $`x`$, $`L`$ with $`x/L`$ arbitrary the profile can be written $`m(x)=m(L/2)\mathrm{\Psi }_{++}(x/L)`$ where the mid-point magnetisation scales as $`m(L/2)L^{\beta /\nu }`$ and $`\mathrm{\Psi }_{++}(x/L)`$ is normalised such that $`\mathrm{\Psi }_{++}(1/2)=1`$. Finally, finite-size scaling theory also predicts that the asymptotic decay of correlations along the critical thin-film is described by a (true) correlation length $`\xi ^{(0)}`$, with $`L/\xi ^{(0)}`$ a universal critical amplitude ratio. It is well known that in $`d=2`$ these predictions are beautifully confirmed by conformal invariance and are consistent with results obtained from exactly solvable models. In particular by using the standard logarithmic function to conformally map the semi-infinite plane to the finite-size strip the magnetisation follows as $`m(x)=C(\frac{L}{\pi }\mathrm{sin}\frac{\pi x}{L})^{1/8}`$ implying $`\mathrm{\Psi }_{++}(z)=(\mathrm{sin}\pi x)^{1/8}`$. Applying the same mapping to the semi-infinite two-point function calculated by Burkhardt and Xue we obtain $`G^{++}(x_1,x_2;y)`$ $`=`$ $`{\displaystyle \frac{C^2}{\sqrt{2}}}\left({\displaystyle \frac{L}{\pi }}\right)^{1/4}\left(\mathrm{sin}{\displaystyle \frac{\pi x_1}{L}}\mathrm{sin}{\displaystyle \frac{\pi x_2}{L}}\right)^{1/8}\times `$ (2) $`\left[\sqrt{u+u^1}\sqrt{2}\right]`$ $`\text{where}u=\left[{\displaystyle \frac{\mathrm{cosh}\pi y/L\mathrm{cos}\pi (x_2+x_1)/L}{\mathrm{cosh}\pi y/L\mathrm{cos}\pi (x_2x_1)/L}}\right]^{1/4}`$ (3) and $`y`$ denotes the scalar distance between the spins measured along the strip. From the asymptotic $`y/L\mathrm{}`$ expansion of (2) it follows that the true correlation length in the $`(++)`$ strip is $`\xi ^{(0)}=L/2\pi `$ identifying the universal critical amplitude ratio predicted by finite-size scaling theory. More generally, this expansion identifies the spectrum of correlation lengths $`\xi ^{(n)}`$ determining the decay of the exponential terms as $$\xi ^{(n)}=\frac{L}{(n+2)\pi };n=0\text{}1\text{}2\text{}3\text{}\mathrm{}.$$ (4) Thus the explicit conformal invariance prediction for the leading-order and next-to-leading-order asymptotic $`y/L\mathrm{}`$ decay of $`G^{++}`$ is $`G^{++}(x_1,x_2;y)={\displaystyle \frac{C^2}{4}}\left({\displaystyle \frac{L}{\pi }}\right)^{1/4}\left(\mathrm{sin}{\displaystyle \frac{\pi x_1}{L}}\mathrm{sin}{\displaystyle \frac{\pi x_2}{L}}\right)^{15/8}\times `$ (5) $`e^{2\pi y/L}\left\{1+4\mathrm{cos}{\displaystyle \frac{\pi x_1}{L}}\mathrm{cos}{\displaystyle \frac{\pi x_2}{L}}e^{\pi y/L}+\mathrm{}\right\}`$ (6) which shows a detailed position dependence across the strip. We wish to compare these exact predictions with the results of the BU local functional theory. The desirable scaling properties of quite general local functionals of the magnetisation have been considered by Fisher and Upton but in application to the $`(++)`$ thin-film (where there is no zero of the equilibrium profile) the appropriate model is simplicity itself. BU assume translational invariance of the trial profiles along the thin-film and base their analysis on the excess (reduced) free-energy functional (ignoring boundary terms) $$[m]=_0^L𝑑x\left[\frac{K}{2}m^{\eta \nu /\beta }\dot{m}^2+\frac{u}{\delta +1}m^{\delta +1}\right]$$ (7) where $`\dot{m}dm/dx`$, $`\eta `$ and $`\delta `$ are standard bulk exponents and $`K`$, $`u`$ are parameters that must be regarded as inputs into the theory, although their values do not influence the determination of the scaling function $`\mathrm{\Psi }_{++}(z)`$ nor the structural properties of $`G`$. The surface fields enter through boundary terms but these play no role in determining the universal quantities of interest and may be replaced by fixed boundary conditions $`m(0)=m(L)=\mathrm{}`$ as appropriate to the scaling limit. Specialising to $`d=2`$ and using the bulk critical exponents $`\beta =1/8`$, $`\nu =1`$, $`\eta =1/4`$ and $`\delta =15`$ the Euler-Lagrange equation for the equilibrium profile is $$Km^2\ddot{m}Km^3\dot{m}^2=um^{15}$$ (8) which has a first integral $$\frac{Km^2\dot{m}^2}{2}=\frac{um^{16}}{16}\mathrm{\Delta }p$$ (9) where the constant of integration $`\mathrm{\Delta }p=f^\times /L=A_{++}L^2`$. From (5) it is straight forward to derive $$m(x)=\left(\frac{K}{8u}\right)^{1/16}\left(\frac{L}{\pi }\mathrm{sin}\frac{\pi x}{L}\right)^{1/8}$$ (10) in exact accordance with conformal invariance. From the profile it follows that $`\mathrm{\Delta }p=\pi ^2K/128L^2`$ implying $`A_{++}=\pi ^2K/128`$. Using scaling arguments BU relate $`K`$ to known bulk critical exponent and amplitude ratios and yield a numerical value for $`A_{++}`$ extremely close to the exact result $`A_{++}=\pi /48`$ and with an error that may even be attributed to uncertainties in the values of bulk critical amplitude ratios used to determine $`K`$. We shall determine $`K`$ another way and show how the exact result for $`A_{++}`$ may be obtained within the local functional theory. With these preliminaries over we can now turn to the main part of our work and describe the calculation of $`G^{++}`$ within the local functional theory. As it stands the BU functional (4) is not capable of yielding the full correlation function because it assumes translational invariance of the magnetisation along the strip. However it is reasonable to argue from isotropy that the full functional is given by $$[m]=\frac{1}{A}𝑑y𝑑x\left[\frac{K}{2}m^2(m)^2+\frac{u}{16}m^{16}\right]$$ (11) and recall that we have adopted fixed boundary conditions as appropriate to the scaling limit. Minimisation of (8) leads to the extended Euler-Lagrange equation $$Km^2^2mKm^3(m)^2=um^{15}$$ (12) which of course reduces to (5) as soon as we impose translational invariance along the strip. We mention here that the solution to (9) for the magnetisation profile interior to a disk of radius R with fixed up-spins at the boundary is $`m(r)=C(\frac{R}{2}(1r^2/R^2))^{1/8}`$ which is identical to the exact conformal invariance result . The same is also true in all higher dimensions $`d`$ if we adopt the appropriate values for the bulk critical exponents in (9). Within the framework of density-functional theory the two-point correlation function follows from solution of the Ornstein-Zernike equation $$\mathrm{𝐝𝐫}_\mathrm{𝟑}\frac{\delta ^2[m]}{\delta m(𝐫_\mathrm{𝟏})\delta m(𝐫_\mathrm{𝟑})}G(𝐫_\mathrm{𝟐},𝐫_\mathrm{𝟑})=\delta (𝐫_\mathrm{𝟏}𝐫_\mathrm{𝟐}).$$ (13) For our functional this reduces to the differential equation $`[{\displaystyle \frac{135}{8}}um^{14}+2m^2\mathrm{\Delta }pK{\displaystyle \frac{}{x_1}}m^2{\displaystyle \frac{}{x_1}}`$ (14) $`Km^2{\displaystyle \frac{^2}{y^2}}]G^{++}(x_1,x_2;y)=\delta (x_1x_2)\delta (y)`$ (15) which we seek to solve using the spectral expansion $$G^{++}(x_1,x_2;y)=\underset{n=0}{\overset{\mathrm{}}{}}\psi _n(x_1)\psi _n^{}(x_2)e^{E_ny}$$ (16) with energy levels $`E_n`$ equivalent to the inverse correlation lengths $`1/\xi ^{(n)}`$ discussed earlier. Orthonormality implies $$2K_0^L𝑑x\psi _n^{}(x)m(x)^2\psi _m(x)=\frac{\delta _{nm}}{E_n}$$ (17) which is needed to determine the correct normalisation factors of the eigenstates. The spectrum is thus obtained through solution of the eigenvalue problem $`\left[{\displaystyle \frac{135}{8}}um^{16}+2\mathrm{\Delta }pKm^2{\displaystyle \frac{}{x}}m^2{\displaystyle \frac{}{x}}\right]\psi _n(x)`$ (18) $`=KE_n^2\psi _n(x)`$ (19) subject to the boundary condition $`\psi _n(0)=\psi _n(L)=0`$. Substituting for $`m(x)`$ and the value of $`\mathrm{\Delta }p`$ quoted earlier (without having to fix $`K`$) this reduces to $`\psi _n^{\prime \prime }(x){\displaystyle \frac{\pi }{4L}}\left(\mathrm{cot}{\displaystyle \frac{\pi x}{L}}\right)\psi _n^{}(x)+`$ (20) $`\left[{\displaystyle \frac{135\pi ^2}{64L^2}}\left(\mathrm{sin}{\displaystyle \frac{\pi x}{L}}\right)^2+{\displaystyle \frac{\pi ^2}{64L^2}}\right]\psi _n(x)`$ $`=`$ $`E_n^2\psi _n(x).`$ (21) which can be solved using standard methods . For the eigenvalue spectrum we find $$E_n=\frac{(n+2)\pi }{L};n=0\text{}1\text{}2\text{}3\text{}\mathrm{}$$ (22) whilst the eigenvectors can be written as $`\psi _n(x)=N_n(\mathrm{sin}\frac{\pi x}{L})^{15/8}\varphi _n(x)`$ with $$\varphi _n(x)=F[\frac{n}{2},2+\frac{n}{2},\frac{1}{2};\mathrm{cos}^2\frac{\pi x}{L}]$$ (23) for the even states and $$\varphi _n(x)=\mathrm{cos}\frac{\pi x}{L}F[\frac{1}{2}\frac{n}{2},\frac{5}{2}+\frac{n}{2},\frac{3}{2};\mathrm{cos}^2\frac{\pi x}{L}]$$ (24) for the odd states. Here $`F[\alpha ,\beta ,\gamma ;z]`$ denotes the usual hypergeometric function and $`N_n`$ is an appropriate normalisation factor which can be found using (13). The eigenvalue spectrum (16) is identical to the result obtained from conformal invariance (2) and exactly identifies all the universal critical amplitude ratios for the correlation lengths. From the above the asymptotic decay of the correlation function within the local functional theory is given by $`G^{++}(x_1,x_2;y)={\displaystyle \frac{C^2}{4}}{\displaystyle \frac{8}{3K\pi }}\left(\mathrm{sin}{\displaystyle \frac{\pi x_1}{L}}\mathrm{sin}{\displaystyle \frac{\pi x_2}{L}}\right)^{15/8}\times `$ (25) $`e^{2\pi y/L}\left\{1+\left({\displaystyle \frac{N_1}{N_0}}\right)^2\mathrm{cos}{\displaystyle \frac{\pi x_1}{L}}\mathrm{cos}{\displaystyle \frac{\pi x_2}{L}}e^{\pi y/L}+\mathrm{}\right\}`$ (26) where $`C=(K/8u)^{1/16}`$ is the non-universal constant for the magnetisation while from (13) we find $`N_1/N_0=2`$ in precise agreement with the conformal result (3)! Moreover it is clear that the parameter $`K`$ simply sets the overall scale of the correlation function and that in $`d=2`$ the appropriate value is $`K=3\pi /8`$. Utilising this identification we are led to the prediction $`A_{++}=\pi /48`$ which is identical to the exact conformal invariance result . We comment here that only small discrepancies with the conformal result are found at higher orders which are surely of no practical importance. It is also possible to calculate the moments $`G_{2n}^{++}(x_1,x_2)𝑑yy^{2n}G^{++}(x_1,x_2;y)`$ using the local functional theory and compare these with the expressions obtained from conformal invariance. Consider for example the zeroth moment $`G_0(x_1,x_2)`$ from which we can construct a universal two-point scaling-function $`\mathrm{\Lambda }_0^{++}(x_1/L,x_2/L)=G_0^{++}(x_1,x_2)/G_0^{++}(L/2,L/2)`$. We omit the details of this calculation and simply quote our final result calculated within the BU functional $`\mathrm{\Lambda }_0^{++}({\displaystyle \frac{x_1}{L}},{\displaystyle \frac{x_2}{L}})=\left(\mathrm{sin}{\displaystyle \frac{\pi x_1}{L}}\mathrm{sin}{\displaystyle \frac{\pi x_2}{L}}\right)^{1/8}\times `$ (27) $`\left[{\displaystyle \frac{\pi x_1}{L}}\mathrm{cot}{\displaystyle \frac{\pi x_1}{L}}\right]\left[\pi \left({\displaystyle \frac{x_2}{L}}1\right)\mathrm{cot}{\displaystyle \frac{\pi x_2}{L}}1\right]`$ (28) which turns out to be identical to the conformally invariant result obtained by (numerically) integrating (2) provided that $`x_1=x_2`$. There are small discrepancies between the expressions if $`x_1x_2`$ but even here the local functional theory is an excellent qualitative description of the structure of correlations across the strip. Nevertheless this implies that the exact conformal result for $`G_0^{++}(x_1,x_2)`$ does not satisfy the correlation function algebra necessitated by all local functional theories . Given that it is therefore not possible to construct a truly exact local functional theory consistent with conformal invariance predictions for $`G^{++}(x_1,x_2;y)`$, the level of agreement achieved by the theory is all the more striking. Next we turn our attention to the $`d`$ dimensional thin-film and focus on the decay of $`G^{++}(x_1,x_2;𝐲)`$ for asymptotic parallel displacements $`\|𝐲\|\mathrm{}`$. Within the local functional theory this behaves as $$G^{++}(x_1,x_2,𝐲)C^2\mathrm{\Xi }_d\left(\frac{x_1}{L}\right)\mathrm{\Xi }_d\left(\frac{x_2}{L}\right)\frac{e^{\|𝐲\|/\xi ^{(0)}}}{\|𝐲\|^{(d2)/2}}$$ (29) and recall that we anticipate that $`L/\xi ^{(0)}`$ is universal. Here we have introduced a universal scaling function $`\mathrm{\Xi }_d(x/L)`$ describing the position dependence across the thin-film which we normalise (analogous to $`\mathrm{\Psi }_{++}`$) such that $`\mathrm{\Xi }_d(1/2)=1`$. Both $`\mathrm{\Xi }_d(x/L)`$ and $`\xi ^{(0)}`$ are determined by an eigenvalue equation analogous to (15) which depends only on $`d`$, $`\beta /\nu `$ and $`\mathrm{\Psi }_{++}(x/L)`$. No information concerning $`K`$ or $`u`$ is required. We omit details of the calculation which are obtained assuming the same numerical values for the critical exponents $`\beta =0.328`$ and $`\nu =0.632`$ in $`d=3`$ as used by BU and the standard mean-field results $`\beta =\nu =1/2`$ for $`d4`$. For the universal correlation length critical amplitude we obtain $`L/\xi ^{(0)}7.34`$ and $`L/\xi ^{(0)}8.40`$ in $`d=3`$ and $`d4`$ respectively and recall $`L/\xi ^{(0)}=2\pi `$ in $`d=2`$. The numerical results for the scaling functions $`\mathrm{\Xi }_d(x/L)`$ are quite remarkable and are shown in Fig. 1. To extremely high precision (see inset) and for all practical purposes these functions are completely indistinguishable from the analytic expression $$\stackrel{~}{\mathrm{\Xi }}_d(x/L)=\left(\mathrm{sin}\frac{\pi x}{L}\right)^{d^{}\beta /\nu }$$ (30) where $`d^{}=\mathrm{min}\{d,4\}`$. Notice this is exact in $`d=2`$ and embodies the correct short-distance expansion behaviour close to each wall $`d`$. Using $`\stackrel{~}{\mathrm{\Xi }}_d(x/L)`$ we can derive an extremely accurate expression for the ground state eigenvalue and hence obtain for the universal ratio $$\frac{L}{\xi ^{(0)}}=\sqrt{2d^{}\left(2+\frac{\beta }{\nu }\right)J_{++}(1)^2+\pi ^2\left(d^{}\frac{\beta }{\nu }\right)}$$ (31) where $`J_{++}(1)=_0^1𝑑u(1u^d^{})^{1/2}`$ is the elliptic function introduced by BU. This is exact in $`d=2`$ and within $`0.3\%`$ of the numerical values quoted above for $`d=3`$ and $`d4`$. For all intents and purposes this may be regarded as the exact analytical local functional prediction for the universal critical amplitude ratio. In summary we have established that the BU local functional model is an excellent description of correlations in the $`d=2`$ critical $`++`$ thin-film and used the theory to make predictions for the universal structure of the correlation function in higher dimensional systems. EDM and CR acknowledge support from the EPSRC and EC (ERBFM-BICT983229) respectively.
warning/0004/hep-th0004077.html	ar5iv	text	# 1 Introduction ## 1 Introduction The study of black hole solutions in $`N=2`$ five-dimensional supergravity coupled to vector and hypermultiplets plays an important role in the understanding of the non-perturbative structure of string and M-theory . In this setting the interplay between classical and quantum results is exemplified at its best. In this paper we consider general charged black holes of the $`D=5`$, $`N=2`$ theories, not necessarily those obtained from compactification of eleven-dimensional supergravity on a Calabi-Yau threefold. The analysis is simplified by the rich geometric structure of the $`N=2`$ theories. Black hole solutions are given in terms of a rescaled cubic homogeneous prepotential which defines very special geometry . In the extremal BPS case, half of the vacuum supersymmetries are preserved, while at the horizon supersymmetry is fully restored . In five dimensions the supergravity action contains a Chern-Simon term which allows the existence of black holes with nonvanishing angular momenta (but still nonrotating horizon) . These issues have been the object of recent, extensive studies. In the present paper we focus on the asymptotic symmetries of the near-horizon geometry of the general near-extremal solution: the aim is the computation of the entropy from a counting of microstates to be compared to the macroscopic, thermodynamical entropy. We will see that the calculation of the microscopic entropy of small excitations above extremality is equivalent to a microstate counting for certain black holes in two-dimensional anti-de Sitter space. The latter, however, is problematic: $`AdS_2`$ has two timelike boundaries, but when applying Cardy’s formula for the density of states only one boundary is taken into account. This procedure leads to a statistical entropy result which is off by a $`\sqrt{2}`$ factor with respect to the Bekenstein-Hawking entropy . Up to now, no satisfactory explanation of this mismatch is known. Our results support the point of view that only the ground state has an effective description in terms of a quantum-mechanical system , whereas the excitations above extremality are described by a two-dimensional conformal field theory . In our work we address these issues in a constructive approach. The main result that we present is an explicit duality transformation, which realizes an invariance of the $`N=2`$ supergravity action. This duality turns the $`AdS_2\times S^3`$ near horizon geometry of the extremal black hole solution into $`AdS_3\times S^2`$. The key point underlying the duality is the fact that the three-sphere can be written as a Hopf fibration over the base $`S^2`$. For $`AdS_3,`$ the counting of microstates is performed using Cardy’s formula and it is shown that it reproduces correctly the Bekenstein-Hawking entropy, thus resolving the discrepancy previously found for $`AdS_2`$. In the case where the $`D=5`$, $`N=2`$ supergravity action is obtained by Calabi-Yau (CY) compactification of M-theory, the considered duality transformation, which maps electrically charged black holes onto magnetically charged black strings, corresponds to the duality between M2 branes wrapping CY two-cycles and M5 branes wrapping CY four-cycles. According to , M-theory compactified on $`AdS_3\times S^2\times M`$, where $`M`$ denotes some Calabi-Yau threefold, is dual to a $`(0,4)`$ superconformal field theory living on an M5 brane wrapping some holomorphic CY four-cycle. This fact has been used in to compute the entropy of five-dimensional BPS black holes<sup>1</sup><sup>1</sup>1The work in includes as a special case also the results obtained in .. We stress that our method for microstate counting applies to any near-extremal black hole in $`N=2`$, $`D=5`$ supergravity, independent of whether it is obtained by CY compactification or not. Our paper is organized as follows: in section 2, the basic notions of $`N=2`$, $`D=5`$ supergravity and very special geometry relevant to our analysis are summarized. In section 3 we review the black hole solutions and consider the $`STU`$ model as a simple example, which nonetheless retains all the interesting features of the general solutions. In section 4 we present the isometry superalgebra which arises in the near horizon limit, while in section 5 we show that the motion of a particle which moves near the horizon of the extremal rotating black hole is described by conformal quantum mechanics. This indicates that the ground state may have a description in terms of conformal quantum mechanics , even when rotation is included. In section 6 we compute the statistical entropy of small excitations near extremality, using the $`AdS_2`$ central charge , and find a $`\sqrt{2}`$ factor of discrepancy as compared to the thermodynamical Bekenstein-Hawking result. In section 7 we construct the duality transformation for the supergravity action, and in the following section we finally perform the state counting, using the fact that the near-horizon geometry of the dual solution includes an $`AdS_3`$ factor. In this way, we obtain a microscopic entropy which agrees precisely with the corresponding thermodynamical result. We conclude with some final remarks. ## 2 $`D=5`$, $`N=2`$ Supergravity and Very Special Geometry The theory of $`N=2`$ supergravity theory coupled to an arbitrary number $`n`$ of Maxwell supermultiplets was first considered in . In the analysis of , it was established that the scalar fields of the vector multiplets parametrize a Riemannian space. The homogeneous symmetric spaces take the form $$=\frac{Str_0(J)}{Aut(J)},$$ where $`Str_0`$($`J)`$ is the reduced structure group of a formally real unital Jordan Algebra of degree three, $`Aut(J)`$ is its automorphism group. The scalar manifold can be regarded as a hypersurface, with vanishing second fundamental form of an $`(n+1)`$-dimensional Riemannian space $`𝒢`$ whose coordinates $`X`$ are in correspondence with the vector multiplets including that of the graviphoton. The equation of the hypersurface is $`𝒱=1`$ where $`𝒱`$, the prepotential, is a homogeneous cubic polynomial in the coordinates of $`𝒢`$, $$𝒱(X)=\frac{1}{6}C_{IJK}X^IX^JX^K.$$ (2.1) Non-simple Jordan algebras of degree three are of the form $`\mathbb{\Sigma }_n`$, where $`\mathbb{\Sigma }_n`$ is the Jordan algebra associated with a quadratic form. The corresponding symmetric scalar manifolds are $$=SO(1,1)\times \frac{SO(n1,1)}{SO(n1)}.$$ (2.2) In this case, $`𝒱(X)`$ is factorizable into a linear times a quadratic form in $`(n1)`$ scalars, which for the positivity of the kinetic terms in the Lagrangian, must have a Minkowski metric. For simple Jordan algebras, one obtains four sporadic locally symmetric spaces related to the four simple unital formally real Jordan algebras over the four division algebras , , , 𝕆. For more details we refer the reader to . For M-theory compactification on a Calabi-Yau threefold with Hodge numbers $`h_{(1,1)}`$ and $`h_{(2,1)},`$ the five dimensional theory contains the gravity multiplet, $`h_{(1,1)}1`$ vector multiplets and $`h_{(2,1)}+1`$ hypermultiplets. The $`(h_{(1,1)}1)`$-dimensional space of scalar components of the abelian vector supermultiplets coupled to supergravity can be regarded as a hypersurface of a $`h_{(1,1)}`$-dimensional manifold whose coordinates $`X^I(\varphi )`$ are in correspondence with the vector bosons (including the graviphoton). The defining equation of the hypersurface is as in (2.1) $$𝒱(X)=\frac{1}{6}C_{IJK}X^IX^JX^K=X^IX_I=1,I,J,K=1,\mathrm{},h_{(1,1)}.$$ (2.3) Here $`C_{IJK}`$ are the topological intersection numbers of the Calabi-Yau, $`X_I`$ are the so called “dual” special coordinates. The bosonic part of the ungauged supersymmetric $`N=2`$ Lagrangian which describes the coupling of vector multiplets to supergravity is given by $$e^1=\frac{1}{2}R\frac{1}{4}G_{IJ}F_{\mu \nu }{}_{}{}^{I}F_{}^{\mu \nu J}\frac{1}{2}𝒢_{ij}_\mu \varphi ^i^\mu \varphi ^j+\frac{e^1}{48}ϵ^{\mu \nu \rho \sigma \lambda }C_{IJK}F_{\mu \nu }^IF_{\rho \sigma }^JA_\lambda ^K.$$ (2.4) The corresponding vector and scalar metric are completely encoded in the function $`𝒱(X)`$, $`G_{IJ}`$ $`={\displaystyle \frac{1}{2}}_I_J\mathrm{ln}𝒱(X)\|_{𝒱=1}`$ $`𝒢_{ij}`$ $`=G_{IJ}_iX^I_jX^J\|_{𝒱=1}`$ (2.5) where $`_i`$ and $`_I`$ refer, respectively, to partial derivatives with respect to the scalar fields $`\varphi ^i`$ and $`X^I=X^I(\varphi ^i)`$. Further useful relations are $$_iX_I=\frac{2}{3}G_{IJ}_iX^J,X_I=\frac{2}{3}G_{IJ}X^J.$$ (2.6) It is worth pointing out that for Calabi-Yau compactifications, $`𝒱`$ represents the intersection form, $`X^I`$ and $`X_I=\frac{1}{6}C_{IJK}X^JX^K`$ correspond, respectively, to the size of the two- and four-cycles of the Calabi-Yau threefold. ## 3 Black Holes in the $`STU=1`$ Model In the last few years considerable progress has been made in the study of BPS black hole states of the low-energy effective actions of compactified string and $`M`$-theory. This was mainly motivated by the important role that these states play in the understanding of the non-perturbative structure of string theory. The magnetic and electric BPS solutions of five-dimensional $`N=2`$ supergravity models coupled to vector and hypermultiplets can be regarded as solitons interpolating between two vacua: Minkowski flat space at infinity and $`AdS_3\times S^2`$ and $`AdS_2\times S^3`$ near the horizon. At a generic point in space-time, the BPS solution breaks half of supersymmetry. However, near the horizon supersymmetry is enhanced and fully restored. In M-theory compactified on a Calabi-Yau threefold, electrically charged point-like and magnetically charged string-like BPS states correspond to the two and five-branes of M-theory wrapped around the two- and four-cycles of the Calabi-Yau space respectively. Though the details of the low-energy Lagrangian depend very much on the geometric and topological data of the compactified Calabi-Yau space, the analysis of the BPS solutions is considerably simplified by the rich geometric structure based on “very special geometry” underlying the $`N=2`$ five-dimensional theories with vector supermultiplets . The metric for the BPS black hole solutions can be brought to the form $`ds^2`$ $`=`$ $`e^{4V}(dt+w_mdx^m)^2+e^{2V}d\stackrel{}{x}^2,`$ $`F_{mn}^I`$ $`=`$ $`_m(X^I𝒬_n)_n(X^I𝒬_m),`$ $`F_{tm}^I`$ $`=`$ $`_m(e^{2V}X^I),`$ $`e^{2V}X_I`$ $`=`$ $`{\displaystyle \frac{1}{3}}H_I,`$ (3.1) where $`H_I`$ are harmonic functions, $`H_I=h_I+\frac{Q_I}{r^2}`$, $`h_I`$ are constants and $`Q_I`$ denote the electric charges. Furthermore one has $`𝒬_n=e^{2V}w_{n\text{ }}`$, and the field strength of $`w_m`$ is self-dual. If one defines the rescaled coordinates $`Y^I=e^VX^I,`$ then the underlying very special geometry implies that $$e^{3V}=\frac{1}{6}C_{IJK}Y^IY^JY^K.$$ As a general magnetic string solution of $`D=5`$, $`N=2`$ supergravity, one obtains $`ds^2`$ $`=`$ $`e^W(dt^2+dz^2)+e^{2W}d\stackrel{}{x}^2,`$ $`F_{mn}^I`$ $`=`$ $`ϵ_{mnp}_pH^I,\text{ }e^WX^I=H^I,`$ $`e^{3W}`$ $`=`$ $`{\displaystyle \frac{1}{6}}C_{IJK}H^IH^JH^K,`$ (3.2) with the harmonic functions $$H^I=h^I+\frac{P^I}{r},$$ (3.3) where $`h^I`$ are constants and $`P^I`$ are magnetic charges. In the electric case, the near-horizon geometry is given by $`AdS_2\times S^3`$ and the black hole entropy, related to the horizon volume $`S^3`$, is given in terms of the extremized electric central charge. For the magnetically charged $`D=5`$ BPS black string, with near-horizon geometry $`AdS_3\times S^2`$, one similarly finds that the extremized value of the BPS tension is related to the volume of the $`S^2`$. As an example we consider the $`STU=1`$ model . This can be obtained by compactification of heterotic string theory on $`K_3\times S^1`$ . The tree-level prepotential of this model is given by $$𝒱=STU=1,$$ (3.4) and corresponds to the scalar manifold (2.2) for $`n=2`$. Taking $`S=X^0`$, $`T=X^1`$ and $`U=X^2`$, one gets for the matrix $`G^{IJ}`$ $$G^{IJ}=2\left(\begin{array}{ccc}S^2& 0& 0\\ 0& T^2& 0\\ 0& 0& U^2\end{array}\right).$$ (3.5) Considering $`S`$ as the dependent field, i. e. $`S=1/(TU)`$, we find $$𝒢_{ij}=\left(\begin{array}{cc}\frac{1}{T^2}& \frac{1}{2TU}\\ \frac{1}{2TU}& \frac{1}{U^2}\end{array}\right),𝒢^{ij}=\frac{4}{3}\left(\begin{array}{cc}T^2& \frac{TU}{2}\\ \frac{TU}{2}& U^2\end{array}\right).$$ (3.6) The field equations following from the action (2.4) admit the non-extremal static black hole solution $`ds^2`$ $`=`$ $`e^{4V}fdt^2+e^{2V}(f^1dr^2+r^2d\mathrm{\Omega }_3^2),`$ $`F_{rt}^I`$ $`=`$ $`H_I^2_r\stackrel{~}{H}_I,`$ (3.7) $`X^I`$ $`=`$ $`H_I^1e^{2V},`$ where $$d\mathrm{\Omega }_3^2=d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2+\mathrm{cos}^2\theta d\psi ^2$$ (3.8) denotes the metric on the three sphere $`S^3`$. The $`H_I`$ are harmonic functions given by $$H_I=1+\frac{Q_I}{r^2},$$ (3.9) and $`V`$ reads $$e^{2V}=(H_0H_1H_2)^{1/3}.$$ (3.10) Furthermore we have $$f=1\frac{\mu }{r^2}$$ (3.11) with the nonextremality parameter $`\mu `$, and $$\stackrel{~}{H}_I=1+\frac{\stackrel{~}{Q}_I}{r^2},$$ (3.12) where the $`\stackrel{~}{Q}_I`$ denote the physical electric charges. They are related to the $`Q_I`$ appearing in (3.9) by the equations $`Q_I`$ $`=`$ $`{\displaystyle \frac{\mu }{2}}\mathrm{sinh}\beta _I\mathrm{tanh}{\displaystyle \frac{\beta _I}{2}},`$ $`\stackrel{~}{Q}_I`$ $`=`$ $`{\displaystyle \frac{\mu }{2}}\mathrm{sinh}\beta _I.`$ (3.13) The extremal (BPS) limit is reached when $`\beta _I\mathrm{}`$, $`\mu 0`$, with $`\mu \mathrm{sinh}\beta _I`$ kept fixed. For the ADM mass $`M_{ADM}`$, the Bekenstein-Hawking entropy $`S_{BH}`$, and the Hawking temperature $`T_H`$ one obtains $`M_{ADM}`$ $`=`$ $`{\displaystyle \frac{\pi }{4G_5}}({\displaystyle \underset{I}{}}Q_I+{\displaystyle \frac{3}{2}}\mu ),`$ (3.14) $`S_{BH}`$ $`=`$ $`{\displaystyle \frac{A_{hor}}{4G_5}}={\displaystyle \frac{\pi ^2}{2G_5}}{\displaystyle \underset{I}{}}(\mu +Q_I)^{1/2},`$ (3.15) $`T_H`$ $`=`$ $`{\displaystyle \frac{\mu }{\pi _I(\mu +Q_I)^{1/2}}}.`$ (3.16) In the extremal case, also a rotating generalization of (3.7) can be obtained from the general form (3.1). Its metric is given by $$ds^2=e^{4V}(dt+w_\varphi (r,\theta )d\varphi +w_\psi (r,\theta )d\psi )^2+e^{2V}(dr^2+r^2d\mathrm{\Omega }^2),$$ (3.17) where $`w_\varphi (r,\theta )`$ $`=`$ $`{\displaystyle \frac{\alpha \mathrm{sin}^2\theta }{r^2}},`$ $`w_\psi (r,\theta )`$ $`=`$ $`{\displaystyle \frac{\alpha \mathrm{cos}^2\theta }{r^2}}.`$ (3.18) The gauge fields are $$A_t^I=e^{2V}X^I,A_\varphi ^I=e^{2V}X^Iw_\varphi ,A_\psi ^I=e^{2V}X^Iw_\psi .$$ (3.19) The moduli $`X^I`$ and the functions $`V`$ and $`H_I`$ are as in (3.7), (3.10) and (3.9) respectively, and the ADM mass is given by (3.14) for $`\mu =0`$. The Bekenstein-Hawking entropy and the angular momenta read $`S_{BH}`$ $`=`$ $`{\displaystyle \frac{A_{hor}}{4G}}={\displaystyle \frac{\pi ^2}{2G_5}}\sqrt{Q_0Q_1Q_2\alpha ^2},`$ $`J^\varphi `$ $`=`$ $`J^\psi ={\displaystyle \frac{\alpha \pi }{4G_5}}.`$ ## 4 Near-Horizon Limit and Isometry Superalgebra In the following two sections, we shall be particularly interested in the near-horizon limit of (3.17). For $`r0`$ we can write $$e^{2V}=\frac{Z_{hor}}{3r^2},$$ (4.1) where $`Z=Q_IX^I`$ is the central charge, and $`Z_{hor}=3(Q_0Q_1Q_2)^{1/3}`$ is its value at the horizon. Introducing the horospherical coordinates $`(\tau ,\rho )`$, $$\tau =\frac{t}{2}\sqrt{\frac{3}{Z_{hor}}},\rho =2r^2\sqrt{\frac{3}{Z_{hor}}},$$ (4.2) one gets for the near-horizon metric $`ds^2`$ $`=`$ $`\rho ^2d\tau ^2+{\displaystyle \frac{Z_{hor}}{12\rho ^2}}d\rho ^2+{\displaystyle \frac{6\alpha }{Z_{hor}}}\rho d\tau (\mathrm{sin}^2\theta d\varphi \mathrm{cos}^2\theta d\psi )`$ (4.3) $`+{\displaystyle \frac{Z_{hor}}{3}}\left(d\mathrm{\Omega }^2{\displaystyle \frac{27\alpha ^2}{Z_{hor}^3}}(\mathrm{sin}^2\theta d\varphi \mathrm{cos}^2\theta d\psi )^2\right).`$ We observe that, in contrast to the case of vanishing rotation parameter, the spacetime does not split into a product $`AdS_2\times S^3`$. Although the $`AdS_2`$ part is the same as without rotation, there are nondiagonal elements, and the three-sphere is distorted. The isometry superalgebra of the near-horizon supergravity configuration was determined in , where the fact that the residual isometry supergroup can be determined (modulo bosonic factors) from a knowledge of the Killing spinors has been used. In this way, one obtains that the near-horizon geometry is invariant under the superalgebra $`su(1,1\|2)u(1)`$ in the rotating case, and under $`su(1,1\|2)su(2)`$ for $`\alpha =0`$ . Thus, for $`\alpha 0`$, the bosonic subalgebra is $`su(1,1)su(2)_Lu(1)_R`$. In fact, the near-horizon spacetime is a homogeneous manifold of the form $`[SO(2,1)\times SU(2)_L\times U(1)_R]/[U(1)\times U(1)]`$ . The conformal algebra $`su(1,1)so(2,1)`$ is generated by the Killing vectors $`h`$ $`=`$ $`_\tau ,`$ $`d`$ $`=`$ $`\rho _\rho \tau _\tau ,`$ (4.4) $`k`$ $`=`$ $`{\displaystyle \frac{Z_{hor}}{12\rho ^2}}(1{\displaystyle \frac{27\alpha ^2}{Z_{hor}^3}}_\tau )\tau ^2_\tau +2\tau \rho _\rho {\displaystyle \frac{3\alpha }{2Z_{hor}\rho }}(_\varphi _\psi ),`$ satisfying $$[d,h]=h,[d,k]=k,[h,k]=2d.$$ (4.5) Thus, although the manifold is not a product $`AdS_2\times S^3`$, we find the $`so(2,1)`$ symmetry inherent to $`AdS_2`$. In the following section, we will see that this symmetry, which is the conformal symmetry in $`0+1`$ dimensions, occurs also in the action of a particle charged under the vectors moving in the near-horizon regime. ## 5 Particle Motion near the Horizon We now consider a particle of mass $`m`$, carrying the charges $`q_I`$ under the abelian vectors, which moves in the background (4.3). Like in , we introduce the new coordinate $`q`$, $$\rho =\frac{Z_{hor}}{3q^2}.$$ (5.1) We use a Hamiltonian formalism, and define $$=\frac{1}{2}g^{\mu \nu }(\mathrm{\Pi }_\mu q_IA_\mu ^I)(\mathrm{\Pi }_\nu q_IA_\nu ^I),$$ (5.2) where the $`\mathrm{\Pi }_\mu `$ denote generalized momenta. For our configuration, this leads to $``$ $`=`$ $`{\displaystyle \frac{9q^4}{2Z_{hor}^2}}\mathrm{\Pi }_\tau ^2\left(1{\displaystyle \frac{27\alpha ^2}{Z_{hor}^3}}\right)+{\displaystyle \frac{27q^2\alpha }{Z_{hor}^3}}\mathrm{\Pi }_\tau \left[\mathrm{\Pi }_\varphi \mathrm{\Pi }_\psi +{\displaystyle \frac{Z_{hor}^2}{9\alpha }}q_IX_{hor}^I\right]`$ (5.3) $`{\displaystyle \frac{1}{2}}(q_IX_{hor}^I)^2+{\displaystyle \frac{3q^2}{2Z_{hor}}}\mathrm{\Pi }_q^2+{\displaystyle \frac{3L^2}{2Z_{hor}}},`$ where $$L^2=\mathrm{\Pi }_\theta ^2+\frac{\mathrm{\Pi }_\varphi ^2}{\mathrm{sin}^2\theta }+\frac{\mathrm{\Pi }_\psi ^2}{\mathrm{cos}^2\theta }$$ (5.4) denotes the conserved angular momentum. As the coordinates $`\tau `$, $`\varphi `$ and $`\psi `$ are cyclic, the associated conjugate momenta are constants of motion. If $``$ solves the mass-shell constraint $`2=m^2`$, $`\mathrm{\Pi }_\tau `$ is to be identified with the particle Hamiltonian $`H`$. Setting $`\mathrm{\Pi }_q=p`$ and defining $`u=pq`$, one obtains $$H=\frac{p^2}{2F(u)}+\frac{mg}{2q^2F(u)},$$ (5.5) with $$mg=L^2+\frac{Z_{hor}}{3}(m^2(q_IX_{hor}^I)^2),$$ (5.6) and the function $`F(u)`$ given by $$F(u)=\frac{1}{2}\left[\sqrt{C^2+(1\frac{27\alpha ^2}{Z_{hor}^3})(\frac{3}{Z_{hor}}(u^2+L^2)+m^2(q_IX_{hor}^I)^2)}+C\right].$$ (5.7) In (5.7), the constant $`C`$ is defined by $$C=\frac{9\alpha }{Z_{hor}^2}(\mathrm{\Pi }_\varphi \mathrm{\Pi }_\psi )+q_IX_{hor}^I.$$ (5.8) One observes that in the limit $$Z_{hor}\mathrm{},mq_IX_{hor}^I0,$$ (5.9) with $`Z_{hor}(mq_IX_{hor}^I)`$ kept fixed, we have $`F(u)m`$, and (5.5) reduces to the DFF model $$H=\frac{p^2}{2m}+\frac{g}{2q^2}.$$ (5.10) Note that also the general Hamiltonian (5.5) describes a model of conformal mechanics. To see this, we write it in the form $$H=\frac{p^2}{2f(u)},$$ (5.11) with $$f(u)=\frac{u^2F(u)}{u^2+mg}.$$ (5.12) The generators of the conformal group are then given by $$H=\frac{p^2}{2f},D=\frac{1}{2}u,K=\frac{1}{2}q^2f,$$ (5.13) satisfying the Poisson bracket algebra $$[D,H]=H,[D,K]=K,[H,K]=2D.$$ (5.14) ## 6 Statistical Entropy from $`AdS_2`$ Central Charge In order to determine the central charge of the boundary CFT, we proceed along the lines of , and reduce the bosonic part of the $`D=5`$, $`N=2`$ supergravity action to two dimensions. In this section we shall only consider nonrotating black holes carrying electric charge. This means that we can consistently truncate the Chern-Simons term in (2.4), so that the bosonic part of the action in five dimensions reads $$I=\frac{1}{16\pi G_5}d^5x\sqrt{g}\left[R\frac{1}{2}G_{IJ}F_{\mu \nu }^IF^{J\mu \nu }𝒢_{ij}_\mu \varphi ^i^\mu \varphi ^j\right],$$ (6.1) where $`G_5=l^3`$ denotes Newton’s constant. The matrices $`G_{IJ}`$ and $`𝒢_{ij}`$ for the $`STU`$ model were given in section 3. The reduction ansatz for the metric is $$ds^2=ds_{(2)}^2+l^2\mathrm{\Phi }^2d\mathrm{\Omega }^2,$$ (6.2) where $`\mathrm{\Phi }`$ denotes the dilaton and $`d\mathrm{\Omega }^2`$ is given by (3.8). One now assumes that the gauge fields, scalars and dilaton do not depend on the coordinates on the internal $`S^3`$. In this way, one arrives at the two-dimensional effective action $$I=\frac{\mathrm{\Omega }}{16\pi }d^2x\sqrt{g}\left[\mathrm{\Phi }^3R+6\mathrm{\Phi }(\mathrm{\Phi })^2+\frac{6\mathrm{\Phi }}{l^2}\mathrm{\Phi }^3𝒢_{ij}_\alpha \varphi ^i^\alpha \varphi ^j\frac{1}{2}\mathrm{\Phi }^3G_{IJ}F_{\alpha \beta }^IF^{J\alpha \beta }\right],$$ (6.3) where $`\mathrm{\Omega }=2\pi ^2`$ denotes the volume of the unit $`S^3`$, and early greek indices $`\alpha ,\beta ,\mathrm{}`$ refer to two-dimensional spacetime. We now wish to integrate out the field strength $`F_{\alpha \beta }^I`$ (which in two dimensions must be a multiple of the volume form $`ϵ_{\alpha \beta }`$) from the action. This can be done using the Lagrange multiplier method of . Let us briefly sketch how this works: Instead of looking at the gauge field action $$I_g=\frac{\mathrm{\Omega }}{16\pi }d^2x\sqrt{g}\frac{1}{2}\mathrm{\Phi }^3G_{IJ}F_{\alpha \beta }^IF^{J\alpha \beta },$$ (6.4) one looks at the formally extended action $$\stackrel{~}{I}_g=\frac{\mathrm{\Omega }}{16\pi }d^2x\left[\frac{1}{2}\sqrt{g}\mathrm{\Phi }^3G_{IJ}F_{\alpha \beta }^IF^{J\alpha \beta }+\lambda _I(F_{\alpha \beta }^I_\alpha A_\beta ^I+_\beta A_\alpha ^I)ϵ^{\alpha \beta }\right],$$ (6.5) where the definition of $`F^I`$ as a field strength associated with $`A^I`$ is implemented by means of the Lagrange multiplier $`\lambda _I`$. Note that the three variables $`F^I,A^I`$ and $`\lambda _I`$ are considered as independent in this setting. Variation with respect to $`A^I`$ yields $$_\beta (\lambda _Iϵ^{\alpha \beta })=0,$$ (6.6) so that $`\mathrm{\Lambda }_I:=\lambda _I/\sqrt{g}`$ are constants. Due to (6.6), the term $`\lambda _I(_\alpha A_\beta ^I+_\beta A_\alpha ^I)ϵ^{\alpha \beta }`$ in the action (6.5) is a boundary term and can be dropped. We can then integrate out the field strength $`F^I`$, using its equation of motion $$\mathrm{\Phi }^3G_{IJ}F^{J\alpha \beta }+\mathrm{\Lambda }_Iϵ^{\alpha \beta }=0.$$ (6.7) This yields $$\stackrel{~}{I}_g=\frac{\mathrm{\Omega }}{16\pi }d^2x\sqrt{g}\frac{G^{IJ}\mathrm{\Lambda }_I\mathrm{\Lambda }_J}{\mathrm{\Phi }^3},$$ (6.8) so that we are left with the total two-dimensional action $$I=\frac{\mathrm{\Omega }}{16\pi }d^2x\sqrt{g}\left[\mathrm{\Phi }^3R+6\mathrm{\Phi }(\mathrm{\Phi })^2+\frac{6\mathrm{\Phi }}{l^2}\mathrm{\Phi }^3𝒢_{ij}_\alpha \varphi ^i^\alpha \varphi ^j\frac{G^{IJ}\mathrm{\Lambda }_I\mathrm{\Lambda }_J}{\mathrm{\Phi }^3}\right].$$ (6.9) The dilaton kinetic term in (6.9) can be eliminated by a conformal rescaling $$\overline{g}_{\alpha \beta }=\mathrm{\Phi }^2g_{\alpha \beta }.$$ (6.10) Defining $`\overline{\mathrm{\Phi }}=\mathrm{\Phi }^3`$, we obtain $$I=\frac{\mathrm{\Omega }}{16\pi }d^2x\sqrt{\overline{g}}\left[\overline{\mathrm{\Phi }}\overline{R}+\frac{6}{l^2\overline{\mathrm{\Phi }}^{1/3}}\overline{\mathrm{\Phi }}𝒢_{ij}_\alpha \varphi ^i^\alpha \varphi ^j\frac{G^{IJ}\mathrm{\Lambda }_I\mathrm{\Lambda }_J}{\overline{\mathrm{\Phi }}^{5/3}}\right].$$ (6.11) Let us now consider the nonextremal black hole solution (3.7) of the action (6.1), and expand it near extremality. To this end, we introduce an expansion parameter $`ϵ`$ ($`ϵ0`$), and set $`t`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{t}}{ϵ}},r=\sqrt{ϵ}\stackrel{~}{r},\mu =\mu _0ϵ,`$ $`\overline{\mathrm{\Phi }}`$ $`=`$ $`\overline{\mathrm{\Phi }}_0+ϵ\phi ,\varphi ^i=\varphi _0^i+ϵ\stackrel{~}{\varphi }^i,`$ (6.12) where $$\overline{\mathrm{\Phi }}_0=\frac{(\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2)^{1/2}}{l^3},\varphi _0^i=\stackrel{~}{Q}_i^1(\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2)^{1/3}.$$ (6.13) Introducing the new coordinate $$x=\frac{(\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2)^{1/6}}{2l^2}(\stackrel{~}{r}^2\frac{\mu _0}{2}),$$ (6.14) we arrive at $$d\overline{s}^2=(\lambda ^2x^2a^2)d\stackrel{~}{t}^2+(\lambda ^2x^2a^2)^1dx^2$$ (6.15) for the rescaled two-dimensional metric, with $`\lambda `$ and $`a`$ given by $`\lambda `$ $`=`$ $`{\displaystyle \frac{2l}{(\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2)^{1/3}}},`$ (6.16) $`a^2`$ $`=`$ $`{\displaystyle \frac{\mu _0^2}{4l^2(\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2)^{1/3}}}.`$ (6.17) Defining a new dilaton $`\eta `$ by $$\eta =\frac{\mathrm{\Omega }ϵ\phi }{8\pi },$$ (6.18) we obtain for the action at lowest order in the expansion parameter $`ϵ`$, $$I=\frac{1}{2}d^2x\sqrt{\overline{g}}\eta [\overline{R}+2\lambda ^2],$$ (6.19) so the leading order is governed by the Jackiw-Teitelboim (JT) model . (6.15), together with the linear dilaton $`\eta `$ $`=`$ $`\eta _0\lambda x,`$ (6.20) $`\eta _0`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }ϵ}{16\pi l^2}}(\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2)^{2/3}{\displaystyle \underset{I}{}}\stackrel{~}{Q}_I^1,`$ represents a black hole solution of this model , with mass and thermodynamical entropy given by $`M_{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\eta _0a^2\lambda ,`$ $`S_{(2)}`$ $`=`$ $`2\pi \eta _{hor}=2\pi \eta _0a.`$ (6.21) This black hole spacetime has constant curvature, i. e. it is locally $`AdS_2`$. Now it is known that the asymptotic symmetries of two-dimensional anti-de Sitter space form a Virasoro algebra , similar to the case of $`AdS_3`$, where one has two copies of Virasoro algebras as asymptotic symmetries . When realized canonically in the Hamiltonian formulation of JT gravity, this algebra was shown to exhibit a central charge $$c=24\eta _0.$$ (6.22) Using this central charge in Cardy’s formula, the authors of were able to give a microscopic derivation of the entropy of the two-dimensional black holes (6.15) in the JT model. Our aim is now to perform a similar calculation for the near-extremal five dimensional black hole under consideration, making use of the fact that the dimensionally reduced supergravity action coincides with the JT model at leading order in the nonextremality parameter, and that the relevant two-dimensional metric is given by (6.15)<sup>2</sup><sup>2</sup>2Cf. for similar computations in the case of heterotic 4D string black holes.. First of all, we expand the ADM mass $`M_{ADM}`$ (3.14) and Bekenstein-Hawking entropy $`S_{BH}`$ (3.15) of the black hole (3.7) in five dimensions for $`\mu 0`$, yielding $`M_{ADM}`$ $`=`$ $`{\displaystyle \frac{\pi }{4l^3}}{\displaystyle \underset{I}{}}\stackrel{~}{Q}_I(1+{\displaystyle \frac{\mu ^2}{8\stackrel{~}{Q}_I^2}}),`$ (6.23) $`S_{BH}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }}{4l^3}}(\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2)^{1/2}(1+{\displaystyle \frac{\mu }{4}}{\displaystyle \underset{I}{}}\stackrel{~}{Q}_I^1),`$ (6.24) so the small excitations above extremality have the energy $$\mathrm{\Delta }M_{ADM}=\frac{\pi \mu ^2}{32l^3}\underset{I}{}\stackrel{~}{Q}_I^1$$ (6.25) and entropy $$\mathrm{\Delta }S_{BH}=\frac{\mathrm{\Omega }\mu }{16l^3}(\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2)^{1/2}\underset{I}{}\stackrel{~}{Q}_I^1.$$ (6.26) Comparing this with the two-dimensional results (6.21), one finds $`\mathrm{\Delta }S_{BH}=S_{(2)}`$ and $`\mathrm{\Delta }M_{ADM}=ϵM_{(2)}`$. The factor $`ϵ`$ appearing in the relation between the two masses stems from the fact that $`M_{ADM}`$ was computed with respect to the Killing vector $`_t`$, whereas $`M_{(2)}`$ is related to $`_{\stackrel{~}{t}}=ϵ_t`$. This means that up to these normalizations the five- and two-dimensional energies and entropies match. Expanding also the Hawking temperature (3.16) for small values of the nonextremality parameter $`\mu `$, one finds for the temperature dependence of $`\mathrm{\Delta }M_{ADM}`$ $$\mathrm{\Delta }M_{ADM}=\frac{\pi ^3T_H^2}{32l^3}\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2\underset{I}{}\stackrel{~}{Q}_I^1,$$ (6.27) so the energy of the excitations above extremality is that of an ideal gas of massless particles in $`1+1`$ dimensions. This suggests that the microstates should be described by a two-dimensional field theory rather than a quantum mechanical system. Let us now proceed with the computation of the statistical entropy, using the central charge (6.22). The Virasoro generator $`L_0`$ for the black hole (6.15) is given by $$L_0=\frac{M_{(2)}}{\lambda }=\frac{\mathrm{\Omega }ϵ\mu _0^2}{128\pi l^4}(\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2)^{1/3}\underset{I}{}\stackrel{~}{Q}_I^1.$$ (6.28) Inserting this together with the central charge (6.22) into Cardy’s formula, we get for the statistical entropy $$S_{stat}=2\pi \sqrt{\frac{cL_0}{6}}=\frac{\mathrm{\Omega }\mu }{8\sqrt{2}l^3}(\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2)^{1/2}\underset{I}{}\stackrel{~}{Q}_I^1,$$ (6.29) which agrees, up to a factor $`\sqrt{2}`$, with the thermodynamical entropy $`\mathrm{\Delta }S_{BH}`$ of the small excitations above extremality. The same mismatch by a factor $`\sqrt{2}`$ has been found in , where the authors proposed an explanation of this for the case when the model (6.19) comes from dimensional reduction of three-dimensional $`AdS`$ gravity. Although in our case $`AdS_2`$ arises as near-horizon geometry of a higher-dimensional black hole with no intermediate $`AdS_3`$ geometry involved, we shall see in the next section that by means of a duality transformation the near-horizon geometry $`AdS_2\times S^3`$ of the extremal black hole becomes $`AdS_3\times S^2`$. We will then be able to use Strominger’s counting of microstates in order to reproduce correctly the Bekenstein-Hawking entropy of the black hole. ## 7 Duality Invariance of the Supergravity Action In this section we will show that in presence of a Killing vector field $`_z`$, the supergravity action (6.1) is invariant under a certain generalization of T-duality<sup>3</sup><sup>3</sup>3By considering (6.1) we assumed that the Chern-Simons term does not contribute. One can easily generalize the discussion below to nonvanishing CS term. This results in a $`\theta `$ term in four dimensions, which does not spoil the considered duality invariance.. The key observation is then that the three sphere $`S^3`$ appearing in the black hole geometry can be written as a Hopf fibration, i. e. as an $`S^1`$ bundle over $`\text{}P^1S^2`$. Performing then a duality transformation along the Hopf fibre untwists the $`S^3`$, and transforms the electrically charged black hole into a magnetically charged black string, which has $`AdS_3\times S^2`$ as near-horizon limit in the extremal case. To begin with, we reduce the action (6.1) to four dimensions, using the usual Kaluza-Klein reduction ansatz for the five-dimensional metric, $$ds^2=e^{k/\sqrt{3}}ds_4^2+e^{2k/\sqrt{3}}(dz+𝒜_\alpha dx^\alpha )^2,$$ (7.1) where $`k`$ denotes the dilaton, and early greek indices $`\alpha ,\beta ,\mathrm{}`$ refer to four-dimensional spacetime. Assuming that the fields appearing in (6.1) are independent of $`z`$, one arrives at the four-dimensional action $$I_4=\frac{L}{16\pi G_5}d^4x\sqrt{g_4}\left[R_4\frac{1}{2}(k)^2\frac{1}{4}e^{\sqrt{3}k}^2\frac{1}{2}e^{k/\sqrt{3}}F^2𝒢_{ij}_\alpha \varphi ^i^\alpha \varphi ^j\right],$$ (7.2) where $`L`$ denotes the length of the circle parametrized by $`z`$, $``$ is the field strength associated to the Kaluza-Klein vector potential $`𝒜`$, and $$^2=_{\alpha \beta }^{\alpha \beta },F^2=G_{IJ}F_{\alpha \beta }^IF^{J\alpha \beta }.$$ (7.3) We now dualize both $``$ and $`F^I`$, using again the Lagrange multiplier method of . Dropping boundary terms, we arrive at the dualized action $`I_4`$ $`=`$ $`{\displaystyle \frac{L}{16\pi G_5}}{\displaystyle }d^4x\sqrt{g_4}[R_4{\displaystyle \frac{1}{2}}(k)^2{\displaystyle \frac{1}{4}}e^{\sqrt{3}k}(^{})^2`$ (7.4) $`{\displaystyle \frac{1}{2}}e^{k/\sqrt{3}}{\displaystyle \frac{1}{4}}G^{IJ}{}_{}{}^{}F_{I\alpha \beta }^{}F_J^{\alpha \beta }𝒢_{ij}_\alpha \varphi ^i^\alpha \varphi ^j],`$ where we defined $`{}_{}{}^{}_{\alpha \beta }^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\sqrt{3}k}ϵ_{\alpha \beta \gamma \delta }^{\gamma \delta },`$ (7.5) $`{}_{}{}^{}F_{I\alpha \beta }^{}`$ $`=`$ $`e^{k/\sqrt{3}}G_{IJ}ϵ_{\alpha \beta \gamma \delta }F^{J\gamma \delta }.`$ (7.6) Comparing (7.4) with (7.2), we observe that the gravitational and gauge field parts of the four-dimensional action, as well as the dilaton kinetic energy, are invariant under the $`\text{}_4`$ transformation $$kk,_{\alpha \beta }^{}_{\alpha \beta },F_{\alpha \beta }^I^{}F_{I\alpha \beta },G_{IJ}\frac{1}{4}G^{IJ}.$$ (7.7) The $`\text{}_4`$ is actually a subgroup of the usual symplectic $`Sp(2m+2,\text{})`$ duality group of $`D=4`$, $`N=2`$ supergravity (coupled to $`m`$ vector multiplets) generated by $$S=\left(\begin{array}{cc}\mathrm{𝟎}& \mathrm{𝟏}\\ \mathrm{𝟏}& \mathrm{𝟎}\end{array}\right).$$ (7.8) Note that the transformation $`G_{IJ}G^{IJ}/4`$ means that $`X^I`$ $``$ $`3X_I={\displaystyle \frac{1}{2}}C_{IJK}X^JX^K,`$ $`X_I`$ $``$ $`{\displaystyle \frac{1}{3}}X^I,`$ (7.9) so essentially the special coordinates go over into their duals. The fact that this dualization implies $`G_{IJ}G^{IJ}/4`$ can be shown using the expression $$G_{IJ}=\frac{9}{2}X_IX_J\frac{1}{2}C_{IJK}X^K,$$ (7.10) as well as the ”adjoint identity” $$C_{IJK}C_{J^{}(LM}C_{PQ)K^{}}\delta ^{JJ^{}}\delta ^{KK^{}}=\frac{4}{3}\delta _{I(L}C_{MPQ)}$$ (7.11) of the associated Jordan algebra . It can also be seen that this duality transformation is consistent with the relations (2.6). Furthermore, making use of the equation $$𝒢_{ij}_\alpha \varphi ^i^\alpha \varphi ^j=G_{IJ}_\alpha X^I^\alpha X^J,$$ (7.12) one checks that (7.9) does not change the kinetic term of the scalar fields, so (7.7), (7.9) represent in fact a duality invariance of the four-dimensional action (7.2). In the special case of the $`STU=1`$ model, (7.9) implies that the moduli $`\varphi ^i`$ go over into their inverse, $$\varphi ^i\frac{1}{\varphi ^i}.$$ (7.13) We now wish to apply the duality (7.7), (7.9) to the black hole solution (3.7). To this end, we consider the $`S^3`$ as an $`S^1`$ bundle over $`S^2`$, and write for its metric $$d\mathrm{\Omega }^2=\frac{1}{4}\left[d\vartheta ^2+\mathrm{sin}^2\vartheta d\phi ^2+(d\zeta +\mathrm{cos}\vartheta d\phi )^2\right],$$ (7.14) where $`\zeta `$ ($`0\zeta 4\pi `$) parametrizes the $`S^1`$ fibre. Introducing the coordinate $`z=\lambda \zeta `$, where $`\lambda `$ denotes an arbitrary length scale, one can write the 5d metric in the KK form (7.1), where $`ds_4^2`$ $`=`$ $`{\displaystyle \frac{re^V}{2\lambda }}\left[e^{4V}fdt^2+e^{2V}f^1dr^2+e^{2V}{\displaystyle \frac{r^2}{4}}(d\vartheta ^2+\mathrm{sin}^2\vartheta d\phi ^2)\right],`$ $`e^{k/\sqrt{3}}`$ $`=`$ $`{\displaystyle \frac{re^V}{2\lambda }},`$ (7.15) $`𝒜`$ $`=`$ $`\lambda \mathrm{cos}\vartheta d\phi .`$ (Note that $`=d𝒜`$ is essentially the Kähler form on $`S^2`$). We now dualize in 4d according to (7.7), and then relift the solution to five dimensions. This yields the configuration $`ds^2`$ $`=`$ $`e^{2V}\left[{\displaystyle \frac{\mu }{4\lambda ^2}}dt^2+2dzdt+{\displaystyle \frac{4\lambda ^2}{r^2}}dz^2\right]+{\displaystyle \frac{r^2}{4\lambda ^2}}e^{4V}\left[f^1dr^2+{\displaystyle \frac{r^2}{4}}d\mathrm{\Omega }_2^2\right],`$ $`F_{\vartheta \phi }^I`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{Q}_I}{4\lambda }}\mathrm{sin}\vartheta ,`$ (7.16) $`X^I`$ $`=`$ $`H_Ie^{2V}.`$ One effect of the duality transformation is thus the untwisting of the Hopf fibration<sup>4</sup><sup>4</sup>4The fact that Hopf bundles can be untwisted by T-dualities was observed in . The idea of untwisting and twisting fibres to relate strings and black holes, and thus to gain new insights into black hole microscopics, was also explored in .. Although the metric in (7.16) contains nondiagonal elements proportional to $`dzdt`$, there is no rotation present. To see this, one observes that the nondiagonal elements come from the vector potential $`𝒜`$ in four dimensions, which gives rise to the field strength $``$. The equations of motion for $``$ following from the action (7.2) read $$_\alpha (e^{k\sqrt{3}}^{\alpha \beta })=0,$$ (7.17) so there exists an associated conserved charge $$J=_{S_{\mathrm{}}^2}d^2S_{\alpha \beta }e^{k\sqrt{3}}^{\alpha \beta }.$$ (7.18) For the solution (7.16) under consideration, however, one easily verifies that $`J`$ (which, up to a normalization factor, represents the angular momentum) vanishes. One can further simplify (7.16) by an $`SL(2,\text{})`$ transformation $$\left(\begin{array}{c}t^{}\\ z^{}\end{array}\right)=\left(\begin{array}{cc}0& \frac{2\lambda }{\sqrt{\mu }}\\ \frac{\sqrt{\mu }}{2\lambda }& \frac{2\lambda }{\sqrt{\mu }}\end{array}\right)\left(\begin{array}{c}t\\ z\end{array}\right).$$ (7.19) Introducing also the new radial coordinate $`\rho =r^2/(4\lambda )`$, we then get for the metric $$ds^2=e^{2V}(fdt^{}{}_{}{}^{2}+dz^{}{}_{}{}^{2})+e^{4V}(f^1d\rho ^2+\rho ^2d\mathrm{\Omega }_2^2).$$ (7.20) (7.20), together with the gauge and scalar fields given in (7.16), represents a nonextremal generalization of the supersymmetric magnetic black string found in . The duality (7.7) thus maps electrically charged black holes onto magnetically charged black strings. Now a short comment on the $`SL(2,\text{})`$ transformation (7.19) is in order. The orbits of the Killing vector $$_z^{}=\frac{2\lambda }{\sqrt{\mu }}_t$$ (7.21) are non-compact since the time coordinate is non-compact. This means that globally the spacetimes in (7.16) and (7.20) are not equivalent. To make the transformation (7.19) a symmetry, we have to compactify the orbits of $`_z^{}`$. We shall see below however, that the temperature and entropy of one black string can be deduced from the other, which indicates that the two solutions (7.16) and (7.20) are in the same universality class . The Bekenstein-Hawking entropy of the black string (7.20) results to coincide precisely with that of the dual black hole given by (3.15), if we assign to $`z^{}`$ the period $`\mathrm{\Delta }z2\lambda /\sqrt{\mu }`$, where $`\mathrm{\Delta }z=4\pi \lambda `$ denotes the period of $`z`$. The Hawking temperature can be computed by requiring the absence of conical singularities in the Euclidean metric, yielding $$T_H=\frac{2\lambda \sqrt{\mu }}{\pi _I(\mu +Q_I)^{1/2}},$$ (7.22) i. e. $`2\lambda /\sqrt{\mu }`$ times the black hole temperature (3.16). The factor $`2\lambda /\sqrt{\mu }`$ stems from the rescaling of the time coordinate contained in (7.19). Thus, up to this normalization, the temperature and entropy of the black string (7.20) coincide with that of the dual black hole (3.7), i. e. they are duality invariant. ## 8 Microstate Counting from $`AdS_3`$ Gravity We now want to use the near-horizon geometry of the dual solution (7.20) to count the microstates giving rise to the Bekenstein-Hawking entropy. In it was shown that in the extremal case, the geometry becomes $`AdS_3\times S^2`$ near the event horizon. The idea is now to use the central charge of $`AdS_3`$ gravity in Cardy’s formula, in order to compute the statistical entropy, like it was done by Strominger for the BTZ black hole<sup>5</sup><sup>5</sup>5Cf. also , where similar computations for black strings in six dimensions with $`BTZ\times S^3`$ near-horizon geometry were performed.. As only the $`AdS_3`$ part is relevant, we would like to reduce the supergravity action from five to three dimensions. To this end, we first Hodge-dualize the magnetic two-form field strength in (7.16). This yields for the action (6.1) $$I=\frac{1}{16\pi G_5}d^5x\sqrt{g}\left[R\frac{1}{2}G^{IJ}H_{I\mu \nu \rho }H_J^{\mu \nu \rho }𝒢_{ij}_\mu \varphi ^i^\mu \varphi ^j\right],$$ (8.1) where $$H_{I\mu \nu \rho }=\frac{1}{2\sqrt{3}}G_{IJ}ϵ_{\mu \nu \rho \lambda \sigma }F^{J\lambda \sigma }.$$ (8.2) Note that for the solution under consideration, the $`H_I`$ do not depend on the coordinates of the internal $`S^2`$. Furthermore, in 3d the three-forms $`H_I`$ are proportional to the volume form and can be integrated out. For the metric, we use the reduction ansatz $$ds^2=ds_3^2+l^2\mathrm{\Phi }^2d\mathrm{\Omega }_2^2,$$ (8.3) where $`G_5=l^3`$ as before, and $`d\mathrm{\Omega }_2^2`$ denotes the standard metric on the unit $`S^2`$. This gives the reduced action $$I=\frac{1}{4l}d^3x\sqrt{g_3}\mathrm{\Phi }^2\left[R_3+\frac{2}{l^2\mathrm{\Phi }^2}+\frac{2}{\mathrm{\Phi }^2}(\mathrm{\Phi })^2\frac{1}{2}G^{IJ}H_{I\alpha \beta \gamma }H_J^{\alpha \beta \gamma }𝒢_{ij}_\alpha \varphi ^i^\alpha \varphi ^j\right],$$ (8.4) where early greek indices $`\alpha ,\beta ,\mathrm{}`$ refer to three-dimensional spacetime. Using the procedure described in section 6, the three-forms $`H_I`$ can be integrated out. In this way, one finally obtains $$I=\frac{1}{4l}d^3x\sqrt{g_3}\mathrm{\Phi }^2\left[R_3+\frac{2}{l^2\mathrm{\Phi }^2}+\frac{2}{\mathrm{\Phi }^2}(\mathrm{\Phi })^2\frac{G_{IJ}P^IP^J}{\mathrm{\Phi }^4l^4}𝒢_{ij}_\alpha \varphi ^i^\alpha \varphi ^j\right],$$ (8.5) where we introduced the magnetic charges $$P^I=\frac{\stackrel{~}{Q}_I}{4\lambda }$$ (8.6) of the black string (7.16). We find it convenient to conformally rescale the metric, $$\overline{g}_{\alpha \beta }=\mathrm{\Phi }g_{\alpha \beta },$$ (8.7) yielding $$I=\frac{1}{4l}d^3x\sqrt{\overline{g}_3}\mathrm{\Phi }^{\frac{3}{2}}\left[\overline{R}_3+\frac{2}{l^2\mathrm{\Phi }^3}\frac{3}{2\mathrm{\Phi }^2}(\mathrm{\Phi })^2\frac{G_{IJ}P^IP^J}{\mathrm{\Phi }^5l^4}𝒢_{ij}_\alpha \varphi ^i^\alpha \varphi ^j\right]$$ (8.8) for the action. The conformally rescaled 3d metric reads $$d\overline{s}_3^2=\frac{\rho }{l}(fdt^{}{}_{}{}^{2}+dz^{}{}_{}{}^{2})+e^{6V}\frac{\rho d\rho ^2}{lf}.$$ (8.9) The idea is now to expand this metric near the horizon and near extremality. This can be done by setting $$t^{}=\frac{t^{\prime \prime }}{\sqrt{ϵ}}(2\lambda )^4\sqrt{\frac{l}{\mu _0\lambda \stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2}},z^{}=\frac{z^{\prime \prime }}{\sqrt{ϵ}}\frac{(2\lambda )^2}{\sqrt{\mu _0}},\rho =ϵ\stackrel{~}{r}^2\frac{\mu _0l}{(2\lambda )^4},\mu =\mu _0ϵ,$$ (8.10) and taking the limit $`ϵ0`$. This leads to the metric $$d\overline{s}_3^2=\frac{\stackrel{~}{r}^2\stackrel{~}{r}_+^2}{l_{eff}^2}dt^{\prime \prime }{}_{}{}^{2}+\stackrel{~}{r}^2dz^{\prime \prime }{}_{}{}^{2}+\frac{l_{eff}^2d\stackrel{~}{r}^2}{\stackrel{~}{r}^2\stackrel{~}{r}_+^2},$$ (8.11) where we introduced $`\stackrel{~}{r}_+^2`$ $`=`$ $`{\displaystyle \frac{4\lambda ^3}{l}},`$ $`l_{eff}^2`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2}{16l\lambda ^3}}.`$ (8.12) We recognize (8.11) as the BTZ black hole , with event horizon at $`\stackrel{~}{r}=\stackrel{~}{r}_+`$. One easily verifies that the period of the coordinate $`z^{\prime \prime }`$ is $`2\pi `$. $`\mathrm{\Lambda }_{eff}=1/l_{eff}^2`$ is the effective cosmological constant. The effective 3d Newton constant can be read off from the action (8.8), yielding $$\frac{1}{16\pi G_{eff}}=\frac{1}{4l}\mathrm{\Phi }_{hor}^{3/2},$$ (8.13) where the subscript indicates that the dilaton $`\mathrm{\Phi }`$ is to be evaluated at the horizon. In this way, we get $$\frac{1}{G_{eff}}=\frac{4\pi }{l^{5/2}}e^{3V_{hor}}\rho _{hor}^{3/2}.$$ (8.14) The Bekenstein-Hawking entropy of the BTZ black hole (8.11) is given by $$S_{(3)}=\frac{A_{hor}}{4G_{eff}}=\frac{\pi ^2}{2l^3}\underset{I}{}(\mu +Q_I)^{1/2},$$ (8.15) which, as it should be, equals the entropy (3.15) of the five-dimensional black hole we started with. The BTZ black hole mass $`M_{(3)}`$ can be computed using the formula $$\stackrel{~}{r}_+^2=8G_{eff}M_{(3)}l_{eff}^2,$$ (8.16) which yields $$M_{(3)}=\frac{\lambda ^3}{2lG_{eff}l_{eff}^2}.$$ (8.17) We can now apply Strominger’s counting of microstates to reproduce the Bekenstein-Hawking entropy. To this end, one first observes that the central charge appearing in the asymptotic symmetry algebra of $`AdS_3`$ in our case reads $$c=\frac{3l_{eff}}{2G_{eff}}.$$ (8.18) Furthermore, we have the relations $`M_{(3)}`$ $`=`$ $`{\displaystyle \frac{1}{l_{eff}}}(L_0+\overline{L}_0),`$ (8.19) $`J`$ $`=`$ $`L_0\overline{L}_0`$ (8.20) for the mass and angular momentum. For (8.11) one has $`J=0`$, so $`L_0=\overline{L}_0=\frac{1}{2}l_{eff}M_{(3)}`$. Plugging this, together with the central charge (8.18), into Cardy’s formula $$S_{stat}=2\pi \sqrt{\frac{cL_0}{6}}+2\pi \sqrt{\frac{c\overline{L}_0}{6}}$$ (8.21) yields the statistical entropy $$S_{stat}=\frac{\pi \lambda ^{3/2}}{l^{1/2}G_{eff}}=\frac{\pi ^2}{2l^3}\underset{I}{}(\mu +Q_I)^{1/2},$$ (8.22) which coincides precisely with the thermodynamical entropy (3.15) of the 5d black hole (3.7). ## 9 Final Remarks The conclusions we have drawn are valid for general black holes of $`D=5`$, $`N=2`$ supergravities. In particular they apply also to the case of theories obtained from compactifications on Calabi-Yau spaces. In different contexts there has been a discussion of dualities which connect various black hole solutions. We have exhibited an explicit duality transformation which is an invariance of the action: it turns the $`AdS_2\times S^3`$ near horizon geometry into $`AdS_3\times S^2`$. Our calculation shows that the correct statistical entropy is given by the counting of microstates from $`AdS_3`$, where both $`L_0`$ and $`\overline{L}_0`$ are different from zero. Using instead the central charge of the $`AdS_2`$ Virasoro algebra, with only right-movers, gives a factor $`\sqrt{2}`$ mismatch between statistical and thermodynamical entropy. Within the $`AdS_2`$ approach we were able (up to the mentioned factor $`\sqrt{2}`$) to capture only the entropy of the small excitations above extremality, not that of the ground state itself. The reason for this was the fact that in two dimensions the Einstein-Hilbert term is a topological invariant, and does not contribute to the central charge computed in . In the extremal limit $`ϵ0`$ the $`AdS_2`$ central charge (6.22) vanishes, whereas the central charge (8.18) for $`AdS_3`$ is given by $$c=\frac{3\pi }{16l^3\lambda ^3}\stackrel{~}{Q}_0\stackrel{~}{Q}_1\stackrel{~}{Q}_2.$$ (9.1) This is in agreement with Strominger’s observation that the $`AdS_2`$ Virasoro algebra is related to the right-moving $`AdS_3`$ Virasoro algebra by a topological twist which shifts the central charge to zero. It might be that the degeneracy of the ground state itself is effectively captured by a model of conformal quantum mechanics . However, our results support the point of view that the excitations above extremality are described by a two-dimensional conformal field theory .
warning/0004/gr-qc0004004.html	ar5iv	text	# The nonlinear evolution of de Sitter space instabilities ## I Introduction and Summary ### A Anti-evaporation and proliferation Schwarzschild-de Sitter black holes have unusual quantum properties and instabilities . They are of cosmological interest because they can be produced during inflation . They are also theoretically significant because they change the global structure of de Sitter space fundamentally. Unlike their asymptotically flat cousins, Schwarzschild-de Sitter black holes are surrounded by a cosmological horizon. This limits their size and puts them in a thermal bath. Their temperature is always larger than that of the cosmological horizon . Nevertheless, Hawking and one of the authors argued in Ref. that some black holes accrete so much quantum radiation that they will grow, or ‘anti-evaporate,’ instead of evaporating (see also ). In the maximal solution, the black hole and cosmological horizon are of equal size, and the spatial geometry will be $`S^1\times S^2`$, with constant two-sphere radius. This geometry can nucleate semiclassically in de Sitter space through a gravitational tunneling process. Its evolution is unstable to the formation of $`n`$ de Sitter regions, distributed around the $`S^1`$ and connected by $`n`$ black hole throats . This configuration may be visualized as a ‘doughnut’ with $`n`$ ‘wobbles,’ whose minima and maxima represent the black hole and cosmological horizons, respectively. For $`n=1`$, this reduces to the ordinary, submaximal Schwarzschild-de Sitter solution. For higher $`n`$, it resembles a ring-like sequence of $`n`$ Schwarzschild-de Sitter solutions. As the de Sitter regions form between the cosmological horizons, the black holes are expected to evaporate and shrink. When they finally disappear, the doughnut is pinched in $`n`$ places, leaving $`n`$ disconnected pieces. This corresponds to the fragmentation of space into $`n`$ large daughter universes . Thus de Sitter space is globally unstable. Locally, however, the daughter universes are indistinguishable from de Sitter. Therefore, they will harbour more maximal black holes, whose evaporation can lead to further fragmentation. Iteratively, an unbounded number of disconnected de Sitter universes is produced: de Sitter space proliferates. This drastic conclusion depends crucially on the nonlinear evolution of the instabilities of the maximal Schwarzschild-de Sitter solution (also referred to as the ‘Nariai’ solution below). ### B Analytical shortcomings In Refs. and , the anti-evaporation and proliferation effects were derived analytically, using a number of approximations. The stability of the Nariai geometry was studied using linear perturbation theory. The evolution of the horizons was found only for very early times, as a power series in the time variable. This led to the discovery of anti-evaporation for a class of perturbations. But a power series can prove only initial anti-evaporation. It cannot answer an important question: Will such black holes continue to grow, asymptotically approaching the maximal (Nariai) solution? Or will their growth eventually stall, and evaporation set in after all? This problem is highlighted by recent claims that anti-evaporation may permanently stabilize black holes nucleated during inflation<sup></sup><sup></sup>*We point out, however, that even stable black holes nucleated semiclassically during inflation would be diluted by the cosmological expansion. Typically, they will lie outside the current horizon and will not be observable . . At least for Euclidean boundary conditions, it was shown that the black holes do evaporate at late times . In fact, arguments might be made that the evaporating mode is attractive at late times for all initial conditions . In any case, however, the intermediate period containing the putative transition between anti-evaporation and evaporation is beyond the reach of the analytical methods employed. Can there be a smooth transition, and how does it proceed? To understand this crucial phase of the black hole evolution, it should be explored numerically. If anti-evaporation is only temporary, then how can we be sure that the evaporating phase, once entered, continues until the black hole disappears? This question is particularly important in the context of multiple black hole formation and proliferation . The fragmentation of space is possible only if all black holes become small and disappear like Schwarzschild black holes. To show that black holes really evaporate completely, it is necessary to trace the backreaction at least until they are significantly smaller (perhaps by a factor of 10) than the cosmological horizon. Then they can be assumed to behave like Schwarzschild black holes, which presumably radiate all of their mass away. But linear perturbation theory requires that the small two-spheres inside the black hole and the large two-spheres between the cosmological horizons still be of nearly equal size. Thus, the analytic ‘late-time’ solution, showing black hole evaporation, is not really valid for arbitrarily late times. A full, non-linear numerical calculation is needed to interpolate to the Schwarzschild regime. ### C Numerical analysis In Sec. III we consider $`n=1`$ perturbations of the Nariai geometry, which can be parametrized by their initial amplitude, $`ϵ`$, and phase, $`\vartheta `$. Depending on $`\vartheta `$, the black holes undergo initial periods of anti-evaporation or evaporation. We demonstrate the transitory nature of anti-evaporation for neutral black holes by showing that evaporation sets in at sufficiently late times, independently of $`\vartheta `$. The evaporation rate is compared to an approximate solution of the linear equations in order to verify that this solution is an attractor. Beyond the linear regime, the black holes are found to shrink to a size much smaller than the cosmological horizon before our computation breaks down. Then they live in approximately flat space, and should simply continue to evaporate until they disappear. In Sec. IV we show that maximal black holes become stable if their charge exceeds a critical value, in agreement with thermodynamic expectations. This confirms analytic arguments given in Ref. . In particular, we show that supercritical black holes, independently of their initial behavior, always end up anti-evaporating. This is the condition for a novel type of daughter-universe production proposed in Ref. . In Sec. V we show that a higher mode ($`n>1`$) perturbation behaves as predicted by analytical arguments: First it oscillates, then it freezes out, forming black hole interiors and growing de Sitter regions. Depending on initial conditions, a period of anti-evaporation may follow. Finally, the black holes evaporate and become much smaller than the cosmological horizons. ## II Action and perturbations ### A Including back-reaction The four-dimensional Lorentzian Einstein-Hilbert action with a cosmological constant, $`\mathrm{\Lambda }`$, and a Maxwell field, $`F_{\mu \nu }`$, is given by: $$S=\frac{1}{16\pi }d^4x(g^{\mathrm{IV}})^{1/2}\left[R^{\mathrm{IV}}2\mathrm{\Lambda }F_{\mu \nu }F^{\mu \nu }\right],$$ (1) where $`R^{\mathrm{IV}}`$ and $`g^{\mathrm{IV}}`$ are the four-dimensional Ricci scalar and metric determinant. Restricting to spherically symmetric fields and quantum fluctuations, the metric may be written as $$ds^2=e^{2\rho }\left(dt^2+dx^2\right)+e^{2\varphi }d\mathrm{\Omega }^2,$$ (2) where $`x`$ is the coordinate on the $`S^1`$, with period $`2\pi `$. Using this ansatz, and the on-shell condition for magnetic fields, $$F_{\mu \nu }F^{\mu \nu }=2Q^2e^{4\varphi },$$ (3) the angular coordinates and the Maxwell field can be integrated out in Eq. (1), which reduces the action to $$S=\frac{1}{16\pi }d^2x(g)^{1/2}e^{2\varphi }\left[R+2(\varphi )^2+2e^{2\varphi }2\mathrm{\Lambda }2Q^2e^{4\varphi }\right],$$ (4) In order to include back-reaction effects, a Polyakov term, which arises in the one-loop effective action of a two-dimensional scalar field, will be included. (It would be preferable to work with dilaton-coupled scalars , or even better, with a four-dimensional effective action reduced to two-dimensions. This would complicate the numerical computation enormously and will not be attempted here. For small perturbations, it has been shown that extra terms from dilaton-coupling do not affect results . Thus one would not expect qualitative changes even for black holes noticeably smaller than the cosmological horizon. In order to corroborate our results for nearly-Schwarzschild black holes, however, a full four-dimensional treatment would be desirable in future work.) In the large $`N`$ limit, the contribution from the quantum fluctuations of the scalars dominates over that from the metric fluctuations. In order for quantum corrections to be small, one should take $`N\mathrm{\Lambda }1`$. One can obtain a local form of this action by introducing an independent scalar field $`Z`$ which mimics the trace anomaly . The on-shell equivalence of the equations of motion can be seen by choosing a conformal gauge for the two-dimensional metric, as in Eq. (2). Thus the action of the one-loop model will be given by: $$S=\frac{1}{16\pi }d^2x(g)^{1/2}\left[\left(e^{2\varphi }+\frac{N}{3}Z\right)R\frac{N}{6}\left(Z\right)^2+2+2e^{2\varphi }\left(\varphi \right)^22e^{2\varphi }\mathrm{\Lambda }2Q^2e^{2\varphi }\right].$$ (5) ### B Equations of motion Differentiation with respect to $`t`$ ($`x`$) will be denoted by an overdot (a prime). For any functions $`f`$ and $`g`$, define: $$fg\dot{f}\dot{g}+f^{}g^{},^2g\ddot{g}+g^{\prime \prime },$$ (6) $$\delta f\delta g\dot{f}\dot{g}+f^{}g^{},\delta ^2g\ddot{g}+g^{\prime \prime }.$$ (7) Variation with respect to $`\rho `$, $`\varphi `$ and $`Z`$ yields the following equations of motion: $$^2\varphi +2(\varphi )^2+\frac{N}{6}e^{2\varphi }^2Z+e^{2\rho +2\varphi }\left(\mathrm{\Lambda }e^{2\varphi }+Q^2e^{2\varphi }1\right)=0;$$ (8) $$^2\rho ^2\varphi +(\varphi )^2+\mathrm{\Lambda }e^{2\rho }Q^2e^{2\rho +4\varphi }=0;$$ (9) $$^2Z2^2\rho =0.$$ (10) The constraint equations are: $$\left(\delta ^2\varphi 2\delta \varphi \delta \rho \right)(\delta \varphi )^2=\frac{N}{12}e^{2\varphi }\left[(\delta Z)^2+2\delta ^2Z4\delta Z\delta \rho \right];$$ (11) $$\left(\dot{\varphi }^{}\dot{\rho }\varphi ^{}\rho ^{}\dot{\varphi }\right)\dot{\varphi }\varphi ^{}=\frac{N}{12}e^{2\varphi }\left[\dot{Z}Z^{}+2\dot{Z}^{}2\left(\dot{\rho }Z^{}+\rho ^{}\dot{Z}\right)\right].$$ (12) From Eq. (10), it follows that $`Z=2\rho +\eta `$, where $`\eta `$ satisfies $`^2\eta =0`$. The remaining freedom in $`\eta `$ can be used to satisfy the constraint equations for any choice of $`\rho `$, $`\dot{\rho }`$, $`\varphi `$ and $`\dot{\varphi }`$ on an initial spacelike section . ### C Metric and horizon perturbation Maximal (Nariai) black holes nucleate semiclassically in de Sitter space. This process is mediated by gravitational instantons and has been described in detail in Refs. . Here we are interested not in the nucleation, but in the further evolution of the Nariai solution. Its metric is given by $$e^{2\rho }=\frac{1}{A}\frac{1}{\mathrm{cos}^2t},e^{2\varphi }=B,$$ (13) where $`B`$ is given by the cubic equation $$B\left[1Q^2B\left(1+\frac{N}{3}B\right)\right]=\left[1\frac{N}{3}B\right]\mathrm{\Lambda };$$ (14) the physical solution is the one that limits to $`B=\mathrm{\Lambda }`$ for $`N=Q=0`$). $`A`$ is given by $$A=\mathrm{\Lambda }Q^2B^2.$$ (15) Quantum fluctuations will perturb this solution, so that the two-sphere radius, $`e^\varphi `$, will vary slightly along the one-sphere coordinate, $`x`$. Decomposition into Fourier modes on the $`S^1`$ yields the perturbation ansatz $$e^{2\varphi }=\mathrm{\Lambda }_2\left[1+2ϵ\underset{n}{}\left(\sigma _n(t)\mathrm{cos}nx+\stackrel{~}{\sigma }_n(t)\mathrm{sin}nx\right)\right],$$ (16) where $`ϵ`$ is taken to be small. This will be referred to as the metric perturbation, characterized by the $`\sigma _n`$, $`\stackrel{~}{\sigma }_n`$ at the time $`t=0`$. More generally, one should also consider perturbations of the time derivative of the two-sphere radius, expressed by $`\dot{\sigma }_n`$. As in Ref. , we parameterize the initial conditions for the perturbation as $$\sigma _n(0)=\mathrm{sin}\vartheta ,\dot{\sigma }_n(0)=\mathrm{cos}\vartheta ,$$ (17) where $`\vartheta `$ represents the phase of the initial perturbation. Based on linear perturbation theory, one would expect such perturbations to lead to a classical, as well as a quantum instability . Classically, the regions on the $`S^1`$ where the two-spheres are smaller than the Nariai value should collapse to form black hole interiors. The larger two-spheres, on the other hand, should grow exponentially, developing into asymptotic de Sitter regions. If the first mode dominates, this simply leads to a nearly-maximal Reissner-Nordström-de Sitter solution. But if higher modes are strongly excited, a whole sequence of Reissner-Nordström-de Sitter solutions can develop around the same one-sphere. This may be thought of as a necklace of de Sitter regions, strung together by black hole throats. The expected quantum instability is the evaporation of these black holes. Actually, we will show below that black holes of sufficient charge are stable. But all other black holes would be expected to evaporate and get smaller, since their temperature is higher than that of the surrounding cosmological horizon. The black hole evolution can be calculated numerically using Eqs. (8)–(10). It is important to stress that the black hole and cosmological horizons do not in general correspond to the minimal and maximal two-spheres along the $`S^1`$, although such a slicing can always be found. In general, one must first find the positions of the horizons on the one-sphere; this can be done by finding the points where the gradient of the two-sphere size is null . The two-sphere sizes at those locations give the size of the black hole. The black hole evolution can be monitored by following the horizon location and plotting the horizon size vs. time. It will be convenient to define the horizon perturbation $`\delta `$, for a black hole located at $`x_\mathrm{b}`$, by $$r_\mathrm{b}(t)^2=e^{2\varphi [t,x_\mathrm{b}(t)]}=B\left[1+2ϵ\delta (t)\right].$$ (18) Thus $`\delta `$ corresponds to the fractional difference between the current black hole size and the size of a maximal black hole of equal charge. Evaporation corresponds to increasing values of $`\delta `$. Because the $`S^1`$ expands exponentially, the black hole and cosmological horizons of a Reissner-Nordström-de Sitter solution appear to move ever more closely together in comoving $`S^1`$ coordinates (see Fig. 1). ### D Numerical technique The equations of motion (9) – (10) were solved numerically using the standard method of characteristics for second order quasi-linear hyperbolic equations . For this purpose, the auxiliary field $`Z`$ can be eliminated by inserting Eq. (10) into Eq. (9) and the remaining equations can be rewritten in manifestly hyperbolic form: $$^2\varphi =𝒞\left[e^{2\rho }\left(3\mathrm{\Lambda }e^{2\varphi }(3+N\mathrm{\Lambda })+3e^{4\varphi }Q^2+e^{6\varphi }NQ^2\right)(e^{2\varphi }N6)(\varphi )^2\right]$$ (19) $$^2\rho =3𝒞\left[e^{2(\rho +\varphi )}\left(2e^{2\varphi }Q^21\right)+(\varphi )^2\right],$$ (20) where $$𝒞=\frac{1}{3e^{2\varphi }N}.$$ (21) The charateristic curves, along which the equations reduce to ordinary differential equations, are $`x_\pm =x\pm t`$. Given the solution at two neighboring points on non-identical characteristics and approximating $`dx_\pm `$ by $`\mathrm{\Delta }x_\pm `$, the solution at the intersection of the characteristic curves going through these points can be found by iteration. Assigning the initial conditions Eqs. (13) – (16) on a $`t=0`$ hypersurface thus allows the advancement of the solution along the $`x_\pm `$-grid, employing periodic boundary conditions at $`x=0`$ and $`x=2\pi `$. For most of our computations, we used a resolution of 10000 grid points on $`x`$-hypersurfaces. The accuracy of the results was verified by comparison with perturbative solutions and with results obtained from an independent second-order finite difference code for the same equations. With the exception of the very last time step, the closest to the conformal time coordinate singularity, excellent agreement of all solutions was found. ## III Anti-evaporation and turnaround Anti-evaporation was first found in Ref. for the $`\vartheta =\pi /2`$, $`n=1`$ perturbation of the maximal (Nariai) neutral black hole solution in de Sitter space. A power series approximation showed that black holes formed by this perturbation will grow initially. A numerical calculation, however, allows us to probe the black hole evolution beyond the range of validity of the power series. Interestingly, this proves the anti-evaporation effect to be transitory. This is demonstrated by the long-dashed line in Fig. 3. The horizon perturbation first decreases quadratically, signaling anti-evaporation, in quantitative agreement with the result of Ref. . At a time $`t.2`$, however, it turns around and starts to increase. This means that the black hole eventually stops growing. Instead, it starts to shrink in size, corresponding to evaporation at late times. For other initial phases, the behavior of the black hole can be quite complicated, going through various periods of evaporation and anti-evaporation (Fig. 3). The important result is, however, that it always ends up evaporating at late times, as conjectured in Ref. . The numerical result for the final evaporation rate agrees with the asymptotic solutions given there to within better than one percent: $$\delta (\mathrm{cos}t)^{1c_+},$$ (22) where $`c_+`$ is the larger root of $$c(c+1)=2\sqrt{1+\frac{2}{3}N\mathrm{\Lambda }+\frac{1}{9}N^2\mathrm{\Lambda }^2}.$$ (23) An initial phase just under $`\vartheta =\pi /2`$ corresponds to a large time-derivative of the metric perturbation directed opposite to the perturbation. In this case, the metric overshoots the Nariai value and a cosmological horizon forms. The corresponding black hole is located where the cosmological horizon would usually sit. Its evaporation is mirrored by the growth of the cosmological horizon. In Fig. 3 this cosmological horizon shows up for $`\vartheta =5\pi /6`$ as a negative value of $`\delta `$ which increases in magnitude. How can we be sure that the evaporation seen in the later stages of Fig. 3 is permanent? By the nature of the conformal time $`t`$, an infinite amount of proper time is contained in the final step of the numerical calculation. No matter how refined the step size, this region cannot be resolved. It is sufficient, however, to trace the evolution until the black hole size differs significantly from the size of the cosmological horizon. This is precisely the condition for non-linearity, so it cannot be verified analytically using the power series for $`\delta `$. Our numerical calculation, however, can reach this regime (see Fig. 3). Once it is so small, the black hole will be much hotter than the cosmological radiation, and will no longer be influenced by it. It will behave like a Schwarzschild black hole, so we can trust that it will continue to evaporate. ## IV Stability of charged black holes Nariai solutions can be electrically or magnetically charged. In this case, their spacelike sections will still be given by a direct product of $`S^1\times S^2`$ (with constant $`S^2`$ radius), but the $`S^1`$ and $`S^2`$ radii will not be equal. The flux runs around the $`S^1`$. The metric is classically unstable to small perturbations. Reissner-Nordström-de Sitter black holes form where the two-sphere size is smaller, and asymptotic de Sitter regions develop where the two-spheres are larger . These black holes are of nearly maximal size. For a small charge, one would expect their initial behavior to be similar to the neutral Schwarzschild-de Sitter black holes. When they are highly charged, however, they become nearly extremal. The temperature of a black hole decreases as it approaches extremality. Thermodynamic arguments were given in Ref. which showed that there is a critical value of the charge $`Q_\mathrm{C}^2=3/(16\mathrm{\Lambda })=3/4Q_{\mathrm{max}}^2`$. Roughly speaking, when a subcritical Nariai solution is perturbed, the temperature of the black hole will be larger than that of the cosmological horizon because it is smaller. When a supercritial Nariai solution is perturbed, the black hole will be colder than the cosmological horizon, even though it is smaller, simply because its mass is already very close to the mass of the extremal, zero-temperature solution. Thus one would expect subcritically charged black holes to evaporate at late times. Supercritically charged black holes, on the other hand, should absorb quantum radiation from the cosmological horizon and grow, their size approaching the Nariai value asymptotically. This is confirmed numerically, as shown in Figs. 5 and 5. The value of the critical charge is confirmed quantitatively. This is a non-trivial check that this simple, two-dimensional model reflects the thermodynamic properties of Reissner-Nordström-de Sitter black holes accurately. The anti-evaporation of supercritially charged black holes has an important consequence for the quantum global structure of de Sitter space. It gives rise to a new type of proliferation effect, as pointed out in Ref. . Because the two-sphere size is forever nearly constant in the Hubble-size region between the two horizons, a small quantum fluctuation can easily invert the role of the horizons. In other words, it can increase the two-sphere size on the black hole horizon, turning it into a cosmological horizon, and vice-versa. This amounts to the insertion of a new black hole ‘bead’ into the $`S^1`$ ‘necklace.’ Such processes repeat endlessly, so that an unbounded number of causally disconnected de Sitter regions develop on the same $`S^1`$. ## V Higher modes and proliferation We now return to uncharged solutions, but consider modes with $`n>1`$. Such perturbations lead to a spatial geometry that can be described as a ‘doughnut with $`n`$ wobbles.’ Like for $`n=1`$, the metric perturbation is classically unstable. The mode will oscillate until the $`S^1`$ expansion has stretched it enough to leave the horizon. Then it will freeze out, and grow exponentially. This was shown in a linear approximation in Ref. , and is demonstrated numerically in Fig. 6 for $`n=3`$. The evolution of the horizon perturbation for $`n=3`$ is shown in Fig. 8. While the metric perturbation oscillates, the $`(\varphi )^2=0`$ surfaces move rapidly and ‘cross over’ (see Fig. 8). They only represent black hole horizons after the metric perturbation freezes out. Fig. 8 shows that in the $`n=3,\vartheta =\pi /2`$ case considered there, the metric perturbation freezes out while it of the opposite sign compared to its initial value. This means that cosmological horizons develop from the initial minimal two-spheres. This is reflected in the figure in the negative values for the ‘black hole horizon’ perturbation. Note that the absolute value is increasing at late times, signalling evaporation. ## Acknowledgments We thank the Aspen Center for Physics where this investigation was initiated, and the German American Academic Council (GAAC) for financial support.
warning/0004/hep-th0004189.html	ar5iv	text	# References STUPP–00–159 April, 2000 Logarithmic Behaviours in the Feigin-Fuchs Construction of the $`c=2`$ Conformal Field Theory Hiroki Hata <sup></sup><sup></sup>* E-mail: hata@krishna.th.phy.saitama-u.ac.jp and Shun-ichi Yamaguchi E-mail: yama@krishna.th.phy.saitama-u.ac.jp Department of Physics, Faculty of Science Saitama University, Urawa, Saitama 338-8570, Japan Abstract We obtain logarithmic behaviours of a four-point correlation function in the $`c=2`$ conformal field theory by using the Feigin-Fuchs construction. It becomes an indeterminate form by a naive evaluation, but is obtained by introducing an appropriate regularization procedure. Conformal field theories whose correlation functions have logarithmic behaviour were first studied by Gurarie in the central charge $`c=2`$ model . This model is one of the simplest systems since the correlation function in the problem consists of only one kind of primary fields. Logarithmic conformal field theories have interesting properties: new operators are needed, which are called logarithmic operators and never appeared for ordinary conformal field theories. The origin of the logarithms of Ref. is a hypergeometric function, which is a solution of the differential equation for the four-point correlation function and equivalent to the complete elliptic integral of the first kind. The origin should also be explained by the Feigin-Fuchs construction , but the approach only to the logarithmic operators is studied . The construction also applies to other models . In this paper the four-point correlation function of Ref. is calculated by using the Feigin-Fuchs construction. In this construction the correlation function is given by an integral representation. We find that its integral value becomes an indeterminate form, $`\frac{0}{0}`$. In order to evaluate this form we introduce an appropriate regularization procedure: we perform an analytic continuation of a parameter in the hypergeometric function, namely, take the limit $`c2`$ , and evaluate the form of the correlation function by using essentially the de l’Hospital theorem. In this way the logarithmic term appears, and our result is in agreement with that of Ref. . Application of our method to generic four-point correlation functions with logarithms is under studying. We now consider the Feigin-Fuchs construction of conformal field theories. The action is given by $$S=\frac{1}{8\pi }d^2\xi \sqrt{g}\left(g^{\mu \nu }_\mu \varphi _\nu \varphi +2i\alpha _0R\varphi \right),$$ (1) where $`\varphi `$ is a real scalar field and $`R`$ is the scalar curvature on a sphere with fixed reference metric $`g_{\mu \nu }`$. The parameter $`2\alpha _0`$ can be interpreted as the background charge. On the complex plane, the energy momentum tensor is of the form $$T(z)=\frac{1}{2}_z\varphi _z\varphi +i\alpha _0_z^{\mathrm{\hspace{0.17em}2}}\varphi ,$$ (2) and two-point function of the field $`\varphi `$ is $`\varphi (z,\overline{z})\varphi (w,\overline{w})=\mathrm{ln}\|zw\|^2`$. Thus the central charge of the system is written in terms of $`\alpha _0`$ as $$c=112\alpha _0^{\mathrm{\hspace{0.17em}2}}.$$ (3) When considering correlation functions of primary fields $`\mathrm{\Phi }`$, we treat the correlation function in terms of the vertex operators $`V_\alpha e^{i\alpha \varphi }`$ instead of $`\mathrm{\Phi }`$. The conformal weight $`h`$ of the operator $`V_\alpha `$ is $$h(V_\alpha )=h(V_{2\alpha _0\alpha })=\frac{1}{2}\alpha (2\alpha _0\alpha ),$$ (4) so two admissible operators exist for one field $`\mathrm{\Phi }`$. For the $`(p,q)`$ primary fields $`\mathrm{\Phi }_{p,q}`$ , the conformal weight takes the discrete value $$h_{p,q}=\frac{1}{2}\alpha _0^{\mathrm{\hspace{0.17em}2}}+\frac{1}{8}(p\alpha _++q\alpha _{})^2,$$ (5) where $`\alpha _\pm =\alpha _0\pm \sqrt{\alpha _0^{\mathrm{\hspace{0.17em}2}}+2}`$. The corresponding operator $`V_\alpha `$ has the parameter $`\alpha `$ given by $$\alpha _{p,q}=\alpha _0\frac{1}{2}\left(p\alpha _++q\alpha _{}\right).$$ (6) We also need the screening charges $$Q_\pm =d^2ue^{i\alpha _\pm \varphi (u,\overline{u})}.$$ (7) A certain number of screening charges should be inserted in the correlation function so that the charge neutrality condition, required by the zero mode integration of $`\varphi `$, is satisfied. Since the screening charges are the integrals of the operators with conformal weight one, the insertion have no effects on the conformal properties of correlation functions. Let us now concretely consider the four-point correlation function of Ref. $$G^{(4)}\mu (z_1,\overline{z}_1)\mu (z_2,\overline{z}_2)\mu (z_3,\overline{z}_3)\mu (z_4,\overline{z}_4)$$ (8) in the $`c=2`$ model. Here $`\mu (z,\overline{z})\mathrm{\Phi }_{1,2}(z,\overline{z})`$ is the primary field with the conformal weight $`h_{1,2}=\frac{1}{8}`$. In the Feigin-Fuchs construction, four-point correlation functions of primary fields $`\mathrm{\Phi }_i`$’s in general take the form $$\mathrm{\Phi }_1(z_1,\overline{z}_1)\mathrm{\Phi }_2(z_2,\overline{z}_2)\mathrm{\Phi }_3(z_3,\overline{z}_3)\mathrm{\Phi }_4(z_4,\overline{z}_4)=\underset{i=1}{\overset{4}{}}e^{i\alpha _i\varphi (z_i,\overline{z}_i)}(Q_+)^m(Q_{})^n,$$ (9) where the charge neutrality condition is $$\underset{i=1}{\overset{4}{}}\alpha _i+m\alpha _++n\alpha _{}=2\alpha _0.$$ (10) In our model with $`\alpha _0=\frac{1}{2},\alpha _+=2,\alpha _{}=1`$ and $`\alpha _i=\alpha _{1,2}=\frac{1}{2}`$ we choose $`m=0`$ and $`n=1`$. The correlation function (8) is, therefore, evaluated as $$G^{(4)}=\underset{i=1}{\overset{4}{}}e^{i\frac{1}{2}\varphi (z_i,\overline{z}_i)}Q_{}=\underset{i<j}{}\|z_{ij}\|^{\frac{1}{2}}d^2u\underset{i=1}{\overset{4}{}}\|z_iu\|^1,$$ (11) where $`z_{ij}=z_iz_j`$. On the other hand, from the SL(2,C) Ward identity , the correlation function (8) can be written as $$G^{(4)}=F(x,\overline{x})\underset{i<j}{}\|z_{ij}\|^{\frac{1}{6}},$$ (12) where $`F`$ is an arbitrary function of the SL(2,C) invariant cross ratios $$x=\frac{z_{12}z_{34}}{z_{13}z_{24}},\overline{x}=\frac{\overline{z}_{12}\overline{z}_{34}}{\overline{z}_{13}\overline{z}_{24}}.$$ (13) From Eqs. (11) and (12), by fixing $`z_1=0,z_2=x,z_3=1,z_4=\mathrm{}`$, $`F(x,\overline{x})`$ can be evaluated and we obtain $$G^{(4)}=\|z_{13}z_{24}\|^{\frac{1}{2}}\|x(1x)\|^{\frac{1}{2}}I(\frac{1}{2},\frac{1}{2},\frac{1}{2};x),$$ (14) where $$I(a,b,c;x)=d^2u\|u\|^{2a}\|1u\|^{2b}\|ux\|^{2c}.$$ (15) The integral $`I(a,b,c;x)`$ can be transformed into a sum of squares of line integrals $$I(a,b,c;x)=\frac{\mathrm{sin}[\pi (a+b+c)]\mathrm{sin}(\pi b)}{\mathrm{sin}[\pi (a+c)]}\|I_1(x)\|^2+\frac{\mathrm{sin}(\pi a)\mathrm{sin}(\pi c)}{\mathrm{sin}[\pi (a+c)]}\|I_2(x)\|^2,$$ (16) where $`I_1(x)`$ $`=`$ $`{\displaystyle _1^{\mathrm{}}}𝑑uu^a(u1)^b(ux)^c`$ (17) $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(abc1)\mathrm{\Gamma }(b+1)}{\mathrm{\Gamma }(ac)}}F(c,abc1,ac;x),`$ $`I_2(x)`$ $`=`$ $`{\displaystyle _0^x}𝑑uu^a(1u)^b(xu)^c`$ (18) $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(a+1)\mathrm{\Gamma }(c+1)}{\mathrm{\Gamma }(a+c+2)}}x^{a+c+1}F(b,a+1,a+c+2;x).`$ The coefficients of $`\|I_1(x)\|^2`$ and $`\|I_2(x)\|^2`$ in Eq. (16) are determined by the monodoromy invariance of $`I(a,b,c;x)`$. We see that the result (14) is naively an indeterminate form: $$I(\frac{1}{2},\frac{1}{2},\frac{1}{2};x)=\frac{1}{\mathrm{sin}(\pi )}\left[\|I_1(x)\|^2+\|I_2(x)\|^2\right]\frac{0}{0},$$ (19) since $`I_1(x)=I_2(x)=\pi F(\frac{1}{2},\frac{1}{2},1;x)`$. To evaluate the above indeterminate form we now introduce a regularization procedure as follows: $`I(\frac{1}{2},\frac{1}{2},\frac{1}{2};x)`$ $``$ $`\underset{a\frac{1}{2}}{lim}I(a,\frac{1}{2},\frac{1}{2};x)`$ (20) $`=`$ $`{\displaystyle \frac{\frac{d}{da}\left\{\mathrm{sin}[\pi (a1)]\|I_1(x)\|^2\mathrm{sin}(\pi a)\|I_2(x)\|^2\right\}}{\frac{d}{da}\mathrm{sin}[\pi (a\frac{1}{2})]}}\|_{a=\frac{1}{2}}.`$ Note that $`I_1(x)`$ and $`I_2(x)`$ are the functions of the regularization parameter $`a`$, which are given by Eqs. (17) and (18) with $`b=c=\frac{1}{2}`$. Since $`I_1(x)=I_2(x)`$ when $`a=\frac{1}{2}`$, Eq. (20) becomes $$I(\frac{1}{2},\frac{1}{2},\frac{1}{2};x)=\frac{1}{\pi }[I_1(\overline{x})\left\{\frac{d}{da}(I_1(x)I_2(x))\|_{a=\frac{1}{2}}\right\}+(x\overline{x})].$$ (21) The factor of differential in the braces is evaluated as $`{\displaystyle \frac{d}{da}}\left(I_1(x)I_2(x)\right)\|_{a=\frac{1}{2}}`$ $`=`$ $`\pi \mathrm{ln}\left({\displaystyle \frac{x}{16}}\right)F(\frac{1}{2},\frac{1}{2},1;x)2\pi M(x)`$ (22) $``$ $`2\pi \stackrel{~}{F}(x),`$ where $$M(x)=\frac{1}{\pi }\underset{n=1}{\overset{\mathrm{}}{}}\left[\frac{\mathrm{\Gamma }\left(n+\frac{1}{2}\right)}{n!}\right]^2\left[\psi (1)\psi (n+1)\psi (\frac{1}{2})+\psi (n+\frac{1}{2})\right]x^n$$ (23) and $`\psi (x)`$ is the digamma function. Anti-holomorphic part can be evaluated in the same way. Therefore we finally obtain $$I(\frac{1}{2},\frac{1}{2},\frac{1}{2};x)=2\pi [F(\frac{1}{2},\frac{1}{2},1;x)\stackrel{~}{F}(\overline{x})+(x\overline{x})].$$ (24) Notice that from Appendix C of Ref. $`\stackrel{~}{F}(x)`$ satisfies the following relation $$\stackrel{~}{F}(x)=\frac{\pi }{2}F(\frac{1}{2},\frac{1}{2},1;1x)=_0^{\frac{\pi }{2}}\frac{d\theta }{\sqrt{1(1x)\mathrm{sin}^2\theta }}.$$ (25) This is just the function $`G(1x)`$ in Eq. (10) of Ref. , which is the origin of logarithm. Thus our result (14) has logarithmic behaviour as $`G^{(4)}`$ $`=`$ $`\pi ^2\left\|z_{13}z_{24}\right\|^{\frac{1}{2}}\left\|x(1x)\right\|^{\frac{1}{2}}`$ (26) $`\times \left[F(\frac{1}{2},\frac{1}{2},1;x)F(\frac{1}{2},\frac{1}{2},1;1\overline{x})+F(\frac{1}{2},\frac{1}{2},1;\overline{x})F(\frac{1}{2},\frac{1}{2},1;1x)\right].`$ This agrees with that of Ref. up to overall constant, which was obtained by directly solving the hypergeometric differential equation. In the above procedure we performed an analytic continuation of the first parameter $`a`$ of the function $`I(a,b,c;x)`$. We can reproduce the same result by using the third parameter $`c`$ but not by the second one $`b`$. The fact depends on the choice of the two independent contours of integrals (17) and (18) since the coefficients of $`\|I_1(x)\|^2`$ and $`\|I_2(x)\|^2`$ in Eq. (16) are determined by monodoromy invariance of Eq. (16). Acknowledgements The authors would like to thank Y. Tanii for careful reading of the manuscript. One of the authors (S.Y.) is also grateful to K. Hida and C.-B. Kim for useful discussions.
warning/0004/physics0004045.html	ar5iv	text	# High-harmonic generation from a confined atom ## Abstract The order of high harmonics emitted by an atom in an intense laser field is limited by the so-called cutoff frequency. Solving the time-dependent Schrödinger equation, we show that this frequency can be increased considerably by a parabolic confining potential, if the confinement parameters are suitably chosen. Furthermore, due to confinement, the radiation intensity remains high throughout the extended emission range. All features observed can be explained with classical arguments. 32.80.Rm, 42.65.Ky, 42.50.Hz Typical features of the emission spectra of an atom in a strong laser field, known as “the plateau” and “the cutoff”, are a wide frequency region with harmonics of comparable intensities, and an abrupt intensity decrease at the high-energy-end of the plateau. For a monochromatic driving field, the cutoff energy is given by $`\epsilon _{\mathrm{max}}=\|\epsilon _0\|+3.17U_p`$, where $`\|\epsilon _0\|`$ and $`U_p`$ are the ionization potential and the ponderomotive energy, respectively. This simple cutoff law, derived by classical means only , or using more refined methods , corresponds to the physical picture referred to as the “three-step model” : A bound electron exposed to the laser field leaves the atom through tunneling at a time $`t_0`$ (step 1), propagates in the continuum, being driven back towards its parent ion at a later time $`t_1`$ (step 2), and finally falls back to a bound state under emission of high harmonics (step 3). This scenario describes the spectral features observed experimentally very well . The cutoff frequency, in quantitative agreement with the experiment, is related to the maximum kinetic energy the electron has upon return, $`E_{\mathrm{kin}}(t_1,t_0)`$. According to this picture, in order to increase the cutoff energy, one must increase the kinetic energy of the returning electron. Indeed, the exisiting proposals to extend the plateau towards higher energies reach a higher value of $`E_{\mathrm{kin}}(t_1,t_0)`$ by different means. However, this does not necessarily imply an efficient generation of high-order harmonics up to this larger cutoff energy. For instance, a rather complex situation with several “cutoffs ” emerges by using bichromatic fields with driving waves of comparable intensities. An illustrative example is presented in , using a driving field of linearly polarized monochromatic light of frequency $`\omega `$ and its second harmonic. Under such conditions the monochromatic cutoff, as a function of the field-strength ratio between the two driving waves, splits into two branches. Thereby, the upper branch extends up to $`\|\epsilon _0\|+5U_p`$. However, the harmonics emerging up to the cutoff of the upper branch are weak compared to those from the lower branch and therefore irrelevant to the emission spectrum. The reason is simple: The intensity of the harmonics is strongly influenced by step 1 which is the tunneling process out of the binding potential under the influence of the field. If the field amplitude is small at the emission time $`t_0`$ (which is the case for the upper branch) then the tunneling barrier is large and the generated harmonics will be weak compared to those which originate from an effective tunneling process (as it is the case for the lower branch). Another idea to increase the cutoff energy is to use a static electric field. It provides an additional force which accelerates the electron towards the atomic core resulting in a higher kinetic energy $`E_{\mathrm{kin}}(t_1,t_0)`$. Indeed, it has been demonstrated that with an electric field whose strength is only a few percent of the amplitude of the laser field one can considerably enlarge the cutoff energy . However, the scheme suffers from two principal limitations. First, the increased kinetic energy occurs mainly for electrons with long excursion times. Due to wave packet spreading, those trajectories have negligible influence on the harmonic spectra. This problem has been overcome by introducing an additional magnetic field to restrict the spreading . A second, more severe limitation is the pronounced bound-state depletion caused by the static electric field: the atom is irreversibly ionized within a few field cycles, such that no appreciable high-harmonic generation takes place. The bound-state depletion which prevents an effective extension of the high-harmonic frequency points to the principle dilemma easily described in the picture of the returning electron: To extend the plateau and increase the cutoff, a kinetic energy of the returning electron, as large as possible, is desirable. On the other hand, an electron with such a high energy will leave the atom and is lost for the possible generation of high harmonics in consecutive laser cycles. Hence, we need a mechanism which brings an electron back to the nucleus, despite the fact that it has a kinetic energy so high that it would be irreversibly driven away from the core. Naively, a simple wall for the electron should already do this. However, one must avoid that the abrupt reflection of the charged electron at a wall leads to Bremsstrahlung which masks the desired high-harmonic generation of the atom. In the following we will show that the idea of bringing back the fast electron by an additional confinement and thereby extending the cutoff for the spectrum without additional depletion does indeed work for a suitably soft confinement potential. We consider a one-dimensional situation, which is a reasonable approximation for linearly polarized light. Atomic units are used throughout. The binding of the electron is described by the potential $$V_\mathrm{a}\left(x\right)=1.1\mathrm{exp}\left(x^2/1.21\right),$$ (1) which supports a single bound state $`\|0`$ at energy $`\epsilon _0=0.58\mathrm{a}.\mathrm{u}`$., corresponding to the Argon ionization potential. The system is exposed to a monochromatic laser field $`E(t)=E_0\mathrm{sin}\omega t`$ and the additional confining potential (Fig.1) $`V_\mathrm{h}(x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_\mathrm{h}^2}{2}}x^2h(x),`$ () $`h(x)`$ $`=`$ $`\{\begin{array}{c}1,\|x\|<x_0\\ \mathrm{cos}\left(\frac{\pi }{2}\theta \right),x_0\|x\|x_{\mathrm{max}}\\ 0,\|x\|>x_{\mathrm{max}}\end{array}`$ () with $`\theta =(\|x\|x_0)/(x_{\mathrm{max}}x_0)`$. The parameter $`x_0=nE_0/\omega ^2`$ is chosen to be a multiple of the electron excursion amplitude, and $`x_{\mathrm{max}}=2x_0.`$ Parabolic potential shapes are taken as a first approximation in several physical systems, as for instance electromagnetic traps or solid-state devices . Note, that for the parameter range chosen, identical emission spectra are obtained with and without truncation of the harmonic potential indicating that even in the truncated potential depletion has negligible influence. Thus, the electron does not reach the edges of $`V_\mathrm{h}\left(x\right)`$, which indicates an effective confinement. Futhermore, this shows that the confining potential does not generate harmonics itself. Therefore, high-harmonic generation still takes place only near the atomic core, for which the coordinate $`x`$ is considerably smaller than the electron excursion amplitude. The evolution of the electronic wave packet is described by the time-dependent Schrödinger equation $$i\frac{d}{dt}\|\psi (t)=\left[\frac{p^2}{2}+V(x)pA(t)\right]\|\psi (t),$$ (5) with $`V(x)=V_\mathrm{a}\left(x\right)+V_\mathrm{h}\left(x\right),`$ and the emission spectra are given by $$\sigma (\omega )=\left\|_0^{\mathrm{}}d(t)\mathrm{exp}[i\omega t]\right\|^2,$$ (6) where the dipole acceleration $$d(t)=\psi (t)\left\|dV(x)/dx+E(t)\right\|\psi (t)$$ (7) is computed by means of Ehrenfest’s theorem . We take the atom initially in the ground state $`\|0`$. Furthermore, in the results to be presented we chose $`x_0=73.87\mathrm{a}.\mathrm{u}.,`$ which corresponds to three times the excursion amplitude of an electron in a monochromatic field with $`E_0=0.08\mathrm{a}.\mathrm{u}.`$ and $`\omega =0.057\mathrm{a}.\mathrm{u}.`$. With these field parameters and a reasonable choice of $`\mathrm{\Omega }_\mathrm{h},`$ one indeed finds that the high-harmonic spectrum extends beyond the cutoff energy $`\epsilon _{\mathrm{max}}=\|\epsilon _0\|+3.17U_p`$ without significant loss of intensity, see Fig. 2. More specifically, we have determined a cutoff energy of $`\|\epsilon _0\|+4.55U_p`$ which is a 50% increase compared to the case without trapping. The classical argument for the cutoff energy applies to the situation with confinement as well and we find very good agreement between the cutoff in the quantum spectra (e.g. Fig. 2) and the classical cutoff. The latter has been determined in analogy to the situation without confinement: Starting with an electron of velocity zero, its trajectory is propagated under the influence of the laser field and the confinement potential $`V_\mathrm{h}`$ (but without the atomic potential $`V(x)=V_\mathrm{a}`$). We vary the initial time $`t_0`$ for which the electron leaves the atom within a field cycle, computing $`E_{\mathrm{kin}}(t_1,t_0)`$ for return times $`t=t_1`$ satisfying the condition $`x(t_1)=0.`$ The local maxima in $`E_{\mathrm{kin}}(t_1,t_0)`$ yield the classical prediction for the cutoffs in the harmonic spectra. The good agreement of the classical cutoff with the one found in the quantum spectra allows us to predict, with the classical model, the behavior of the cutoff as a function of the external parameters, i.e. the confinement constant $`\mathrm{\Omega }_\mathrm{h}`$, the frequency and the amplitude of the external field. We find that in the parameter range of interest the cutoff law can be written in the form $$\epsilon _{\mathrm{max}}=\|\epsilon _0\|+f(\mathrm{\Omega }_\mathrm{h},\omega )U_p,$$ (8) where $`f(\mathrm{\Omega }_\mathrm{h},\omega )`$ in general neither exhibits a simple functional form nor can be derived analytically. However, the linear dependence on the field intensity $`E_0^2`$ through $`U_p`$ in Eq. (8) is preserved just as in the case without confinement, see Fig. 3. Only for large confinement constants or electron trajectories with long excursion times $`f(\mathrm{\Omega }_\mathrm{h},\omega )`$ becomes slightly intensity-dependent. The general behavior of $`f(\mathrm{\Omega }_\mathrm{h},\omega )`$ is rather complex. Nevertheless, asymptotically a simple and familiar behavior is recovered: For very high frequency, the monochromatic cutoff constant is is approached, i.e., $`f(\mathrm{\Omega }_\mathrm{h},\omega \mathrm{})3.17`$, as can be seen in Fig. 4. For finite frequency $`\omega `$ the cutoff energy increases with growing $`\mathrm{\Omega }_\mathrm{h}`$. In fact, the lower the frequency, the more sensitively the cutoff law depends on $`\mathrm{\Omega }_\mathrm{h}`$. This property is the actual reason why one can obtain an increased cutoff energy with a confinement. For very low frequencies, the cutoff energy can be easily extended beyond $`\|\epsilon _0\|+9U_p.`$ In practice, however, there is a lower frequency limit to generate an appreciable intensity of high harmonics in the present context. If the confinement frequency is comparable to the laser frequency, $`\mathrm{\Omega }_\mathrm{h}\omega `$, the confinement potential itself starts to contribute to the harmonic generation process, ceasing to be a passive element. Hence, the condition for HHG under a confinement potential can be written as $`\mathrm{\Omega }_\mathrm{h}/\omega 1`$. However, there is also the usual upper limit in frequency $`\omega `$ which comes from the requirement that the atom in the laser field must be in the tunneling regime . Typical frequencies used in HHG experiments, and for which a long plateau is obtained, are in the vicinity of $`\omega =0.057\mathrm{a}.\mathrm{u}.`$. For this frequency a confinement indeed leads to a larger cutoff energy as demonstrated in Fig. 2. In conclusion, we have presented a new scheme for increasing the cutoff energy of the high-harmonic spectra of an atom under the influence of a strong laser field. Placing the atom in a confining parabolic potential, we have shown that the cutoff energy can be increased by more than fifty percent. An effective increase of the cutoff requires a careful choice of the confinement strength. The confinement curvature $`\mathrm{\Omega }_\mathrm{h}`$ must be strong enough for the electron to be appreciably accelerated towards the parent ion, but weak enough for it to move in a “quasi-continuum”. If $`\mathrm{\Omega }_\mathrm{h}`$ is too weak, the conventional cutoff law $`\|\epsilon _0\|+3.17U_p`$ is not altered by it. If $`\mathrm{\Omega }_\mathrm{h}`$ is too strong, the electron moves as a bound particle that does not generate higher harmonics. In the extreme case, one observes the dipole response of a harmonic oscillator, i.e., equally spaced resonances. A rough indication of whether the electron is in a “quasi-continuum” is given by the ratio of the energy difference between two consecutive levels of the confinement potential, $`\mathrm{\Delta }\epsilon _\mathrm{h}=\mathrm{\Omega }_\mathrm{h},`$ and the ionization potential of the atom in question. If $`\mathrm{\Omega }_\mathrm{h}/\|\epsilon _0\|1,`$ this condition is fulfilled. Also, as already discussed, the ratio between the frequency $`\omega `$ of the external field and the confinement curvature $`\mathrm{\Omega }_\mathrm{h}`$ plays an important role. If $`\mathrm{\Omega }_\mathrm{h}/\omega 1,`$ the parabolic potential contributes too actively to the harmonic generation process, and the plateau and cutoff are not present in the spectra. The best results have been obtained for $`x_0100\mathrm{a}.\mathrm{u}.,`$ $`\mathrm{\Omega }_\mathrm{h}0.02\mathrm{a}.\mathrm{u}.`$ and $`\omega 0.04\mathrm{a}.\mathrm{u}.`$ In this case, the energy difference between two consecutive levels of the confinement potential is still of the order of one tenth of the ionization potential $`\|\epsilon _0\|`$ and $`\mathrm{\Omega }_\mathrm{h}/\omega 0.5`$. For this parameter range, the cutoff energy can be extended until approximately $`\|\epsilon _0\|+6U_p`$. On a more technical level, yet very interesting from the theoretical point of view, we have seen that the cutoff law is given by the classical picture of an electron moving under the influence of the laser field and the confinement potential. Very good agreement between the quantum-mechanical full calculation and the classical model occurs for a wide range of field strengths, frequencies around $`\omega 0.05\mathrm{a}.\mathrm{u}.`$ and confinement curvatures of the order of $`\mathrm{\Omega }_\mathrm{h}10^2\mathrm{a}.\mathrm{u}.`$ Thereby we have found that the cutoff law strongly depends on the confinement curvature $`\mathrm{\Omega }_\mathrm{h}`$ and the frequency $`\omega `$ of the laser field, but only linearly on the field intensity $`E_0^2`$. The proposed setup presents several advantages over the schemes using static fields. For instance, using a confining potential, one can achieve a considerable extension of the cutoff energy already for the trajectories corresponding to short electron excursion times, whereas using static fields one mainly affects electron trajectories with long excursion times. Due to wave-packet spreading, the former trajectories are far more important for the harmonic spectra than the latter. In order to reduce the spreading one needs very strong magnetic fields . Another noteworthy feature of a confinement potential is that one can obtain stronger harmonics than in the static field, or even in the monochromatic case. In fact, a serious disadvantage concerning static electric fields is an appreciable decrease in the harmonic intensities compared to the field free case, due to depletion, i.e. irreversible ionization. This problem is not present in our scheme. However, similarly to the so far proposed extension of the cut off energy by using a combination of a static electric and magnetic fields, we are not aware of a direct possibility for an experimental realisation of our scheme. In the former case the magnetic field necessary is unrealistically large for a laboratory application . For our situation, a true electromagnetic trap is too macroscopic compared to the paramater range we need. On the other hand there might be exciting possibilities in the future to design a confined atom as described in a quantum-dot like device, for instance as an impurity. An important issue here, however, is the limitation in the radiation intensity in order to avoid the damage threshold. Recently, solid-state materials which can survive our parameter range, namely fields of wavelength $`\lambda =790\mathrm{nm}`$ and intensities above $`10^{14}\mathrm{W}/\mathrm{cm}^2`$, have been observed . Acknowledgements: We would like to thank K. Richter, D. B. Milošević, M. L. Du and K. Leo for useful discussions.
warning/0004/hep-th0004120.html	ar5iv	text	# 1 Introduction ## 1 Introduction We take spacetime $`\mathrm{\Sigma }`$ to be two dimensional Minkowski space. Let $`\sigma ^\pm =\tau \pm \sigma `$ be the standard lightcone coordinates and let $`\phi `$ be a massless free scalar field satisfying the wave equation $`_+^2\phi =0`$. The standard abelian duality transformations (for a review see ) are $`_+\stackrel{~}{\phi }`$ $`=`$ $`+_+\phi ,`$ (1.1) $`_{}\stackrel{~}{\phi }`$ $`=`$ $`_{}\phi .`$ (1.2) The integrability conditions for the above are precisely the equations of motion $`_+^2\phi =0`$. This means that we can construct $`\stackrel{~}{\phi }(\sigma ^+,\sigma ^{})`$ and verify that as a consequence of the above $`_+^2\stackrel{~}{\phi }=0`$. A generalization of the above are the pseudochiral models of the type introduced by Zakharov and Mikhailov . Consider a standard sigma model with target space a Lie group $`G`$ and with equations of motions $`^a(g^1_ag)=0`$. Introduce the dual Lie algebra valued field $`\varphi `$ by $$g^1_ag=ϵ_a{}_{}{}^{b}_{b}^{}\varphi .$$ It was shown by Nappi that these two descriptions were not quantum mechanically equivalent . The correct dual models were found by Fridling and Jevicki and by Fradkin and Tseytlin . It was eventually understood that the duality transformation should be a canonical transformation and the transformation in the pseudochiral model is not. Nevertheless the pseudochiral model has a variety of interesting field theoretic features . Motivated by string theory, there is now a vast literature on nonabelian duality and Poisson-Lie duality . Following up on recent work we consider a variant of the duality equations proposed there. These equations are a generalization of the ones of Zakharov and Mikhailov. In light of the introductory paragraph it will be interesting to study the physical and mathematical properties of these equations. Here we point out that there is an interesting transformation that maps solutions of the wave equation on a symmetric space into solutions of the wave equation on a symmetric space with the opposite curvature (see the next paragraph). We note that it was pointed out in that the opposite signs found by Nappi in the beta functions for the dual models of Zakharov and Mikhailov are explained by observing that the generalized curvatures in the models have opposite signs. The results presented here are a generalization<sup>3</sup><sup>3</sup>3I am thankful to E. Ivanov for bringing his work to my attention. of results of E. Ivanov . There is a large body of literature discussing conserved currents in sigma models based on groups or coset spaces. A seminal work was Pohlmeyer’s construction for an infinite number of conserved currents in sigma models with target space $`S^n`$. This was generalized by Eichenherr and Forger who showed that the construction generalized to symmetric spaces. Ivanov expanded on these ideas and in doing so introduced the notion the dual algebra and the dual sigma model. He states, > “… in which we show that the equations of any $`d=2`$ $`\sigma `$ model on a symmetric space simultaneously describe a $`d=2`$ $`\sigma `$ model on some other, dual factor space…” \[29, p. 475\]. He explicitly works out the example of a sigma model with target a compact real Lie group $`G`$ (with zero $`B`$-field). He views $`G`$ as a symmetric space $`G\times G/G`$ and he explicitly shows that the sigma model on $`G`$ is dual to the sigma model on $`G^{}/G`$ where $`G^{}`$ is the complexification of $`G`$. He subsequently asserts that the construction generalizes to a generic symmetric space. Ivanov’s construction for a symmetric space is given in Section 6. The work here generalizes Ivanov’s in that we assume a general riemannian manifold and show that it must be a symmetric space. We first establish some notation. The sigma model with target space $`M`$, metric $`g`$ and $`2`$-form $`B`$ will be denoted by $`(M,g,B)`$ and has lagrangian $$=\frac{1}{2}g_{ij}(x)\left(\frac{x^i}{\tau }\frac{x^j}{\tau }\frac{x^i}{\sigma }\frac{x^j}{\sigma }\right)+B_{ij}(x)\frac{x^i}{\tau }\frac{x^j}{\sigma }$$ (1.3) with canonical momentum density $$\pi _i=\frac{}{\dot{x}^i}=g_{ij}\dot{x}^j+B_{ij}x^j.$$ (1.4) The stress energy tensor for this sigma model is given by $`\mathrm{\Theta }_+=0`$, $$\mathrm{\Theta }_{++}=g_{ij}(x)_+x^i_+x^j\text{and}\mathrm{\Theta }_{}=g_{ij}(x)_{}x^i_{}x^j.$$ (1.5) In general a duality transformation between sigma models $`(M,g,B)`$ and $`(\stackrel{~}{M},\stackrel{~}{g},\stackrel{~}{B})`$ is a canonical transformation between the respective phase spaces that preserves the respective hamiltonian densities. We can study a less restrictive situation where we only have “on shell duality”. By this we mean that we only require that a map exists between solutions to the equations of motion of $`(M,g,B)`$ and solutions of the equations of motion of $`(\stackrel{~}{M},\stackrel{~}{g},\stackrel{~}{B})`$. The “on shell” transformation proposed below will be referred to as pseudoduality. There is a certain naturalness to the pseudoduality transformation because it preserves the stress energy tensor. As in it is convenient to choose an orthonormal frame $`\{\omega ^i\}`$ with the antisymmetric riemannian connection $`\omega _{ij}`$. The Cartan structural equations are $`d\omega ^i`$ $`=`$ $`\omega _{ij}\omega ^j,`$ $`d\omega _{ij}`$ $`=`$ $`\omega _{ik}\omega _{kj}+{\displaystyle \frac{1}{2}}R_{ijkl}\omega ^k\omega ^l.`$ We also define the curvature $`2`$-forms by $`\mathrm{\Omega }_{ij}=\frac{1}{2}R_{ijkl}\omega ^k\omega ^l`$. Consider two sigma models $`(M,g,B)`$ and $`(\stackrel{~}{M},\stackrel{~}{g},\stackrel{~}{B})`$. Using the orthonormal frame we define the appropriate $`\sigma ^a`$ derivative of the maps $`x:\mathrm{\Sigma }M`$ and $`\stackrel{~}{x}:\mathrm{\Sigma }\stackrel{~}{M}`$ by<sup>4</sup><sup>4</sup>4More correctly these equations are pullbacks of the type $`x^{}\omega ^i=x^i{}_{a}{}^{}d\sigma ^a`$. In the field of exterior differential systems the pullback is implicit and not usually written. We adhere to that convention. $$\omega ^i=x^i{}_{a}{}^{}d\sigma ^a\text{and}\stackrel{~}{\omega }^i=\stackrel{~}{x}^i{}_{a}{}^{}d\sigma ^a.$$ (1.6) In order to describe the equations of motion for the sigma model we need to consider second derivatives. The covariant derivatives of $`x^i_a`$ are $`x^i_{ab}`$ and are defined by $$dx^i{}_{a}{}^{}+\omega _{ij}x^j{}_{a}{}^{}=x^i{}_{ab}{}^{}d\sigma ^b.$$ (1.7) Taking the exterior derivative of $`\omega ^i=x^i{}_{a}{}^{}d\sigma ^a`$ we learn that $`x^i{}_{ab}{}^{}=x^i_{ba}`$. The equations of motion coming from (1.3) are $$x^k{}_{+}{}^{}=\frac{1}{2}H_{kij}x^i{}_{+}{}^{}x_{}^{j}{}_{}{}^{},$$ (1.8) where $`H=dB`$. There are similar definitions and equations for $`\stackrel{~}{x}^i`$. The lagrangian version of the hamiltonian duality equations in may be written as $`\stackrel{~}{x}_+(\sigma )`$ $`=`$ $`+T_+(x,\stackrel{~}{x})x_+(\sigma ),`$ (1.9) $`\stackrel{~}{x}_{}(\sigma )`$ $`=`$ $`T_{}(x,\stackrel{~}{x})x_{}(\sigma ).`$ (1.10) The orthogonal matrix valued functions $`T_\pm :M\times \stackrel{~}{M}\mathrm{SO}(n)`$ are not arbitrary but related by $$T_+(I+n)=T_{}(In),$$ (1.11) where the antisymmetric tensor $`n_{ij}`$ on $`M\times \stackrel{~}{M}`$ satisfies some PDEs given in . In this article we relax such restrictions on $`T_\pm `$ and consider orthogonal matrix valued functions $`T_\pm :\mathrm{\Sigma }\mathrm{SO}(n)`$ with the constraint that solutions to the the sigma model $`(M,g,B)`$ are mapped into solutions of $`(\stackrel{~}{M},\stackrel{~}{g},\stackrel{~}{B})`$ and vice versa. We will refer to these equations as the pseudoduality equations. Note that the pseudoduality transformations satisfy $`\stackrel{~}{\mathrm{\Theta }}_{\pm \pm }=\mathrm{\Theta }_{\pm \pm }`$. This is very different than in where $`T_\pm `$ are functions on $`M\times \stackrel{~}{M}`$. In other words, we allow explicit dependence on $`\sigma `$ in $`T`$ where as in the dependence occurs because $`x`$ and $`\stackrel{~}{x}`$ depend on $`\sigma `$. We specialize in this article to the case where $`T_+=T_{}=T:\mathrm{\Sigma }\mathrm{SO}(n)`$ and $`H=0`$, $`\stackrel{~}{H}=0`$. In this case the equations of motion become<sup>5</sup><sup>5</sup>5We remind the reader that solutions to Laplace’s equation are called harmonic functions yet solutions to the wave equation are simply called “waves”. $`x_+=\stackrel{~}{x}_+=0`$ and the duality equations are $$\stackrel{~}{x}_\pm (\sigma )=\pm T(\sigma )x_\pm (\sigma ).$$ (1.12) We point out that when $`H=\stackrel{~}{H}=0`$ the sigma models are worldsheet parity and worldsheet time reversal invariant. The isometry $`T`$ is odd under worldsheet parity and under worldsheet time reversal. Note that in the pseudochiral model $`\stackrel{~}{H}0`$ and therefore it is not covered in this paper. The general case with generic $`H`$ and $`\stackrel{~}{H}`$ will be treated elsewhere . Duality equations (1.12) are mathematically quite interesting. We digress a bit and discuss the question of local riemannian isometries. Assume we have a map $`f:M\stackrel{~}{M}`$ and we wish for this map to be an isometry. If $`\{\omega ^j\}`$ and $`\{\stackrel{~}{\omega }^i\}`$ are local orthonormal frames then we require that the metric be preserved: $`f^{}(\stackrel{~}{\omega }^i\stackrel{~}{\omega }^i)=\omega ^j\omega ^j`$. The solution to this equation is $`f^{}\stackrel{~}{\omega }=T\omega `$ where $`T:M\mathrm{O}(n)`$. This Pfaffian system of equations is integrable if the Riemann curvature tensor and its higher order covariant derivatives in the two spaces agree when identified via $`T`$. For more details see the discussion in the paragraph following (3.16). In our case we begin with maps $`x:\mathrm{\Sigma }M`$, $`\stackrel{~}{x}:\mathrm{\Sigma }\stackrel{~}{M}`$ that satisfy hyperbolic Lorentz invariant equations. We fix a Lorentz frame and we observe that we are interested in maps from $`\{x:\mathrm{\Sigma }M\|x_+=0\}`$ into $`\{\stackrel{~}{x}:\mathrm{\Sigma }\stackrel{~}{M}\|\stackrel{~}{x}_+=0\}`$ that preserve the two independent components of the energy momentum tensor $`\mathrm{\Theta }_{++}`$ and $`\mathrm{\Theta }_{}`$. A linearity assumption and (1.5) tells us that the maps have to be of form (1.9) and (1.10) with the more general $`T_\pm (\sigma )`$. Given a map $`x:\mathrm{\Sigma }M`$ there are two preferred tangent vector fields $`/\sigma ^+`$ and $`/\sigma ^{}`$. The conditions that require $`\mathrm{\Theta }`$ to be preserved are very geometric: the lengths of $`/\sigma ^+`$ and $`/\sigma ^{}`$ are preserved by the map. If $`_\mathrm{\Sigma }`$ denotes the Hodge duality operation on $`\mathrm{\Sigma }`$, our equations may be written as $`\stackrel{~}{\omega }=_\mathrm{\Sigma }(T\omega )`$ where we interpret $`\omega `$ and $`\stackrel{~}{\omega }`$ as pull backs to $`\mathrm{\Sigma }`$. It is the desire to have the $`_\mathrm{\Sigma }`$ operation that introduces a $`1`$ in (1.10) even though in principle the $`1`$ could be absorbed<sup>6</sup><sup>6</sup>6We may have to allow $`T_\pm \mathrm{O}(n)`$. into $`T_{}`$. Here we show that there are interesting maps between “waves” on $`M`$ and “waves” on $`\stackrel{~}{M}`$ that preserve natural geometrical structures. This may be of interest to researchers who study two dimensional wave equations. This paper is organized in the following way. In Section 2 we discuss pseudoduality between the $`2`$-sphere $`S^2`$ and the $`2`$-dimensional hyperbolic space $`H^2`$ very concretely. In particular we explicitly verify the necessity for the the orthogonal matrix valued function $`T`$. In Section 3 we become a bit more abstract but still concrete by working with explicit metrics and show that there is pseudoduality between $`S^n`$ and $`H^n`$. In Section 4 we assume general metrics and show that the manifolds $`M`$ and $`\stackrel{~}{M}`$ must be symmetric spaces with the “opposite curvatures”. In Section 5 we construct many examples by discussing the theory of dual symmetric spaces. Section 7 is a discussion of the results of this article. ## 2 A Pedagogic Example It is worthwhile to be very concrete and to consider the pseudoduality between strings moving on a $`2`$-sphere $`S^2`$ and those moving on a $`2`$-hyperboloid $`H^2`$. This example illustrates the necessity for the matrix $`T`$. We use a different coordinate version of the constant curvature metric than in Section 3 to emphasize that everything is independent of the choice of coordinates. The respective constant curvature metrics on $`S^2`$ and $`H^2`$ in polar normal coordinates are $`ds^2`$ $`=`$ $`dr^2+{\displaystyle \frac{\mathrm{sin}^2\alpha r}{\alpha ^2}}d\theta ^2,`$ $`d\stackrel{~}{s}^2`$ $`=`$ $`d\stackrel{~}{r}^2+{\displaystyle \frac{\mathrm{sinh}^2\stackrel{~}{\alpha }\stackrel{~}{r}}{\stackrel{~}{\alpha }^2}}d\stackrel{~}{\theta }^2.`$ The equations of motion for the sigma models are: $`_+^2r`$ $`=`$ $`(\mathrm{sin}\alpha r\mathrm{cos}\alpha r)/\alpha _+\theta _{}\theta ,`$ (2.1) $`(\mathrm{sin}\alpha r)/\alpha _+^2\theta `$ $`=`$ $`\mathrm{cos}\alpha r(_+r_{}\theta +_{}r_+\theta ),`$ (2.2) $`_+^2\stackrel{~}{r}`$ $`=`$ $`(\mathrm{sinh}\stackrel{~}{\alpha }\stackrel{~}{r}\mathrm{cosh}\stackrel{~}{\alpha }\stackrel{~}{r})/\stackrel{~}{\alpha }_+\stackrel{~}{\theta }_{}\stackrel{~}{\theta },`$ (2.3) $`(\mathrm{sinh}\stackrel{~}{\alpha }\stackrel{~}{r})/\stackrel{~}{\alpha }_+^2\stackrel{~}{\theta }`$ $`=`$ $`\mathrm{cosh}\stackrel{~}{\alpha }\stackrel{~}{r}(_+\stackrel{~}{r}_{}\stackrel{~}{\theta }+_{}\stackrel{~}{r}_+\stackrel{~}{\theta }).`$ (2.4) This elementary example illustrates the importance of the matrix $`T`$. Assume $`T=I`$ then two of the duality equations in this coordinate system would be $`_+\stackrel{~}{r}=_+r`$ and $`_{}\stackrel{~}{r}=_{}r`$. This would imply that $`_+^2r=0`$ and $`_+^2\stackrel{~}{r}=0`$ but these are not the equations of motion. We will see that with a $`TI`$ we can construct a pseudoduality transformation. First we need an orthonormal frame. On $`S^2`$ we have $`\omega ^r=dr`$ and $`\omega ^\theta =\mathrm{sin}\alpha r/\alpha d\theta `$. From this it follows that the connection and the curvature are respectively given by $`\omega _{r\theta }=\mathrm{cos}(\alpha r)d\theta `$ and $`d\omega _{r\theta }=\alpha ^2\omega ^r\omega ^\theta `$. On $`H^2`$ we have $`\stackrel{~}{\omega }^{\stackrel{~}{r}}=d\stackrel{~}{r}`$ and $`\stackrel{~}{\omega }^{\stackrel{~}{\theta }}=\mathrm{sinh}\alpha \stackrel{~}{r}/\alpha d\stackrel{~}{\theta }`$. From this we see that the connection and the curvature are respectively given by $`\stackrel{~}{\omega }_{\stackrel{~}{r}\stackrel{~}{\theta }}=\mathrm{cosh}(\stackrel{~}{\alpha }\stackrel{~}{r})d\stackrel{~}{\theta }`$ and $`d\stackrel{~}{\omega }_{\stackrel{~}{r}\stackrel{~}{\theta }}=\stackrel{~}{\alpha }^2\stackrel{~}{\omega }^{\stackrel{~}{r}}\stackrel{~}{\omega }^{\stackrel{~}{\theta }}`$. We need a $`2\times 2`$ orthogonal matrix that we parametrize by a rotation angle $`\varphi `$. The duality equations are $$\left(\begin{array}{c}_\pm \stackrel{~}{r}\\ \mathrm{sinh}(\stackrel{~}{\alpha }\stackrel{~}{r})/\stackrel{~}{\alpha }_\pm \stackrel{~}{\theta }\end{array}\right)=\pm \left(\begin{array}{cc}\hfill \mathrm{cos}\varphi & \hfill \mathrm{sin}\varphi \\ \hfill \mathrm{sin}\varphi & \hfill \mathrm{cos}\varphi \end{array}\right)\left(\begin{array}{c}_\pm r\\ \mathrm{sin}(\alpha r)/\alpha _\pm \theta \end{array}\right).$$ (2.5) The integrability conditions on the above lead to the equation $$d\varphi =\mathrm{cosh}(\stackrel{~}{\alpha }\stackrel{~}{r})d\stackrel{~}{\theta }+\mathrm{cos}(\alpha r)d\theta .$$ (2.6) This is integrable if $`\alpha =\stackrel{~}{\alpha }`$ where integrable means integrable modulo the equations of motion. What happens to the point particle geodesics on $`S^2`$? One easily verifies that for constant $`a`$, the functions $`r=a(\sigma ^++\sigma ^{})=a(2\tau )`$ and $`\theta =0`$ give a solution of equations (2.1) and (2.2). This corresponds to “particle geodesics” on $`S^2`$. We would like to understand what are the dual solutions to the particle geodesics. If we note that $`_\pm r=a`$ and $`_\pm \theta =0`$ we see that the duality equations become $`_\pm \stackrel{~}{r}`$ $`=`$ $`\pm a\mathrm{cos}\varphi `$ (2.7) $`\mathrm{sinh}(\stackrel{~}{\alpha }\stackrel{~}{r})/\stackrel{~}{\alpha }_\pm \stackrel{~}{\theta }`$ $`=`$ $`a\mathrm{sin}\varphi `$ (2.8) An immediate consequence of the above is that $`(_++_{})\stackrel{~}{r}=0`$ and $`(_++_{})\stackrel{~}{\theta }=0`$. We conclude that $`\stackrel{~}{r}=\stackrel{~}{r}(\sigma ^+\sigma ^{})=\stackrel{~}{r}(2\sigma )`$ and $`\stackrel{~}{\theta }=\stackrel{~}{\theta }(\sigma ^+\sigma ^{})=\stackrel{~}{\theta }(2\sigma )`$. Thus we get static solutions on the hyperboloid. Note that (2.6) immediately tells us that $`\varphi =\varphi (\sigma ^+\sigma ^{})=\varphi (2\sigma )`$. Since everything is a function of $`2\sigma =\sigma ^+\sigma ^{}`$ the situation reduces to functions of a single variable. The transformed solutions will be static solutions ($`\tau `$ independent). Note that equations (2.7) and (2.8) lead to $$(_+\stackrel{~}{r})^2+\left(\frac{\mathrm{sinh}\stackrel{~}{\alpha }\stackrel{~}{r}}{\stackrel{~}{\alpha }}_+\stackrel{~}{\theta }\right)^2=a^2.$$ (2.9) This is the “conservation of energy” equation associated with the particle lagrangian $$L=\frac{1}{2}\left[(_+\stackrel{~}{r})^2+\left(\frac{\mathrm{sinh}\stackrel{~}{\alpha }\stackrel{~}{r}}{\stackrel{~}{\alpha }}_+\stackrel{~}{\theta }\right)^2\right]$$ This is a standard problem in classical mechanics. The canonical momentum $$\stackrel{~}{J}=\frac{L}{(_+\stackrel{~}{\theta })}=\left(\frac{\mathrm{sinh}\stackrel{~}{\alpha }\stackrel{~}{r}}{\stackrel{~}{\alpha }}\right)^2_+\stackrel{~}{\theta }$$ is a constant of the motion. Thus the energy integral (2.9) may be written as $$(_+\stackrel{~}{r})^2+\frac{\stackrel{~}{\alpha }^2\stackrel{~}{J}^2}{\mathrm{sinh}^2\stackrel{~}{\alpha }\stackrel{~}{r}}=a^2$$ (2.10) and this problem is reducible to quadrature. This is interesting because we know that there exists crystallographic subgroups $`\mathrm{\Gamma }`$ of $`\mathrm{SL}(2,)`$ such that $`H/\mathrm{\Gamma }`$ is a genus $`g>1`$ compact Riemann surface. Such a surface has closed geodesics of minimal length and some of these are the static “soliton-like” solutions $`(\stackrel{~}{r}(2\sigma ),\stackrel{~}{\theta }(2\sigma ))`$ we are constructing above. ## 3 Constant Curvature Metric Before presenting the general theory we study the special case of constant curvature spaces. Again we take $`T_+=T_{}=T`$. The two dimensional nonlinear sigma model on a space with constant positive curvature $`k>0`$ is shown to be pseudodual to the nonlinear sigma model on a space with constant negative curvature $`k`$. The pseudoduality equations are used to map solutions of one model into solutions of the other model. Locally, a metric on a space with constant curvature is may be written in the Poincaré form $$ds^2=\frac{dx^idx^i}{\left(1+kx^2\right)^2}.$$ (3.1) An orthonormal frame in this metric is $$\omega ^i=\frac{dx^i}{1+kx^2}.$$ (3.2) The connection one-forms (with respect to the orthonormal frame) are $$\omega _{ij}=2k\frac{x^idx^jx^jdx^i}{1+kx^2}.$$ (3.3) The first Cartan structural equation is $`d\omega ^i=\omega _{ij}\omega ^j`$. The second structural equation gives the the curvature two forms $$\mathrm{\Omega }_{ij}=d\omega _{ij}+\omega _{ik}\omega _{kj}=\frac{1}{2}R_{ijkl}\omega ^k\omega ^l=4k\omega ^i\omega ^j.$$ (3.4) Technically the curvature is $`4k`$. The lagrangian for the sigma model is $$=\frac{2_+x^i_{}x^i}{\left(1+kx^2\right)^2}.$$ (3.5) The equations of motion associated with this lagrangian are $$_+^2x^i=2k\frac{x^j_+x^j_{}x^i+x^j_{}x^j_+x^ix^i_+x^j_{}x^j}{1+kx^2}.$$ (3.6) Assume we have two constant curvature spaces $`M`$ and $`\stackrel{~}{M}`$ with respective curvatures $`k`$ and $`\stackrel{~}{k}`$. We would like to see if it is possible to have a duality transformation between them. Assume we have a solution $`x(\sigma )`$ of the sigma model on $`M`$. We attempt to construct a solution of the sigma model on $`\stackrel{~}{M}`$ by requiring that $`{\displaystyle \frac{_+\stackrel{~}{x}^i}{1+\stackrel{~}{k}\stackrel{~}{x}^2}}`$ $`=`$ $`T^i{}_{j}{}^{}{\displaystyle \frac{_+x^j}{1+kx^2}},`$ (3.7) $`{\displaystyle \frac{_{}\stackrel{~}{x}^i}{1+\stackrel{~}{k}\stackrel{~}{x}^2}}`$ $`=`$ $`T^i{}_{j}{}^{}{\displaystyle \frac{_{}x^j}{1+kx^2}},`$ (3.8) where $`T`$ is an orthogonal matrix, $`detT=1`$. Note that the stress energy tensors satisfy $`\mathrm{\Theta }_{\pm \pm }=\stackrel{~}{\mathrm{\Theta }}_{\pm \pm }`$. The first step towards the integrability conditions for the above system is to differentiate (3.7) with respect to $`\sigma ^{}`$ and (3.8) with respect to $`\sigma ^+`$: $`{\displaystyle \frac{_{}(_+\stackrel{~}{x}^i)}{1+\stackrel{~}{k}\stackrel{~}{x}^2}}2\stackrel{~}{k}{\displaystyle \frac{\stackrel{~}{x}^j_{}\stackrel{~}{x}^j_+\stackrel{~}{x}^i}{(1+\stackrel{~}{k}\stackrel{~}{x}^2)^2}}`$ $`=`$ $`(_{}T^i{}_{j}{}^{}){\displaystyle \frac{_+x^j}{1+kx^2}}`$ (3.9) $`+`$ $`T^i{}_{j}{}^{}{\displaystyle \frac{_+^2x^j}{1+kx^2}}2kT^i{}_{j}{}^{}{\displaystyle \frac{x^k_{}x^k_+x^j}{(1+kx^2)^2}},`$ $`{\displaystyle \frac{_+(_{}\stackrel{~}{x}^i)}{1+\stackrel{~}{k}\stackrel{~}{x}^2}}+2\stackrel{~}{k}{\displaystyle \frac{\stackrel{~}{x}^j_+\stackrel{~}{x}^j_{}\stackrel{~}{x}^i}{(1+\stackrel{~}{k}\stackrel{~}{x}^2)^2}}`$ $`=`$ $`(_+T^i{}_{j}{}^{}){\displaystyle \frac{_{}x^j}{1+kx^2}}`$ (3.10) $`+`$ $`T^i{}_{j}{}^{}{\displaystyle \frac{_+^2x^j}{1+kx^2}}2kT^i{}_{j}{}^{}{\displaystyle \frac{x^k_+x^k_{}x^j}{(1+kx^2)^2}}.`$ By imposing the integrability conditions $`_+(_{}\stackrel{~}{x}^i)=_{}(_+\stackrel{~}{x}^i)`$ we can add the above two equations to eliminate $`_+^2\stackrel{~}{x}`$ and by imposing only the equations of motion for $`x^i`$ and the duality relations we obtain $`(_{}T^i{}_{j}{}^{}){\displaystyle \frac{_+x^j}{1+kx^2}}+(_+T^i{}_{j}{}^{}){\displaystyle \frac{_{}x^j}{1+kx^2}}`$ $`=`$ $`2\stackrel{~}{k}{\displaystyle \frac{\stackrel{~}{x}^j_{}\stackrel{~}{x}^i\stackrel{~}{x}^i_{}\stackrel{~}{x}^j}{1+\stackrel{~}{k}\stackrel{~}{x}^2}}T^j{}_{k}{}^{}{\displaystyle \frac{_+x^k}{1+kx^2}}`$ (3.11) $`+`$ $`2\stackrel{~}{k}{\displaystyle \frac{\stackrel{~}{x}^j_+\stackrel{~}{x}^i\stackrel{~}{x}^i_+\stackrel{~}{x}^j}{1+\stackrel{~}{k}\stackrel{~}{x}^2}}T^j{}_{k}{}^{}{\displaystyle \frac{_{}x^k}{1+kx^2}}`$ $``$ $`2kT^i{}_{j}{}^{}{\displaystyle \frac{x^k_{}x^jx^j_{}x^k}{1+kx^2}}{\displaystyle \frac{_+x^k}{1+kx^2}}`$ $``$ $`2kT^i{}_{j}{}^{}{\displaystyle \frac{x^k_+x^jx^j_+x^k}{1+kx^2}}{\displaystyle \frac{_{}x^k}{1+kx^2}}.`$ The particular combinations in the above follow from the general theory and the form of (3.3). Note that at any point $`\sigma `$ we can make $`_\pm x^i(\sigma )`$ arbitrary thus we need that $`_{}T^i_k`$ $`=`$ $`2\stackrel{~}{k}{\displaystyle \frac{\stackrel{~}{x}^j_{}\stackrel{~}{x}^i\stackrel{~}{x}^i_{}\stackrel{~}{x}^j}{1+\stackrel{~}{k}\stackrel{~}{x}^2}}T^j{}_{k}{}^{}2kT^i{}_{j}{}^{}{\displaystyle \frac{x^k_{}x^jx^j_{}x^k}{1+kx^2}},`$ (3.12) $`_+T^i_k`$ $`=`$ $`2\stackrel{~}{k}{\displaystyle \frac{\stackrel{~}{x}^j_+\stackrel{~}{x}^i\stackrel{~}{x}^i_+\stackrel{~}{x}^j}{1+\stackrel{~}{k}\stackrel{~}{x}^2}}T^j{}_{k}{}^{}2kT^i{}_{j}{}^{}{\displaystyle \frac{x^k_+x^jx^j_+x^k}{1+kx^2}}.`$ (3.13) We interpret the above as the pullback under the map $`\sigma (x(\sigma ),\stackrel{~}{x}(\sigma ))`$ of the $`1`$-form $$dT^i{}_{j}{}^{}=\stackrel{~}{\omega }_{ik}T^k{}_{j}{}^{}+T^i{}_{k}{}^{}\omega _{kj}^{}.$$ (3.14) We can immediately verify that if one inserts (3.12) into equation (3.9) we get the equations for motion for $`\stackrel{~}{x}`$, i.e., equations (3.6) with all quantities with tildes. Thus far we have shown that if one starts with a solution of the sigma model defined by lagrangian (3.5) and we impose the generalized duality equations then we find a solution to the “tilde” sigma model provided that we can satisfy equations (3.13) and (3.12). The integrability conditions for the $`dT`$ equation are found by taking the exterior derivative leading to $$0=\stackrel{~}{\mathrm{\Omega }}_{ik}T^k{}_{j}{}^{}+T^i{}_{k}{}^{}\mathrm{\Omega }_{kj}^{}$$ (3.15) In our case the curvature is very simple and the above reduces to $$4\stackrel{~}{k}\stackrel{~}{\omega }^i\stackrel{~}{\omega }^kT^k{}_{j}{}^{}+4kT^i{}_{k}{}^{}\omega _{}^{k}\omega ^j=0.$$ (3.16) First we discuss what equation (3.15) is not and afterwards we discuss what it is. This equation shows up in the study of local riemannian isometries . If we have two riemannian manifolds $`M`$ and $`\stackrel{~}{M}`$ and a map $`f:M\stackrel{~}{M}`$ then the conditions that $`f`$ be a local isometry is that there exists an orthogonal transformation $`T`$ such that $`\stackrel{~}{\omega }^i=T^i{}_{j}{}^{}\omega _{}^{j}`$. When one works out the integrability conditions for this system one finds that (3.14) must be satisfied along with its integrability condition which is (3.15). This tells us that $`\stackrel{~}{R}_{ijkl}T^k{}_{m}{}^{}T_{}^{l}{}_{n}{}^{}=T^i{}_{k}{}^{}T_{}^{j}{}_{l}{}^{}R_{klmn}^{}`$. There are further integrability conditions which tell us that the covariant derivatives of the curvatures also satisfy similar relations. This is the classical theorem of Christoffel on local isometries between riemannian manifolds as reformulated by E. Cartan. In our constant curvature example the requirement simply becomes $`k=\stackrel{~}{k}`$. Our situation is very different. In effect our setup<sup>7</sup><sup>7</sup>7We avoid discussing jet bundles. is a map from $`\mathrm{\Sigma }`$ to $`M\times \stackrel{~}{M}`$. This tells us that $`\omega ^i={\displaystyle \frac{(_ax^i)d\sigma ^a}{1+kx^2}}\text{and}\stackrel{~}{\omega }^i={\displaystyle \frac{(_a\stackrel{~}{x}^i)d\sigma ^a}{1+\stackrel{~}{k}\stackrel{~}{x}^2}}.`$ Substituting these into (3.16), using duality relations (3.7) and (3.8), and noting that at any $`(\sigma ^+,\sigma ^{})`$ we are free to arbitrarily specify $`_\pm x^i(\sigma )`$ leads to $`k=\stackrel{~}{k}`$. The change in sign is due to the negative sign in (3.8). Thus we discover that the duality equations can be implemented if the constant curvature manifolds have the opposite curvature. ## 4 General Theory This problem is best analyzed in the bundle of orthogonal frames . Because most physicists are not familiar with this approach we work things out on the base manifold $`M\times \stackrel{~}{M}`$. First thing to do is to take the exterior derivative of (1.12) $$d\stackrel{~}{x}_\pm =\pm (dT)x_\pm \pm Tdx_\pm $$ and use the definitions of the covariant derivative (1.7) to obtain $$\stackrel{~}{\omega }\stackrel{~}{x}_\pm +\stackrel{~}{x}_{\pm a}d\sigma ^a=\pm (dT)x_\pm T\omega x_\pm \pm Tx_{\pm a}d\sigma ^a.$$ If we use the duality equations (1.12) we have $$\stackrel{~}{\omega }Tx_\pm +\stackrel{~}{x}_{\pm a}d\sigma ^a=\pm (dT)x_\pm T\omega x_\pm \pm Tx_{\pm a}d\sigma ^a.$$ A little algebra shows that $$\stackrel{~}{x}_{\pm a}d\sigma ^a=\pm (dTT\omega +\stackrel{~}{\omega }T)x_\pm \pm Tx_{\pm a}d\sigma ^a.$$ (4.1) We wish to isolate the integrability conditions so wedge the above with $`d\sigma ^\pm `$. $$\stackrel{~}{x}_\pm d\sigma ^{}d\sigma ^\pm =\pm (dTT\omega +\stackrel{~}{\omega }T)x_\pm d\sigma ^\pm \pm Tx_\pm d\sigma ^{}d\sigma ^\pm .$$ We have two equations $`\stackrel{~}{x}_+d\sigma ^{}d\sigma ^+`$ $`=`$ $`+(dTT\omega +\stackrel{~}{\omega }T)x_+d\sigma ^++Tx_+d\sigma ^{}d\sigma ^+,`$ $`\stackrel{~}{x}_+d\sigma ^+d\sigma ^{}`$ $`=`$ $`(dTT\omega +\stackrel{~}{\omega }T)x_{}d\sigma ^{}Tx_+d\sigma ^+d\sigma ^{}.`$ In principle we wish that the integrability conditions $`\stackrel{~}{x}_+=\stackrel{~}{x}_+`$ are satisfied if the equations of motion $`x_+=0`$ hold. Subsequently we would like that this implies that $`\stackrel{~}{x}_+=0`$. We might as well substitute $`x_+=0`$ and $`\stackrel{~}{x}_+=0`$ into the equations above and find $$0=(dTT\omega +\stackrel{~}{\omega }T)x_\pm d\sigma ^\pm .$$ Since $`x_\pm ^i`$ may be arbitrarily specified at any $`\sigma `$ we have $$0=(dTT\omega +\stackrel{~}{\omega }T)d\sigma ^\pm .$$ Since $`x:\mathrm{\Sigma }M`$ and $`\stackrel{~}{x}:\mathrm{\Sigma }\stackrel{~}{M}`$ are maps of a two dimensional worldsheet we conclude that on the worldsheet, the covariant derivative of $`T`$ vanishes $$dTT\omega +\stackrel{~}{\omega }T=0.$$ (4.2) The reason is that the covariant differential of $`T`$ is a $`1`$-form. All our objects arise from maps with domain $`\mathrm{\Sigma }`$ so the covariant differential of $`T`$ is a $`1`$-form on $`\mathrm{\Sigma }`$. You are pulling back both $`\omega `$ and $`\stackrel{~}{\omega }`$ to $`\mathrm{\Sigma }`$. Note that a covariantly constant tensor is determined its value at one point on the worldsheet; the values elsewhere are determined by parallel transport. In order to construct such a $`T`$ we need to verify the integrability conditions for the above. These lead to important constraints on the geometry of $`M`$ and $`\stackrel{~}{M}`$. Taking the exterior derivative and using the Cartan structural equations leads to $$T^i{}_{k}{}^{}\mathrm{\Omega }_{kj}^{}+\stackrel{~}{\mathrm{\Omega }}_{ik}T^k{}_{j}{}^{}=0$$ with special case (3.15). Expanding the above gives $$\frac{1}{2}T^i{}_{k}{}^{}R_{kjlm}^{}\omega ^l\omega ^m+\stackrel{~}{R}_{iklm}T^k{}_{j}{}^{}\stackrel{~}{\omega }_{}^{l}\stackrel{~}{\omega }^m=0.$$ If we now substitute (1.6) and use the pseudoduality equations (1.12) we see that $$T^i{}_{k}{}^{}T_{}^{j}{}_{l}{}^{}R_{klmn}^{}=\stackrel{~}{R}_{ijkl}T^k{}_{m}{}^{}T_{}^{l}{}_{n}{}^{}.$$ (4.3) Thus we conclude that the manifolds $`M`$ and $`\stackrel{~}{M}`$ “have the opposite curvature”. Next we take the exterior derivative of the above to look for further conditions. If the covariant differential of $`R`$ is defined by $`DR_{ijkl}=R_{ijkl;m}\omega ^m`$ and similarly for $`\stackrel{~}{R}`$ we find $$T^i{}_{k}{}^{}T_{}^{j}{}_{l}{}^{}R_{klmn;p}^{}\omega ^p=\stackrel{~}{R}_{ijkl;p}T^k{}_{m}{}^{}T_{}^{l}{}_{n}{}^{}\stackrel{~}{\omega }_{}^{p}.$$ If we now substitute (1.6) into the above and use the pseudoduality equations (1.12) we obtain two equations since $`d\sigma ^+`$ and $`d\sigma ^{}`$ are independent: $`T^i{}_{k}{}^{}T_{}^{j}{}_{l}{}^{}R_{klmn;q}^{}`$ $`=`$ $`\stackrel{~}{R}_{ijkl;p}T^k{}_{m}{}^{}T_{}^{l}{}_{n}{}^{}T_{}^{p}{}_{q}{}^{},`$ $`T^i{}_{k}{}^{}T_{}^{j}{}_{l}{}^{}R_{klmn;q}^{}`$ $`=`$ $`+\stackrel{~}{R}_{ijkl;p}T^k{}_{m}{}^{}T_{}^{l}{}_{n}{}^{}T_{}^{p}{}_{q}{}^{}.`$ The solution of the above is immediate $$R_{klmn;q}=0\text{and}\stackrel{~}{R}_{ijkl;p}=0.$$ (4.4) The manifolds $`M`$ and $`\stackrel{~}{M}`$ must be locally symmetric spaces with the “opposite curvature”. ## 5 Dual Symmetric Spaces There is a large class of examples of pairs of symmetric spaces with opposite curvature. These pairs are called dual symmetric spaces. The simplest pair is the $`n`$-sphere $`S^n`$ and the $`n`$-dimensional hyperbolic space $`H^n`$. More complicated pairs are given by the Grassmann manifolds<sup>8</sup><sup>8</sup>8$`\mathrm{O}_0(p,q)`$ is the component of $`\mathrm{O}(p,q)`$ connected to the identity. $`\mathrm{SO}(p+q)/\mathrm{SO}(p)\times \mathrm{SO}(q)`$ and $`\mathrm{O}_0(p,q)/\mathrm{SO}(p)\times \mathrm{SO}(q)`$. The general theory of symmetric spaces is very extensive . For our purposes we take a lighter approach \[37, Chapter 11\] and a more restrictive view and consider what are called normal symmetric spaces. These are specified by a triplet of data $`(G/H,\sigma ,Q)`$ where $`G`$ is a real Lie group, $`H`$ is a closed subgroup of $`G`$ and $`\sigma `$ is an involutive automorphism of $`G`$. Let $`𝔤`$ be the Lie algebra of $`G`$ and denote the action of the automorphism $`\sigma `$ on $`𝔤`$ by $`s`$. Let $`𝔥`$ be the $`+1`$ eigenspace of $`s`$ and let $`𝔪`$ be the $`1`$ eigenspace of $`s`$. Since $`\frac{1}{2}(I\pm s)`$ are projectors we have $`𝔤=𝔥𝔪`$. The following are also required in a normal symmetric space: 1. If $`F`$ is the fixed point set of $`\sigma `$ and $`F_0`$ is its identity component then $`F_0HF`$. This is a technical requirement that for our purposes we take it to mean that the Lie algebra of $`H`$ is $`𝔥`$. 2. $`Q`$ is an $`\mathrm{Ad}(G)`$-invariant inner product on $`𝔤`$. An ordinary symmetric space only requires $`Q`$ to be an $`\mathrm{Ad}(H)`$-invariant inner product<sup>9</sup><sup>9</sup>9An $`\mathrm{Ad}(H)`$-invariant inner product on $`𝔪`$ leads to a $`G`$-invariant metric on $`M`$, see \[37, p. 312\]. on $`𝔪`$. Here $`\mathrm{Ad}(G)`$ is the adjoint action of the group $`G`$ on its Lie algebra $`𝔤`$. 3. $`Q`$ is $`s`$-invariant. This is not required for an ordinary symmetric space. It follows that $`[𝔥,𝔥]𝔥`$, $`[𝔥,𝔪]𝔪`$ and $`[𝔪,𝔪]𝔥`$. The $`s`$-invariance of $`Q`$ tells us that the direct sum decomposition is an orthogonal decomposition. Note that $`𝔥`$ and $`𝔪`$ are orthogonal with respect to any $`s`$-invariant quadratic form such as the Killing form $`\mathrm{Tr}(\mathrm{ad}(X)\mathrm{ad}(Y))`$. We pick an origin for the symmetric space $`G/H`$ and associate $`𝔪`$ with the tangent space at that point. For a normal symmetric space the sectional curvature associated to the $`2`$-plane spanned by $`X,Y𝔪`$ is given by the following simple formula \[37, p. 319\] $$K(X,Y)=\frac{Q([X,Y],[X,Y])}{Q(X,X)Q(Y,Y)Q(X,Y)^2}$$ (5.1) that requires the $`\mathrm{Ad}(G)`$-invariance of $`Q`$ in its derivation. We remind the reader that knowing all sectional curvatures is equivalent to knowing the curvature tensor and that they are related by $$K(X,Y)=\frac{R_{ijkl}X^iY^jX^kY^l}{Q(X,X)Q(Y,Y)Q(X,Y)^2}.$$ Normal symmetric spaces $`(G/H,\sigma ,Q)`$ and $`(\stackrel{~}{G}/\stackrel{~}{H},\stackrel{~}{\sigma },\stackrel{~}{Q})`$ are said to be dual symmetric spaces if there exist 1. a Lie algebra isomorphism $`S:𝔥\stackrel{~}{𝔥}`$ such that $`\stackrel{~}{Q}(SV,SW)=Q(V,W)`$ for all $`V,W𝔥`$. 2. a linear isometry $`T:𝔪\stackrel{~}{𝔪}`$ such that $`[TX,TY]=S[X,Y]`$ for all $`X,Y𝔪`$. In item (1) above the Lie algebra isomorphism tells us that brackets in $`𝔥`$ are the same as in $`\stackrel{~}{𝔥}`$. While the isometry in item (2) tells us that inner product on $`𝔪`$ is the same as in $`\stackrel{~}{𝔪}`$. For dual symmetric spaces it is easy to see that the sectional curvatures are related by $`\stackrel{~}{K}(TX,TY)=K(X,Y)`$. It is worthwhile to work this out explicitly for the example of the dual symmetric spaces $`\mathrm{SO}(p+q)/\mathrm{SO}(p)\times \mathrm{SO}(q)`$ and $`\mathrm{O}_0(p,q)/\mathrm{SO}(p)\times \mathrm{SO}(q)`$. For $`gG=\mathrm{SO}(p+q)`$ or $`g\stackrel{~}{G}=\mathrm{O}_0(p,q)`$ we take $`\sigma `$ and $`\stackrel{~}{\sigma }`$ to be given by $$g\left(\begin{array}{cc}I_p& 0\\ 0& I_q\end{array}\right)g\left(\begin{array}{cc}I_p& 0\\ 0& I_q\end{array}\right).$$ One easily verifies that $`𝔥=\stackrel{~}{𝔥}=𝔰𝔬(p)𝔰𝔬(q)`$ where the matrices are of the form $$\left(\begin{array}{cc}a& 0\\ 0& c\end{array}\right)$$ where $`a𝔰𝔬(p)`$ and $`c𝔰𝔬(q)`$. One also sees that $`X𝔪`$ and $`\stackrel{~}{X}\stackrel{~}{𝔪}`$ are of the form $$X=\left(\begin{array}{cc}0& b^t\\ b& 0\end{array}\right)\text{and}\stackrel{~}{X}=\left(\begin{array}{cc}0& b^t\\ b& 0\end{array}\right),$$ where $`b`$ is an arbitrary $`q\times p`$ matrix. For the inner products we take $`Q(X,Y)=\frac{1}{2}\mathrm{Tr}(XY)`$ and $`\stackrel{~}{Q}(\stackrel{~}{X},\stackrel{~}{Y})=+\frac{1}{2}\mathrm{Tr}(\stackrel{~}{X}\stackrel{~}{Y})`$. These will be “riemannian” dual symmetric spaces because the metrics on $`𝔪`$ and $`\stackrel{~}{𝔪}`$ are positive definite. We take the subgroups $`H`$ and $`\stackrel{~}{H}`$ to be $`\mathrm{SO}(p)\times \mathrm{SO}(q)`$ and the Lie algebra isomorphism $`S:𝔥\stackrel{~}{𝔥}`$ is the identity transformation. The isometry $`T:𝔪\stackrel{~}{𝔪}`$ is taken to be $$T:\left(\begin{array}{cc}0& b^t\\ b& 0\end{array}\right)\left(\begin{array}{cc}0& b^t\\ b& 0\end{array}\right).$$ A brief computation shows that $`[TX,TY]=[X,Y]`$. From this we see that $`\mathrm{SO}(p+q)/\mathrm{SO}(p)\times \mathrm{SO}(q)`$ is a space with positive sectional curvature and $`\mathrm{O}_0(p,q)/\mathrm{SO}(p)\times \mathrm{SO}(q)`$ is a space with negative sectional curvature. A more extensive discussion of dual symmetric spaces requires a thorough discussion of the theory of orthogonal involutive Lie algebras (orthogonal symmetric Lie algebras), see . ## 6 Ivanov’s Construction In this article we have limited ourselves to studying sigma models on an arbitrary riemannian manifold with vanishing $`3`$-form $`H_{ijk}`$. Ivanov analyzed a subset of those models. He studied sigma models that are Lie group theoretic in origin. In this case we can use special properties of Lie groups to simplify the analysis. Assume we have a normal symmetric space $`G/H`$ as in Section 5. We can choose an orthonormal basis that respects the decomposition $`𝔤=𝔥𝔪`$. If $`\{T_i\}`$ is such a basis with Lie bracket relations $$[T_i,T_j]=f^k{}_{ij}{}^{}T_{k}^{}$$ then the $`\mathrm{Ad}(G)`$-invariance of $`Q`$ tells us that $`f_{ijk}`$ are totally antisymmetric. The indices $`a,b,c,\mathrm{}`$ are associated to the basis elements that are in $`𝔥`$ and the indices $`\alpha ,\beta ,\gamma ,\mathrm{}`$ are associated to basis elements in $`𝔪`$. With this notation the Lie algebra bracket relations are $`[T_a,T_b]`$ $`=`$ $`f^c{}_{ab}{}^{}T_{c}^{},`$ (6.1) $`[T_a,T_\beta ]`$ $`=`$ $`f^\gamma {}_{a\beta }{}^{}T_{\gamma }^{},`$ (6.2) $`[T_\alpha ,T_\beta ]`$ $`=`$ $`f^c{}_{\alpha \beta }{}^{}T_{c}^{}.`$ (6.3) Let $`\theta ^i`$ be the the Maurer-Cartan forms for the Lie group $`G`$ that satisfy the Maurer-Cartan equations $$d\theta ^i+\frac{1}{2}f^i{}_{jk}{}^{}\theta _{}^{j}\theta ^k=0.$$ Using the decomposition $`𝔤=𝔥𝔪`$ we may write the equations above as $`d\theta ^a+{\displaystyle \frac{1}{2}}f^a{}_{bc}{}^{}\theta _{}^{b}\theta ^c`$ $`=`$ $`{\displaystyle \frac{1}{2}}f^a{}_{\beta \gamma }{}^{}\theta _{}^{\beta }\theta ^\gamma ,`$ (6.4) $`d\theta ^\beta +f^\beta {}_{b\gamma }{}^{}\theta _{}^{b}\theta ^\gamma `$ $`=`$ $`0.`$ (6.5) It is worthwhile discussing the geometric meaning of the above equations. It is well known that $`G`$ is a principal $`H`$-bundle over $`G/H`$. The $`H`$-connection $`\theta ^aT_a`$ on $`G`$ defines the horizontal tangent spaces. Its curvature is given by the right hand side of (6.4). Equation (6.5) states that the covariant differential of the $`𝔪`$-valued $`1`$-form $`\theta ^\beta T_\beta `$ is zero. The equations of motion nonlinear sigma model with target space $`G/H`$ may be formulated in terms of a map $`g:\mathrm{\Sigma }G`$ that satisfies the equations $`d\theta ^a+{\displaystyle \frac{1}{2}}f^a{}_{bc}{}^{}\theta _{}^{b}\theta ^c`$ $`=`$ $`{\displaystyle \frac{1}{2}}f^a{}_{\beta \gamma }{}^{}\theta _{}^{\beta }\theta ^\gamma ,`$ (6.6) $`d\theta ^\beta +f^\beta {}_{b\gamma }{}^{}\theta _{}^{b}\theta ^\gamma `$ $`=`$ $`0,`$ (6.7) $`d(_\mathrm{\Sigma }\theta ^\beta )+f^\beta {}_{b\gamma }{}^{}\theta _{}^{b}(_\mathrm{\Sigma }\theta ^\gamma )`$ $`=`$ $`0.`$ (6.8) Here we implicitly interpret $`\theta ^i`$ as the pullback under $`g:\mathrm{\Sigma }G/H`$. More properly we should have written $`g^{}\theta ^i`$. We use $`_\mathrm{\Sigma }`$ to denote the Hodge duality operation on $`\mathrm{\Sigma }`$. On $`1`$-forms it is given by $`_\mathrm{\Sigma }(d\sigma ^\pm )=\pm d\sigma ^\pm `$. The first two equations above are just the pullbacks to $`\mathrm{\Sigma }`$ of the Maurer-Cartan equations for $`G`$. The third equation (6.8) is essentially the wave equation. For future reference we note that if $`\alpha `$, $`\beta `$ are $`1`$-forms on $`\mathrm{\Sigma }`$ then $$(_\mathrm{\Sigma }\alpha )(_\mathrm{\Sigma }\beta )=\alpha \beta ,$$ (6.9) and that $`(_\mathrm{\Sigma })^2\alpha =\alpha `$. Ivanov observed that if we define $`\stackrel{~}{\theta }^\beta =_\mathrm{\Sigma }\theta ^\beta `$ and $`\stackrel{~}{\theta }^a=\theta ^a`$ then the equations above become $`d\stackrel{~}{\theta }^a+{\displaystyle \frac{1}{2}}f^a{}_{bc}{}^{}\stackrel{~}{\theta }_{}^{b}\stackrel{~}{\theta }^c`$ $`=`$ $`+{\displaystyle \frac{1}{2}}f^a{}_{\beta \gamma }{}^{}\stackrel{~}{\theta }_{}^{\beta }\stackrel{~}{\theta }^\gamma ,`$ (6.10) $`d\stackrel{~}{\theta }^\beta +f^\beta {}_{b\gamma }{}^{}\stackrel{~}{\theta }_{}^{b}\stackrel{~}{\theta }^\gamma `$ $`=`$ $`0,`$ (6.11) $`d(_\mathrm{\Sigma }\stackrel{~}{\theta }^\beta )+f^\beta {}_{b\gamma }{}^{}\stackrel{~}{\theta }_{}^{b}(_\mathrm{\Sigma }\stackrel{~}{\theta }^\gamma )`$ $`=`$ $`0.`$ (6.12) These equations may be interpreted as the equations of motion for a sigma model on the symmetric space $`\stackrel{~}{G}/H`$ where $`\stackrel{~}{G}`$ is a real Lie group with real Lie algebra $`\stackrel{~}{𝔤}=𝔥(i𝔪)`$. This Lie algebra has Lie brackets given by $`[\stackrel{~}{T}_a,\stackrel{~}{T}_b]`$ $`=`$ $`f^c{}_{ab}{}^{}\stackrel{~}{T}_{c}^{},`$ (6.13) $`[\stackrel{~}{T}_a,\stackrel{~}{T}_\beta ]`$ $`=`$ $`f^\gamma {}_{a\beta }{}^{}\stackrel{~}{T}_{\gamma }^{},`$ (6.14) $`[\stackrel{~}{T}_\alpha ,\stackrel{~}{T}_\beta ]`$ $`=`$ $`f^c{}_{\alpha \beta }{}^{}\stackrel{~}{T}_{c}^{}.`$ (6.15) In the above we have $`\stackrel{~}{T}_a=T_a`$ and $`\stackrel{~}{T}_\beta =iT_\beta `$. Note that if you think of $`\stackrel{~}{G}`$ as a principal $`H`$-bundle then the curvature of $`\stackrel{~}{G}`$, given by (6.10), is “opposite” to that of $`G`$, given by (6.6). ## 7 Discussion In it was shown that a necessary condition for target space duality is that the target spaces be parallelizable manifolds. In the scenario presented here where we only require “on shell” pseudoduality and we see that the parallelizable requirement is weakened substantially but we still find natural geometric restrictions on the target spaces. In the models considered here where the $`3`$-forms $`H`$ and $`\stackrel{~}{H}`$ vanish we saw that pseudoduality requires that the target spaces be symmetric spaces with the “opposite curvature” generalizing a result of Ivanov . The class of riemannian dual symmetric spaces provides a wealth of examples. In particular we studied explicitly the example of dual Grassmann manifolds $`\mathrm{SO}(p+q)/\mathrm{SO}(p)\times \mathrm{SO}(q)`$ and $`\mathrm{O}_0(p+q)/\mathrm{SO}(p)\times \mathrm{SO}(q)`$. We also saw the importance of the isometry $`T`$. Note that we never explicitly solved for $`T`$ but we discovered conditions that guarantee its existence. Finally it should be emphasized that $`T`$ depends on the path since it comes from integrating (4.2) along the path. We note that equations (1.12) may also be written as $`\stackrel{~}{x}_\tau `$ $`=`$ $`Tx_\sigma ,`$ (7.1) $`\stackrel{~}{x}_\sigma `$ $`=`$ $`Tx_\tau .`$ (7.2) Thus we immediately see that “particle-like” solutions ($`\sigma `$ independent) on $`M`$ get mapped into static “soliton-like” solutions on $`\stackrel{~}{M}`$ and vice-versa. If say $`\stackrel{~}{M}`$ has noncontractible loops then there will be stable “soliton-like” solutions. Recently Evans and Mountain constructed an infinite number of local conserved commuting charges on sigma models with the target space being a compact symmetric space. It would be interesting to apply the results of this paper to the construction of Evans and Mountain. ## Acknowledgments I would like to thank T.L. Curtright, H. Fenderya, L.A. Ferreira, L. Mezincescu, R. Nepomechie, J. Sánchez Guillén, I.M. Singer and C. Zachos for comments and discussions. I would also thank E. Ivanov for bringing to my attention reference . This work was supported in part by National Science Foundation grant PHY–9870101.
warning/0004/cond-mat0004157.html	ar5iv	text	# Trajectory versus probability density entropy ## I introduction The problem of establishing a connection between the Kolmogorov-Sinai (KS) entropy and the conventional entropy expressed in terms of probability density is an interesting problem that is attracting some attention in literature. Early work on this subject goes back to the discussion of Goldstein and Penrose: These authors, almost twenty years ago, established a connection between the KS entropy and a coarse-grained version of the distribution density entropy. The work of Ref. is based on a formal and rigorous mathematical treatment which for this reason might have eluded the attention of physicists working on this subject. Thus we restate the problem using intuitive arguments which also make it possible for us to account for the more recent literature on the subject. In fact, our heuristic treatment will allow us to relate the results of the more recent work of Latora and Baranger to the earlier work of Zurek and Paz. In addition to revisiting the problem of how to make the KS entropy emerge from a nonequilibrium dynamic picture, we shall touch also the intriguing problem of whether a thermodynamic perspective has to rest on the adoption of trajectories, as implied by the concept itself of KS entropy, or on the use of probability densities, advocated with strong arguments by Petrosky and Prigogine. It is convenient to stress that the KS entropy is a property of a single trajectory. The phase space is divided into cells, each cell being assigned a given label $`\omega _r`$. Then we define a sequence of symbols by means of a single trajectory: The sequence is determined assigning to any time step the label of the cell where the trajectory lies at that time step. The trajectory is supposed to be large enough as to yield reliable values for the probabilities determined through the numerical frequencies. This means that we fix a window of size $`N`$, and we move this window along the sequence. For any window position a string of symbols $`\omega _0,\omega _1,\mathrm{}\omega _{N1}`$ is determined. Moving the window of fixed size $`N`$ along the infinite sequence generated by the trajectory we have to evaluate how many times the same string of symbols appears, thereby leading us to determine the probability $`p(\omega _0,\omega _1,\mathrm{}\omega _{N1})`$. The KS entropy is then defined by $$h_{KS}\underset{N\mathrm{}}{lim}H(N)/N,$$ (1) where $`H(N)`$ is the conventional Shannon entropy of the window of size $`N`$ defined by $$H(N)=\underset{\omega _0,,\omega _1\mathrm{}\omega _{N1}}{}p(\omega _0,\omega _1,\mathrm{}\omega _{N1})ln[p(\omega _0,\omega _1,\mathrm{}\omega _{N1})].$$ (2) It is evident therefore that the KS entropy rests on trajectories, and, more specifically, it implies the adoption of only one trajectory of virtually infinite length. The KS entropy is very attractive because its value turns out to be independent of the repartition into cells of the phase space, due to the crucial role of the so called generating partitions. In the specific case where a natural invariant distribution exists, it is shown that $$h_{KS}=\underset{i}{}𝑑𝐱\rho _{eq}(𝐱)\lambda _i(𝐱)$$ (3) with $`\lambda _i(𝐱)>0`$.Note that $`𝐱`$ denotes the coordinate of a multidimensional phase space, $`\rho _{eq}(𝐱)`$ is the natural invariant distribution and $`\lambda _i(x)`$ is a local Lyapounov coefficient, with $`i=1,d`$, $`d`$ being the dimension of the system under study. From Eq.(3) we see that, as earlier pointed out, the KS entropy is independent of the repartition into cells. The original definition of Eq.(1), with $`N`$ thought of as time, means that the KS entropy, as a property of a single trajectory, is the rate of entropy increase per unit of time. However, since the single trajectory under examination is infinitely long, and explores in time all the phase space available, the KS entropy can also be expressed in the form of an average over the equilibrium distribution density, without any prejudice for the single trajectory nature of this “thermodynamic” property. According to Petrosky and Prigogine, on the contrary, the connection between dynamics and thermodynamics implies the use of the Liouville equation $$\frac{}{t}\rho (𝐱,t)=iL\rho (𝐱,t),$$ (4) where $`L`$ denotes both the classical and the quantum Liouville operator, and $`\rho (𝐱,t)`$ is the nonequilibrium distribution density. The reason for this choice is that the analysis of the Liouville operator, through the “rigged Hilbert” space, allows the appearance of complex eigenvalues which correspond to irreversibility, and to the collapse of trajectories as well. This is the reason why distribution densities are judged to be more fundamental than trajectories. In this paper we limit our analysis to the special case where dynamics are generated by maps rather than by Hamiltonians. We do not address the difficult issue of discussing the thermodynamic limit $`N\mathrm{}`$ which is the subject of very interesting recent discussions, and where, according to Lebowitz, ergodicity and mixing are neither necessary nor sufficient to guarantee the connection between dynamics and thermodynamics. We consider the case of low-dimension chaos, where probability emerges as a consequence of sensitivity to initial conditions. Even in this case, however, according to the perspective established by Petrosky and Prigogine, probability densities are more fundamental than trajectories. The readers interested in knowing more about this perspective, entirely based on probability density, should consult the illuminating work of Driebe. In this case the counterpart of Eq.(4) becomes $$\rho (𝐱,t+1)=\mathrm{\Lambda }\rho (𝐱,t),$$ (5) where $`\mathrm{\Lambda }`$ is referred to as Frobenius-Perron operator. Of course, the operator $`L`$ of Eq.(4) has to be identified with $`i(\mathrm{\Lambda }1)`$. According to the traditional wisdom, the Frobenius-Perron operator is expected to make the distribution densities evolve in the same way as that resulting from the time evolution of a set of trajectories with initial conditions determined by the initial distribution density:The known cases of discrepancy between the two pictures are judged to be more apparent than real . Nevertheless, even in the case of invertible maps, the birth of irreversibility can be studied using the same perspective as that adopted for Hamiltonian systems, with Eq.(4) replaced by Eq.(5), and so using again probability densities rather than trajectories. However, we attempt at digging out the KS entropy from Eq.(5), and this purpose forces us to formulate a conjecture on how to relate entropy to $`\rho (𝐱,t)`$. A plausible choice seems to be $$S(t)=_𝐗\rho (𝐱,t)ln[\rho (𝐱,t)]𝑑𝐱.$$ (6) We share the view of Goldstein and Penrose who consider the KS entropy to be a nonequilibrium entropy. In other words, we may hope to derive the KS entropy from the time derivative of $`S(t)`$ of Eq.(6). As Goldstein and Penrose do, to realize that purpose we have to address a delicate problem: In the case of invertible maps, $`S(t)`$ is time independent, thereby implying a vanishing KS entropy. Yet, the baker’s transformation, which is a well known example of invertible map, thereby yielding a time independent $`S(t)`$, is shown to yield a KS entropy equal to $`ln2`$, a fact suggesting a steady condition of entropy increase. We plan to discuss all this with the joint use of heuristic arguments and of the rigorous theoretical tools of Ref.. The present paper uses as a paradygm of invertible map the two-dimensional baker’s transformation, depending on two coordinates, $`x`$ and $`y`$, the former corresponding to dilatation and the latter to contraction. Using this prototype for invertible dynamics, we aim at proving that the adoption of the distribution density in the case of invertible chaotic maps would lead to an increasing process of fragmentation, depending not only, as the KS entropy does, on the positive Lyapounov coefficient, but also on the negative one. The adoption of a coarse graining has the effect of quenching the action of the negative Lyapunov coefficient, thereby allowing the KS entropy to show up. Then, to go beyond these heuristic arguments we make a trace on the variable $`y`$, namely, on the process responsible for contraction, and we focus our attention on the contracted dynamics. This is equivalent to that produced by the Bernoulli shift map. Here room is only left for dilatation and the problem can be solved with a rigorous mathematical method, without using trajectories. The outline of the paper is as follows. In Section II we shall illustrate our heuristic picture. In Section III we shall address the problem by means of a rigorous treatment resting on the theoretical tools provided by Driebe. In Section IV we shall draw some conclusions. Some delicate mathematical problems behind the theoretical calculations of Section III are detailed in Appendix. ## II heuristic arguments Note that the cases studied by Latora and Baranger are two-dimensional, and our discussion here refers to a two-dimensional case, too. We have in mind the backer’s transformation and $`𝐱(x,y)`$. We denote by $`W(t)`$ the number of cells occupied at a given time $`t`$. Note that $`W(0)<W_{max}`$, where the symbol $`W_{max}`$ denotes the total number of cells into which we have divided the phase space $`𝐗`$. Our heuristic approach is based on the following assumptions. (i) At the initial time only $`W(0)`$ cells are occupied. (ii) At all times the trajectories are equally distributed over the set of occupied cells. This means $$S(t)=lnW(t).$$ (7) (iii) We denote by $`\lambda `$ the positive Lyapounov coefficient, and we set $$W(t)=W(0)exp(\lambda t).$$ (8) All these three assumptions have been borrowed from the recent work of Ref.. The joint use of all them yields $$S(t)=\lambda tlnW(0),$$ (9) which corresponds to the Kolmogorov thermodynamical regime. Note that the positive Lyapounov coefficient in the case of the baker’s transformation is shown to be: $$\lambda =ln2.$$ (10) Note also that according to the arguments of Section I, the connection with the KS entropy is established through the time derivative of $`S(t)`$. Thus, we conclude that $$\frac{dS}{dt}=\lambda =ln2=h_{KS},$$ (11) which corresponds to deriving the KS entropy from the distribution density picture. This Kolmogorov regime is not infinitely extended. It has an upper bound, given by the fact that when equilibrium is reached, even in the merely sense of a coarse-grained equilibrium, then the entropy stops increasing. An estimate of this time is obviously given by the solution of the following equation $$lnW_{max}=\lambda tlnW(0),$$ (12) which yields the following saturation time $$t_S=\frac{1}{\lambda }ln(\frac{W_{max}}{W(0)}).$$ (13) Furthermore a lower bound of validity exists, which will be easily estimated with very simple arguments. If the initial distribution includes a large number of cells and the size of this distribution along the coordinate $`y`$ is $`L`$, and the size of the cells is $`ϵ`$ with $`ϵ<L,`$, then it is evident that, in spite of the coarse graining the total number of cells occupied remains the same for a while. This time is easily estimated using the equation $$Lexp(\lambda t)=ϵ,$$ (14) which in fact defines the time at which the distribution volume, and consequently,the system entropy starts increasing. This time is denoted by the symbol $`t_D`$ and reads $$t_D=\frac{1}{\lambda }ln\left(\frac{L}{ϵ}\right).$$ (15) We denote by $`U(t)`$ the volume of the distribution density at time $`t`$ and by $`V`$ the volume of the phase space, thereby implying that $`U(t)V`$. We note that $$\frac{W_{max}}{W(0)}=\frac{V}{U(0)},$$ (16) where $`V`$ is the total volume of the phase space and $`U(0)`$ is the initial volume of the distribution density. Thus the Kolmogorov regime shows up in the following time interval $$t_D=\frac{1}{\lambda }ln\left(\frac{L}{ϵ}\right)<t<t_S=\frac{1}{\lambda }ln\left(\frac{V}{U(0)}\right).$$ (17) The time duration of the regime of validity of the Kolmogorov regime can be made infinitely extended by making the cell size infinitely small. This means that the conflict between the KS entropy prescription and the time independence of $`S(t)`$ can be bypassed by focusing our attention on the intermediate region, whose time duration tends to infinity with $`ϵ0`$. We note that a choice can be made such that $`V/U(0)=(L/ϵ)^\chi `$, with $`\chi >1`$. This means the time duration of the Kolmogorov regime can be made $`\chi `$ times larger than the time duration of the transition regime. For $`ϵ0`$ both time durations become infinite, thereby showing that a Kolmogorov regime of infinite time duration can be obtained at the price, however, of waiting an infinitely long time for the entropy to increase. The infinite waiting time before the regime of entropy increase fits the observation that the Gibbs entropy of an invertible map is constant. The linear entropy increase showing up “after this infinite waiting time” allows the emergence of the KS entropy from within the probability density perspective. This kind of coarse graining might be criticized as corresponding to arbitrary choices of the observer. It is interesting to remark that there exists another interesting form of coarse graining, produced by weak stochastic forces. Both in the case where this stochastic forces mimic the interaction with the environment or in the case where it happens to be an expression of spontaneous fluctuations this kind of coarse graining can be regarded as being produced by nature. Here we limit ourselves to remarking that according to Zurek and Paz these stochastic forces contribute a fluctuation-dissipation process mimicking the interaction between the system of interest and the environment. These authors studied the inverted stochastic oscillator $$\frac{d^2x}{dt^2}=\lambda ^2x(t)+\gamma \frac{dx}{dt}+f(t),$$ (18) where the friction $`\gamma `$ and the stochastic force $`f(t)`$ are related to one another by the standard fluctuation-dissipation relation $$<ff(t)>=2\gamma <(\frac{dx}{dt})^2>_{eq}\delta (t)2D\delta (t).$$ (19) It is interesting to remark that the proper formulation of the second principle implies that the entropy of a system can only increase or remain constant under the condition of no energy exchange between the system and its environment. In the case of Eq.(19) the energy exchange between system and environment is negligible for any observation made in the time scale $$t<<1/\gamma .$$ (20) To ensure that the system entropy increase to take place with no energy exchange between system and its environment Zurek and Paz set the condition of Eq.(20) and this, in turn, allows them to neglect the friction term in Eq.(18). Then, these authors adopted the modes $$u\frac{dx}{dt}+\lambda x$$ (21) and $$w\frac{dx}{dt}+\lambda x,$$ (22) which make it possible for them to split Eq.(18) into $$\frac{du}{dt}=\lambda u(t)+f(t)$$ (23) and $$\frac{dw}{dt}=\lambda w(t)+f(t).$$ (24) Let us imagine the initial distribution density as a rectangle of size $`\mathrm{\Delta }w(0)`$ along the direction $`w`$ and $`\mathrm{\Delta }u(0)`$ along the direction $`u`$. We keep denoting by $`U(t)`$ the distribution volume at a given time $`t`$. Thus the volume of the initial distribution is $$U(0)=\mathrm{\Delta }u(0)\mathrm{\Delta }w(0).$$ (25) In the absence of the stochastic force $`f(t)`$, Eqs.(23) and Eqs.(23) result in an exponential increase and an exponential decrease, with the same rate $`\lambda `$, respectively. Consequently, the Liouville theorem $`U(t)=U(0)`$ is fulfiled. In the presence of stochastic force, we work as follows. In the former equation, with $`u`$ increasing beyond any limit, the weak stochastic force $`f(t)`$ can be neglected. This is not the case with the latter equation. In fact, $`w`$ is a contracting variable in the absence of the stochastic force. In the presence of the stochastic force the minimum size of the distribution along $`w`$ is given by $$<w^2>_{eq}^{1/2}=(D/\lambda )^{1/2}.$$ (26) This minimum size is reached in a time determined by the solution of the following equation $$\mathrm{\Delta }w(0)exp(\lambda t)=(D/\lambda )^{1/2}$$ (27) yielding $$t_D=\frac{1}{\lambda }ln\left(\frac{\lambda }{D}\right)^{1/2}\mathrm{\Delta }w(0).$$ (28) Due to the fact that deterministic chaos is simulated by Zurek and Paz by means of an inverted parabola, these authors did not consider the entropy saturation effects. However, it is straigthforward to evaluate the saturation effect with heuristic arguments concerning the case where the total volume of the phase space has the finite value $`V`$. From the time $`t=t_D`$ on, the distribution volume $`U(t)`$ increases exponentially in time with the following expression $$U(t)=\mathrm{\Delta }w(0)\mathrm{\Delta }u(0)exp(\lambda t)=(D/\lambda )^{1/2}\mathrm{\Delta }u(t_D)exp(\lambda t).$$ (29) Thus, the saturation time is now given by $$t_S=\frac{1}{\lambda }ln\left[\frac{V}{\mathrm{\Delta }u(0)\mathrm{\Delta }w(0)}\right].$$ (30) Using Eq.(25) we can write this saturation time as $$t_S=\frac{1}{\lambda }ln\left[\frac{V}{U(0)}\right],$$ (31) which coincides with Eqs.(17) and (13). In conclusion, it seems that the emergence of a Kolmogorov regime is made possible by the existence of a form of coarse graining, and that it is independent of whether the coarse graining is realized by the division into cells or by a weak stochastic force. This property seems to make less important the discussion of whether the stochastic force is of environmental origin or rests on some kind of extension of the current physical laws. However, we have to point out that the situation significantly changes, if we move from a strongly to a weakly chaotic classical system. As a relevant example, let us refer ourselves to the work of Ref.. The authors of this work study the asymptotic time limit of a diffusion process generated by using an intermittent map as a dynamic generator of diffusion. If these dynamics are perturbed by a white noise, a transition is provoked, at long times, from anomalous to normal diffusion. When the only source of random behavior is given by the sporadic randomness of the intermittent map, the long-time limit is characterized by Lévy statistics, a physical condition in a striking conflict with the condition of Gaussian statistics produced by the action of fluctuations. Here we limit our attention to the case of strong chaos where the two distinct sources of coarse graining produce equivalent effects. It might be of some interest for the reader to compare the coarse-graining approach of this section to the more formal method recently adopted by Fox to deal with the same problem. It is interesting to stress that to make the regime of validity of the Kolmogorov regime as extended as possible we must make the ratio $`V/U(0)`$ as large as possible (virtually infinite). This means that we have to choose an initial distribution density so sharp as to become apparently equivalent to a single trajectory. This seems to be an attractive way of explaining why in this condition the KS entropy is recovered, since, as stressed in Section I, the KS entropy is a single trajectory property. However, in accordance with the authors of Refs. we must admit that there exists a deep difference between a trajectory and a very sharp distribution. The latter is stable and robust, while the former is not. In Section III we shall show that the rigorous derivation of the Kolmogorov regime requires a non trivial mathematical procedure, and the mathematical effort to make from this side, to derive the KS entropy, serves the useful purpose of proving that the KS entropy of a trajectory is a really wise way of converting into advantages the drawbacks of the trajectory instability. ## III The KS entropy from a reduced Frobenius-Perron equation This Section is devoted to a rigorous discussion resting only on the theoretical tools described in Ref. for a genuine probability density aproach. According to Mackey, if we rule out the possibility that the laws of physics are misrepresented by invertible dynamic prescriptions, there are only two possible sources of entropy increase. The first is the coarse graining discussed in Section II. The second is the adoption of reduced equation of motion, obtained by a trace over “irrelevant” degrees of freedom. In fact here we study the Bernoulli shift map, $$x_{t+1}=2x_t,mod1.$$ (32) The Frobenius-Perron equation of this map is defined by $$\rho (x,t+1)=\mathrm{\Lambda }\rho (x,t)\frac{1}{2}[\rho (\frac{x}{2},t)+\rho (\frac{x+1}{2},t)].$$ (33) It is straigtforward to show that the Frobenius-Perron operator of Eq.(33) stems from the contraction over the variable $`y`$ of the baker’s mapping, acting in fact on the unit square of two-dimensional space (x,y) (see, for instance Ref.). It is shown that the KS entropy of the baker’s transformation is well defined and turns out to be the same as that of the Bernoulli shift map, namely $`h_{KS}=ln2`$. Intuitively, this suggests that the main role of coarse graining is that of making inactive the process of contraction, and with it the negative Lyapunov coeficient. This intuitive argument seems to be plausible and raises the interesting question of how to prove it with a rigorous approach. This is equivalent to deriving the Kolmogorov regime using a rigorous mathematical method rather than the heuristic arguments of Section II. We must observe again that this is made possible by the fact that the tracing has changed the originally invertible map into one that is not invertible. To address this issue we follow the prescription of Ref.. First of all, we express the distribution density at time $`t`$ under the form given by Ref. which reads: $$\rho (x,t)=1+\underset{j=1}{\overset{\mathrm{}}{}}exp(\gamma _jt)\frac{B_j(x)}{j!}[\rho ^{(j1)}(1,0)\rho ^{(j1)}(0,0)].$$ (34) Note that $`\gamma _jjln2`$, $`B_j(x)`$ are the Bernoulli polynomials and $`\rho ^{(n)}(x,t)`$ denotes the $`n`$-th order derivative of $`\rho (x,t)`$ with respect to $`x`$. Hereby, we shall show how to derive from the previous one more tractable expression,which will be checked in appendix. In the case of an initial condition close to equilibrium, resulting from the sum of the equilibrium distribution and the first “excited” state, it is easy to prove that the entropy $`S(t)`$ of Eq. (6) reaches exponentially in time the steady-state condition. This suggests that the Kolmogorov regime, where the entropy $`S(t)`$ is expected to be a linear function of time, must imply an initial condition with infinitely many “excited” states. To deal with a condition of this kind it is convenient to express Eq.(34) in an equivalent form given by $$\rho (x,t)=1+\underset{j=1}{\overset{\mathrm{}}{}}_{\mathrm{}}^+\mathrm{}\frac{B_j(x)}{j!}(iz\omega )^{j1}\widehat{\rho }(\omega )(exp(i\omega )1)\frac{d\omega }{2\pi },$$ (35) where $`zexp[t(ln2)]`$ and $`\widehat{\rho }(\omega )`$ is related to the initial condition $`\rho (x,0)`$ by the Fourier transform $$\rho (x,0)=_{\mathrm{}}^+\mathrm{}\widehat{\rho }(\omega )exp(i\omega x)\frac{d\omega }{2\pi }.$$ (36) The following equation $$\underset{j=0}{\overset{\mathrm{}}{}}\frac{B_j(x)}{j!}z^j=z\frac{exp(zx)}{exp(z)1},$$ (37) is known to generate Bernoulli polynomials. Using this Bernoulli polynomial generatrix, we arrive, after some algebra, at $$\rho (x,t)=z_{\mathrm{}}^+\mathrm{}exp(i\omega zx)\widehat{\rho }(\omega )\frac{exp(i\omega )1}{exp(i\omega z)1}\frac{d\omega }{2\pi }.$$ (38) By expanding the denominator of Eq.(38) into a Taylor series and using Eq.(36), we finally derive the fundamental expression $$\rho (x,t)=z\underset{n=0}{\overset{\mathrm{}}{}}[\rho (zx+zn,0)\rho (zx+zn+1,0)].$$ (39) This important expression makes it possible for us to discuss analytically the entropy time evolution ensuing the preparation of an initially very sharp distribution. Let us consider in fact $$\rho (x,0)=\frac{\alpha }{1exp(\alpha )}exp(\alpha x),0x1.$$ (40) For $`\alpha \mathrm{}`$ this initial distribution becomes a very sharp distribution located at $`x=0`$. By plugging this initial distribution density into Eq.(39) we obtain $$\rho (x,t)=z\alpha \frac{exp(\alpha xz)}{1exp(\alpha z)}.$$ (41) It is evident that this simple analytical expression for the time evolution of the distribution density is exact, and corresponds to the time evolution dictated by the Frobenius-Perron operator of Eq.(33). We are now in a position to discuss the central issue of this paper, namely, the time evolution of the Gibbs entropy of Eq.(6), which, in the case here under study, reads $$S(t)=_X\rho (x,t)ln[\rho (x,t)]𝑑x,$$ (42) with $`X`$ now denoting the interval $`[0,1]`$. By plugging Eq.(40) within Eq.(42) we obtain $$S(t)=1ln(\alpha z)+\mathrm{ln}\left[1exp(\alpha z)\right]\frac{\alpha z}{exp(\alpha z)1}.$$ (43) In the limiting case $`\alpha \mathrm{}`$ this exact prediction is approximated very well by $$S(t)=ln(\alpha )+(ln2)t.$$ (44) It indicates that a sharp initial distribution makes the system evolve according to the KS entropy, with no regime of transition from mechanics to thermodynamics. The third regime of Ref. is still present. It is straigtforward to show that the saturation time $`t_S=ln\alpha /ln2`$ resulting from Eq.(43) is the same as that of Eq.(13) in the case $`V=1`$. In fact using Eq.(16) and $`V=1`$ we obtain that $`W_{max}/W(0)=1/U(0)`$, where $`U(0)`$ is the size of the initial distribution. The size of the initial distribution of Eq.(40), for $`\alpha \mathrm{}`$, becomes prportional to $`1/\alpha `$. Thus, $`ln\alpha ln(W_{max}/W(0)`$ in accordance with Eq. (13). This is an elegant result, involving a modest amount of algebra. However, it refers to an initial distribution located at $`x=0`$. We want to prove that this is a general property, independent of where the initially sharp distribution is located, at the price, as we shall see, of a more complicated mathematical treatment. For this purpose we study the case where the distribution shape is the Lorentzian curve: $$\rho (x,0)=A\frac{\mathrm{\Gamma }}{\left(xx_0\right)^2+\mathrm{\Gamma }^2},$$ (45) with $`x_0`$ being a generic point of the interval $`[0,1]`$ and $`x`$ running in the same interval. Setting the normalization condition yields $$A=\frac{1}{\mathrm{arctan}\left(\frac{x_0}{\mathrm{\Gamma }}\right)+\mathrm{arctan}\left(\frac{1x_0}{\mathrm{\Gamma }}\right)}.$$ (46) We have to set again the condition that the initial distribution is very sharp. Thus we make the assumption $`\mathrm{\Gamma }0`$, yielding $`A1/\pi `$. We plug this approximated value of $`A`$ into Eq. (39), thereby obtaining the following density time evolution $$\rho (x,t)=\frac{z\mathrm{\Gamma }}{\pi }\underset{n=0}{\overset{\mathrm{}}{}}\left[\frac{1}{\left(zx+znx_0\right)^2+\mathrm{\Gamma }^2}\frac{1}{\left(zx+znx_0+1\right)^2+\mathrm{\Gamma }^2}\right].$$ (47) We are now in a position to study the entropy time evolution again. Plugging Eq.(47) into (42) we find $`S\left(t\right)`$ $`=`$ $`{\displaystyle \underset{X}{}}\rho (x,t)\mathrm{ln}\left[{\displaystyle \frac{z}{\pi \mathrm{\Gamma }}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\left(\frac{zx+znx_0}{\mathrm{\Gamma }}\right)^2+1}}{\displaystyle \frac{1}{\left(\frac{zx+znx_0+1}{\mathrm{\Gamma }}\right)^2+1}}\right]𝑑x`$ $`={\displaystyle \underset{X}{}}\rho (x,t)\mathrm{ln}\left[{\displaystyle \frac{z}{\pi \mathrm{\Gamma }}}\right]𝑑x`$ $``$ $`{\displaystyle \underset{X}{}}\rho (x,t)\mathrm{ln}\left[{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\left(\frac{zx+znx_0}{\mathrm{\Gamma }}\right)^2+1}}{\displaystyle \frac{1}{\left(\frac{zx+znx_0+1}{\mathrm{\Gamma }}\right)^2+1}}\right]𝑑x,`$ where $`[y]`$ denotes the integer part of $`y`$. To derive a more tractable expression we note that in the limiting case of $`\mathrm{\Gamma }`$ very small, the quantities $`[\mathrm{}]`$ contributing this series are almost zero except for $`n=\left[x\right]+\left[\frac{x_0}{z}\right]=\left[\frac{x_0}{z}\right]`$ in the first term, and for the possible contribution $`n=\left[x\right]\left[\frac{1x_0}{z}\right]=\left[\frac{1x_0}{z}\right]`$ in the second term. The latter condition cannot be realized, since $`n`$ is a positive integer. Thus, using only the first class of contributions, we get for the entropy $`S(t)`$ the following approximate expression $$S\left(t\right)\mathrm{ln}\mathrm{\Gamma }+\left(\mathrm{ln}2\right)t\mathrm{ln}\pi \underset{0}{\overset{z/\mathrm{\Gamma }}{}}\frac{\mathrm{ln}\left(y^2+1\right)}{y^2+1}𝑑y,$$ (48) which, in the limiting case $`\frac{z}{\mathrm{\Gamma }}\mathrm{}`$ becomes $$S\left(t\right)\mathrm{ln}\mathrm{\Gamma }+\left(\mathrm{ln}2\right)t\mathrm{ln}\pi \underset{0}{\overset{\mathrm{}}{}}\frac{\mathrm{ln}\left(y^2+1\right)}{y^2+1}𝑑y=\mathrm{ln}\mathrm{\Gamma }+\left(\mathrm{ln}2\right)t\mathrm{ln}\pi \pi \mathrm{ln}2\mathrm{ln}\mathrm{\Gamma }+\left(\mathrm{ln}2\right)t.$$ (49) As in the earlier case, the validity of the approximation yielding the linear dependence of $`S(t)`$ on time, is broken at the time $`t\mathrm{ln}\mathrm{\Gamma }/\mathrm{ln}2`$. In conclusion, the Kolmogorov condition is realized by very sharp initial distributions. Our derivation of this interesting property was done adhering to the recommendation of Ref. of resting only on densities rather than on trajectories. ## IV concluding remarks This paper shows that there exists a subtle difference between a Liouville density and a probability distribution. The Liouville distribution density of a chaotic map, which is at the same time invertible, becomes increasingly fragmented with time. If the initial distribution has a volume which is much smaller than the volume of the phase space, the highly fragmented distribution density at large times has a volume identical to the initial, but the impression afforded by a coarse-grained observation is that the initial volume increases with time till to become equal to that of the whole phase space. If we do the calculation of the time evolution of the Liouville density observing the motion of many trajectories with slightly different initial conditions, we are forced to adopt a coarse-graining procedure, dictated by the need itself of counting how many trajectories are found at a given time in a given small region of the phase space. This has the effect of making the Gibbs entropy increase. If, on the contrary, the calculation genuinely rests on the motion of the Liouville density, thereby implying that a quantum-like formalism is adopted, the Gibbs entropy is constant. On the other hand, the KS entropy is a trajectory property and this, to first sight, might lead us to believe that, being a trajectory property, cannot be recovered from within an approach genuinely resting on the distribution density. The heuristic arguments used in Section II subtly rest on the trajectory properties and thus on the assumption that the two perspectives are equivalent in spite of the warnings of the authors of Refs.. This is the reason why we judge the theoretical calculations of Section III to be of significant interest. These results have been obtained without having any direct recourse to the trajectory instability and only using the theoretical tools illustrated in Ref., which in fact address the problem of the dynamics of map using the quantum mechanical language of eigenstate and eigenvalue. It has to be pointed out that to establish this rigorous connection between the linear increase of the Gibbs entropy and the KS entropy we need to use infinitely many “excited” states, namely, a condition very far from equilibrium. The adoption of a “reduced” Liouville-like equation has been essential for the success of this enterprise. The Frobenius-Perron operator of the Bernoulli shift map is obtained from the Frobenius-Perron operator of the baker’s transformation via contraction over the variable $`y`$, corresponding to the contraction process. As a consequence, room is only left for dilatation. This is the reason why the heuristic argument of Eq.(8) holds true with no restriction and the use of a method rigorously based on densities lead to the same result. After all, the adoption of the theoretical tools of Ref. as a rigorous way to evaluate the regression to equilibrium, must reflect in a way the trajectory instabilities behind the Pesin theorem of Eq.(3). In the case where the action of the negative Lyapounov coefficient is quenched by a coarse-graining process, in addition to the saturation and to the Kolmogorov regime, also a short-time regime of transition to thermodynamics appears, so that, as found in Ref. three distinct regimes can be detected, the regime of transition to thermodynamics, the Kolmogorov regime, and the saturation regime. Here we show that the three regimes discussed by Latora and Baranger are exhibited also in the case of a coarse graining produced by random fluctuations. In conclusion, this paper contributes to deepening our understanding of the connection between density and trajectory distribution. The density entropy exhibits a regime of increase linear in time with a rate equivalent to that of the trajectory entropy if the initial distribution is very sharp. It has to stressed that it cannot be infinitely sharp. This would make the distribution density useless, since an infinitely sharp initial distribution would manifest the lack of robustness pointed out by Driebe. Thus, the connection is established using a genuine density. At the same time the accordance between the heuristic approach yielding the derivation of Eq.(44) can be interpreted as a rigorous support of the assumptions made by some some authors to derive the KS entropy from within a probability distribution approach. We want to remark that we are not aware of earlier work where the Lyapounov coefficient is derived analytically from a probability density approach with no use of heuristic arguments but that of Pattanayak and Brumer which, however, seems to be limited to the study of the transition regime. A problem open to future research concerns the role of coarse-graining. We have seen that in the case of strong chaos studied in this paper there is no essential difference between the coarse graining resulting from the repartition of the phase space into cells and the coarse graining caused by fluctuations. In the case of dynamic systems like the weakly chaotic billiards studied by Zaslavsky, which are easily proved to be statistically equivalent to the intermittent maps of Ref., the two different coarse-graining sources result in different physical effects at long times. On the other hand, it is expected that the dynamic process of transition to Lévy statistics is characterized by thermodynamic properties of non-extensive nature. The study of the time evolution of non-extensive entropy under the actions of the two different sources of coarse graining studied in this paper would an interesting program for future research work. APPENDIX The purpose of this Appendix is to check the main result of Section III. Our theoretical reference on this issues is given by the book of Ref. . We note, however, that Eqs. (41) and Eq.(47) are not directly derived from Eq.(34), which is a theoretical finding of Ref. , but they are derived from Eq.(39), wich is the result of a further development of the theory of Ref. . We feel therefore the need of proving that Eq.(39) fits the main requirement of keeping the norm unchanged and of being an exact solution of Frobenius-Perron equation of Eq.(33). To double check our results, we shall prove also that Eqs. (41) and Eq.(47) fit the same property. As a general remark about the content of this appendix we note that a function of the variable $`x`$, defined only within the finite interval $`[0,1]`$, admits a treatment based on its Fourier transform if it is thought of as being defined on the whole interval $`[\mathrm{},\mathrm{}]`$ with vanishing value outside $`[0,1]`$. Similarly we can define the Fourier series of this function assuming it to be periodically repeated all over the real axis. We shall adopt this approach throughout the whole appendix. Let us check Eq.(39) first. We note that the argument of the density of Eq.(39) can be arbitrary with the only condition that the variable $`x`$ is in the interval . Furthermore, since Eq.(39) is derived from Eq.(38), it is enough for us to prove that Eq.(38) is properly normalized and is a solution of the Frobenius-Perron equation of Eq.(33). In conclusion, we have to check: $$\rho (x,t)=z_{\mathrm{}}^+\mathrm{}e^{i\omega zx}\widehat{\rho }(\omega )\frac{e^{i\omega }1}{e^{i\omega z}1}\frac{d\omega }{2\pi }.$$ (A-1) First we check that this equation is norm conserving or: $$_0^1\rho (x,t)𝑑x=1.$$ (A-2) To do so, we integrate Eq.(A-1 with respect to the variable x from $`\mathrm{}`$ to $`+\mathrm{}`$. Thus we obtain: $`{\displaystyle _0^1}z{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\omega zx}\widehat{\rho }(\omega ){\displaystyle \frac{e^{i\omega }1}{e^{i\omega z}1}}{\displaystyle \frac{d\omega }{2\pi }}`$ $`dx`$ (A-3) $`={\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{e^{i\omega z}1}{ı\omega }}\widehat{\rho }(\omega ){\displaystyle \frac{e^{i\omega }1}{e^{i\omega z}1}}{\displaystyle \frac{d\omega }{2\pi }}.`$ (A-4) This means the expression $$_{\mathrm{}}^+\mathrm{}\frac{e^{i\omega }1}{ı\omega }\widehat{\rho }(\omega )\frac{d\omega }{2\pi }_0^1_{\mathrm{}}^+\mathrm{}e^{i\omega y}\widehat{\rho }(\omega )\frac{d\omega }{2\pi }𝑑y.$$ (A-5) The integral over $`\omega `$ is by definition the Fourier transform of $`\rho (y,0)`$, yielding thereby (A-6) $`{\displaystyle _0^1}\rho (x,t)𝑑x={\displaystyle _0^1}\rho (y,0)𝑑y=1t,`$ (A-7) due to the fact that the initial condition is assumed to be normalized. We want now to prove that the distribution density $`\rho (x,t)`$ of Eq. (A-1) is a solution of the Frobenius-Peron operator of Eq.(33), namely, that: $`\rho (x,t+1)={\displaystyle \frac{1}{2}}\left[\rho ({\displaystyle \frac{x}{2}},t)+\rho ({\displaystyle \frac{x+1}{2}},t)\right].`$ (A-8) Remembering that $`z=2^t`$ and $`\frac{z}{2}=2^{t1}`$ we can write Eq.(A-1)as: $`\rho (x,t+1)={\displaystyle \frac{z}{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\omega zx/2}\widehat{\rho }(\omega ){\displaystyle \frac{e^{i\omega }1}{e^{i\omega z/2}1}}{\displaystyle \frac{d\omega }{2\pi }}.`$ (A-10) Plugging Eq.(A-1) into it, the r.h.s. of Eq.(A-8) becomes: $$\frac{1}{2}\left[\rho (\frac{x}{2},t)+\rho (\frac{x+1}{2},t)\right]$$ $`={\displaystyle \frac{z}{2}}\left[{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\omega zx/2}\widehat{\rho }(\omega ){\displaystyle \frac{e^{i\omega }1}{e^{i\omega z}1}}{\displaystyle \frac{d\omega }{2\pi }}+{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\omega z(x+1)/2}\widehat{\rho }(\omega ){\displaystyle \frac{e^{i\omega }1}{e^{i\omega z}1}}{\displaystyle \frac{d\omega }{2\pi }}\right]`$ (A-11) and after a little algebra we get: $$\frac{1}{2}\left[\rho (\frac{x}{2},t)+\rho (\frac{x+1}{2},t)\right]=\frac{z}{2}_{\mathrm{}}^+\mathrm{}e^{i\omega zx/2}\widehat{\rho }(\omega )\frac{\left[e^{i\omega }1\right]\left[e^{i\omega zx/2}+1\right]}{e^{i\omega z}1}\frac{d\omega }{2\pi }.$$ (A-12) By decomposing the denominator as follows $$\frac{1}{e^{i\omega z}1}=\frac{1}{\left[e^{i\omega z/2}1\right]\left[e^{i\omega z/2}+1\right]}$$ (A-13) and simplifying, we obtain $$\frac{1}{2}\left[\rho (\frac{x}{2},t)+\rho (\frac{x+1}{2},t)\right]=\frac{z}{2}_{\mathrm{}}^+\mathrm{}e^{i\omega zx/2}\widehat{\rho }(\omega )\frac{e^{i\omega }1}{e^{i\omega z/2}1}\frac{d\omega }{2\pi }$$ (A-14) that coincide with Eq.(A-10). To check the property $`\rho (x,t)1`$ for $`t\mathrm{}`$ or $`z0`$ we use (A-1) and we write: $`\rho (x,t)z{\displaystyle _{\mathrm{}}^+\mathrm{}}\widehat{\rho }(\omega ){\displaystyle \frac{e^{i\omega }1}{i\omega z}}{\displaystyle \frac{d\omega }{2\pi }}`$ (A-15) $`={\displaystyle _{\mathrm{}}^+\mathrm{}}\widehat{\rho }(\omega ){\displaystyle \frac{e^{i\omega }1}{i\omega }}{\displaystyle \frac{d\omega }{2\pi }}={\displaystyle _{\mathrm{}}^+\mathrm{}}\widehat{\rho }(\omega ){\displaystyle _0^1}e^{i\omega x}𝑑x{\displaystyle \frac{d\omega }{2\pi }}=`$ $`{\displaystyle _0^1}\rho (x)𝑑x=1.`$ (A-16) Now we shall test directly the Eq.(41) using Eq.(34). Plugging directly the initial distribution: $$\rho (x,0)=\frac{\alpha }{1e^\alpha }e^{\alpha x},0x1$$ (A-17) into Eq.(34) we obtain: $`\rho (x,t)={\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}e^{\gamma _jt}{\displaystyle \frac{B_j(x)}{j!}}[\rho ^{(j1)}(1,0)\rho ^{(j1)}(0,0)]`$ (A-18) $`={\displaystyle \frac{\alpha }{1e^\alpha }}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}e^{\gamma _jt}{\displaystyle \frac{B_j(x)}{j!}}[(\alpha ^{j1})e^\alpha (\alpha ^{j1})].`$ (A-19) Using the Bernoulli polynomials generatrix we get $$\rho (x,t)=z\alpha \frac{e^{\alpha xz}}{1e^{\alpha z}}.$$ (A-20) that coincides with Eq.(41) obtaned using the formula Eq.(39) Finally let us to check the norm conservation of Eq.(47). Without any approximation, using the value of $`A`$ of Eq.(46), and a little algebra, we get: $`{\displaystyle _0^1}\rho (x,t)𝑑x=\left[\mathrm{arctan}\left({\displaystyle \frac{x_0}{\mathrm{\Gamma }}}\right)+\mathrm{arctan}\left({\displaystyle \frac{1x_0}{\mathrm{\Gamma }}}\right)\right]^1`$ (A-21) $`{\displaystyle _0^1}z\mathrm{\Gamma }{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{1}{\left(zx+znx_0\right)^2+\mathrm{\Gamma }^2}}{\displaystyle \frac{1}{\left(zx+znx_0+1\right)^2+\mathrm{\Gamma }^2}}]dx`$ (A-22) $`=\left[\mathrm{arctan}\left({\displaystyle \frac{x_0}{\mathrm{\Gamma }}}\right)+\mathrm{arctan}\left({\displaystyle \frac{1x_0}{\mathrm{\Gamma }}}\right)\right]^1{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}\left[\mathrm{arctan}\left({\displaystyle \frac{xz+nzx_0}{\mathrm{\Gamma }}}\right)\mathrm{arctan}\left({\displaystyle \frac{xz+nz+1x_0}{\mathrm{\Gamma }}}\right)\right]_{x=0}^{x=1}.`$ (A-23) This yields: $`{\displaystyle _0^1}\rho (x,t)𝑑x=\left[\mathrm{arctan}\left({\displaystyle \frac{x_0}{\mathrm{\Gamma }}}\right)+\mathrm{arctan}\left({\displaystyle \frac{1x_0}{\mathrm{\Gamma }}}\right)\right]^1`$ $`{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}[\mathrm{arctan}\left({\displaystyle \frac{z+nzx_0}{\mathrm{\Gamma }}}\right)\mathrm{arctan}\left({\displaystyle \frac{z+nz+1x_0}{\mathrm{\Gamma }}}\right)\mathrm{arctan}\left({\displaystyle \frac{nzx_0}{\mathrm{\Gamma }}}\right)+\mathrm{arctan}\left({\displaystyle \frac{nz+1x_0}{\mathrm{\Gamma }}}\right)].`$ Examining the expression we can note that only the terms with $`n=0`$ survive in the sum but that terms simplify with the external factor (the constant $`A`$) so finally $$_0^1\rho (x,t)𝑑x=1t.$$ (A-24) Checking that Eq.(47) fulfils the Frobenius-Perron operator involves some extended but straightoforward algebra.
warning/0004/cond-mat0004462.html	ar5iv	text	# Transport properties of heterogeneous materials derived from Gaussian random fields: Bounds and Simulation. ## I Introduction The calculation of the effective transport properties of random composite media is important in many scientific and engineering applications. Several techniques (effective medium approximations and cluster expansions) have been developed for predicting the effective properties of such materials (briefly reviewed in ref. ). However difficulties encountered in such methods have provided the impetus for the development of rigorous bounds . Such bounds rely on statistical descriptions of the microstructure of the material which are available for relatively few classes of media. The advancement of computing technology has also made direct simulation of effective properties feasible. It is the latter two approaches which we shall discuss here, in the context of the effective conductivity of a three dimensional amorphous isotropic two phase material. There has been significant advances in the evaluation of the Beran-Milton (BM) bounds in the past decade. The key parameter $`\zeta _1`$ (or $`\zeta _2=1\zeta _1`$) which incorporates microstructural information regarding the composite has been evaluated primarily for materials comprised of statistically independent cells or dispersions of regularly shaped inclusions . The simulation of the effective conductivity of continuum random media is a computationally intensive process and has only recently been studied for the second class of materials. Such an approach provides a basis for testing the bounding theories and for generating outright predictions of the effective properties of composites. A class of materials which is not, in general, well described by cellular or particulate models is that of amorphous composites. Such materials arise in certain alloys , microemulsions and other systems . The model which best captures some of the salient features of such composites is the spatially uncorrelated penetrable sphere (or the identical overlapping sphere) model . Due to the simplicity in evaluating the statistical correlation functions of such a material it has served as a useful ‘model’ amorphous medium . However specific features of this model restrict its generality. The inclusion (sphere) phase and the matrix phase are topologically very different, the small scale structure of the phase boundaries is spherical and there are no long range correlations in the model. An alternative approach is to empirically measure the specific correlation functions of a sample and to apply the results in the evaluation of bounds . This approach is complicated and subject to error. It is therefore interesting to seek a more complex model of amorphous composites, yet simple enough so that the correlation functions can be calculated. Another method of modeling random composites is to define the interface between the phases as a level cut of some random field (see ref. for a review). Recent progress in the theory of interfaces of level-cut Gaussian random fields has made it possible to calculate the statistical information necessary for the evaluation of the BM bounds. There is evidence that the Gaussian random interface model is a good approximation to certain oil-water microemulsions and we conjecture that it is a reasonable model for amorphous alloys. In this paper we investigate the effective conductivity of such media using the above mentioned bounding techniques and computer simulations. The results are compared with the previously studied models to demonstrate the differences that arise. The paper is organised as follows. In section II we describe the equations governing the electric field in a composite medium and the bounds on the effective conductivity. In sections III & IV we derive the statistical correlation functions for the random media and apply them in the calculation of the microstructure parameter. Sections V-VII are concerned with the generation of the random materials, the simulation of the effective conductivity and comparison of the data with the bounds. ## II Bounds on the effective properties of composite materials The relationship between the current density $`j`$ and the electric field $`E=\varphi `$ is given by Ohms Law, $$𝐣=\sigma \varphi .$$ (1) Where, due to charge conservation, $`\varphi `$ satisfies, $$^2\varphi =0$$ (2) throughout the material. At the boundary of different regions of the material with conductivities $`\sigma _1`$ and $`\sigma _2`$ we have, $$\varphi _1=\varphi _2,\sigma _1\varphi _1.𝐧=\sigma _2\varphi _2.𝐧$$ (3) The effective conductivity is defined by a macroscopic form of Ohm’s law, $$\sigma _e=\frac{<\sigma \varphi >}{<\varphi >}.$$ (4) Now consider a composite material made up of two components with conductivities $`\sigma _1`$ and $`\sigma _2`$ with volume fractions $`p`$ and $`q=1p`$. The effective conductivity will then depend on the $`\sigma _i`$, their respective volume fractions, and the spatial distribution (microstructure) of each phase . The first bounds on $`\sigma _e`$ were calculated by Wiener who proved that $`<\sigma ^1>^1\sigma _e<\sigma >`$. These bounds assume no details about the microstructure and are hence valid for a general composite. As more statistical information regarding the composite is included in the calculation of the bounds they become more restrictive. If the sample is assumed to be isotropic and macroscopically homogeneous then the 2nd order bounds of Hashin and Sthrikman are applicable. To distinguish between such materials the third order bounds of Beran are necessary. (The term $`n`$th order bounds refers to the fact that the bounds are exact to $`O(\sigma _1\sigma _2)^n`$). The Beran bounds were derived using variational principles and were subsequently simplified by Milton . Following the notation of Milton we define $`<a>=pa_1+qa_2`$, $`<\stackrel{~}{a}>=qa_1+pa_2`$ (interchanging $`p`$ and $`q`$) and $`<a>_\zeta =\zeta _1a_1+\zeta _2a_2`$. Here $`a_i=\sigma _i`$ or $`1/\sigma _i`$. In these terms the lower bound on $`\sigma _e`$ is, $$\sigma _l=\left[<\sigma ^1>\frac{2pq(\sigma _1^1\sigma _2^1)^2}{2<\stackrel{~}{\sigma }^1>+<\sigma ^1>_\zeta }\right]^1$$ (5) while the upper bound is, $$\sigma _u=\left[<\sigma >\frac{pq(\sigma _1\sigma _2)^2}{<\stackrel{~}{\sigma }>+2<\sigma >_\zeta }\right].$$ (6) The so called microstructure parameter $`\zeta _1`$ is given by a number of equivalent integrals, of which the formulation due to Brown is the best for our purposes, $`\zeta _1`$ $`=`$ $`{\displaystyle \frac{9}{2pq}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dr}{r}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}{\displaystyle _1^1}duP_2(u)\times `$ (8) $`\left(p_3(r,s,t){\displaystyle \frac{p_2(r)p_2(s)}{p}}\right)`$ where $`t^2=r^2+s^22rsu`$ and $`P_2(u)=(3u^21)/2`$ is the Legendre polynomial of order 2. The functions $`p_n`$ are $`n`$-point solid-solid correlation functions (see section III) where the ‘solid’ is phase 1 and the ‘void’ is phase 2. As Milton notes these bounds converge when $`\zeta _1=0,1`$ and are equal to one of the 2nd order Hashin-Sthrikman bounds in each case. An improved lower bound has been derived by Milton for the case $`\sigma _2>\sigma _1`$. In later sections we consider materials with $`\sigma _1>\sigma _2`$ for which this bound is (see ref. ), by interchanging the roles of the materials, $$\sigma _l=\sigma _2\frac{1+(1+2p)\beta _{12}2(q\zeta _1p)\beta _{12}^2}{1+q\beta _{12}(2q\zeta _1+p)\beta _{12}^2},$$ (9) where $`\beta _{12}=(\sigma _1\sigma _2)/(\sigma _1+2\sigma _2)`$. By way of mathematical analogy these bounds also apply to to the effective dielectric, diffusion and magnetic permeability coefficients of composite materials. ## III Correlation Functions for the Gaussian Random Interface model There is an extensive literature on the calculation of statistical correlation functions . The case of the three point solid-solid correlation function has been considered empirically , and theoretically for cellular materials and spherical inclusions to name a few. Here we take the interface between the phases to be defined by a level cut of a random field . Now consider a Gaussian random field $`y(𝐫)`$ (see section V) and let the level sets $`y(𝐫)=\alpha `$ define the interface (with the region $`y>\alpha `$ being phase 1). Then the $`n`$ point correlation function is given by the volume average, $$p_n(𝐫_1,𝐫_2,\mathrm{},𝐫_n)=<H(y_1\alpha )\mathrm{}H(y_n\alpha )>$$ (10) where $`H(y)`$ is the Heavyside function and $`y_i=y(𝐫_i)`$. $`p_n`$ is then the probability that the $`n`$ points will lie in phase 1. For a macroscopically homogeneous isotropic material $`p_n`$ only depends on the distances $`r_{ij}=\|𝐫_i𝐫_j\|`$ between the points. Since volume and ensemble averages are equivalent in such a medium we can use the latter to evaluate eqn. (10). The joint probability density of $`y_i`$ is, $$P_n(y_1y_2\mathrm{}y_n)=\frac{1}{\sqrt{(2\pi )^n\|G\|}}\mathrm{exp}(\frac{1}{2}𝐲^TG^1𝐲),$$ (11) where the elements of $`G`$ are $`g_{ij}=<y_iy_j>`$ . The latter quantity we refer to as the field-field correlation function, $$g_{ij}=_0^{\mathrm{}}4\pi k^2\rho (k)\frac{\mathrm{sin}k\|𝐫_i𝐫_j\|}{k\|𝐫_i𝐫_j\|}𝑑k.$$ (12) where $`\rho (k)`$ is the spectral density of the field. Berk and Teubner have derived the one point function (volume fraction), $$p=\frac{1}{\sqrt{2\pi }}_\alpha ^{\mathrm{}}\mathrm{exp}(\frac{1}{2}t^2)𝑑t$$ (13) and the two point function, $$p_2(g_{ij})=\frac{1}{2\pi }_0^{g_{ij}}\mathrm{exp}\left(\frac{\alpha ^2}{1+t}\right)\frac{dt}{\sqrt{1t^2}}+p^2.$$ (14) The three point function is calculated using the techniques described in ref. . The following identities are used , $$<\mathrm{exp}(i𝐲.𝐰)>=\mathrm{exp}(\frac{1}{2}𝐰^TG𝐰)$$ (15) and $$H(y\alpha )=\frac{1}{2\pi i}_Ce^{iw(y\alpha )}\frac{dw}{w}$$ (16) where the contour $`C`$ is directed along the real axis except near the origin where it crosses the imaginary axis in the upper half plane. This leads (after algebra) to, $$\frac{p_3^T}{g_{12}}=\frac{p_2(g_{12})}{g_{12}}\frac{1}{\sqrt{2\pi }}_{\alpha F_{12}}^\alpha \mathrm{exp}(\frac{1}{2}t^2)𝑑t,$$ (17) where, $$F_{12}=\frac{\sqrt{1g_{12}}}{\sqrt{1+g_{12}}}\frac{1+g_{12}g_{13}g_{23}}{\sqrt{\|G\|}},$$ (18) and $`\|G\|=1g_{12}^2g_{13}^2g_{23}^2+2g_{12}g_{13}g_{23}`$. Similar expressions can be derived for $`p_3^T/g_{13}`$ and $`p_3^T/g_{23}`$. Defining $`A_{ij}=p_3^T/g_{ij}`$ we have therefore, $`p_3^T(g_{12},g_{13},g_{23})`$ $`=`$ $`g_{12}{\displaystyle _0^1}A_{12}(tg_{12},tg_{13},tg_{23})𝑑t`$ (19) $`+`$ $`g_{13}{\displaystyle _0^1}A_{13}(tg_{12},tg_{13},tg_{23})𝑑t`$ (20) $`+`$ $`g_{23}{\displaystyle _0^1}A_{23}(tg_{12},tg_{13},tg_{23})𝑑t.`$ (21) The truncated three point correlation function $`p_3^T`$ is related to the $`p_3`$ by the expression, $`p_3(g_{12},g_{13},g_{23})`$ $`=`$ $`p_3^T(g_{12},g_{13},g_{23})+pp_2(g_{12})`$ (22) $`+`$ $`pp_2(g_{13})+pp_2(g_{23})2p^3.`$ (23) To examine the limit $`r_{12}0`$ $`(g_{12}1,g_{23}g_{13})`$ set $`f(g_{13})=p_3^T(1,g_{13},g_{13})`$ then $`df(g_{13})/dg_{13}=(12p)dp_2(g_{13})/dg_{13}`$ and $`f(0)=0`$ implying $`p_3^T(1,g_{13},g_{13})=(12p)(p_2(g_{13})p^2)`$ as it should. (Similarly in the other limits). The X-ray spectra of these materials can be calculated from $`p_2`$ and hence they can be related to physical composites. Furthermore it has been shown that the surface to volume ratio is given by , $$\frac{S}{V}=\frac{2}{\pi }e^{\frac{1}{2}\alpha ^2}\sqrt{\frac{1}{3}<k^2>}.$$ (24) As the evaluation of the integrals in eqns. (14) & (22) are computationally intensive it is useful to derive various approximations. Rigorous approximations for $`p_2`$ for the cases $`\|\alpha \|1`$ and $`\|\alpha \|1`$ are derived in appendix A. A useful non-rigorous approximation to $`p_{123}^T`$ can be developed by requiring that the approximation have similar properties to the actual function for $`r_{ij}1`$ and satisfy the known consistency conditions in various limits . Using the compact notation $`p_{ij}^T=p_2(g_{ij})p^2`$ and $`p_{ijk}^T=p_3^T(g_{ij},g_{ik},g_{jk})`$ we have ($`r_{ij}1`$), $`p_{ij}^T`$ $``$ $`g_{ij}`$ (25) $`p_{123}^T`$ $``$ $`g_{12}g_{13}+g_{12}g_{23}+g_{13}g_{23}.`$ (26) Using this information, and including a higher order term for consistency ($`p_3(r_{12},r_{12},0)=p_2(r_{12})`$) we construct, $`p_{123}^T`$ $``$ $`{\displaystyle \frac{12p}{2p(1p)}}\left(p_{12}^Tp_{13}^T+p_{12}^Tp_{23}^T+p_{13}^Tp_{23}^T\right)`$ (28) $`{\displaystyle \frac{12p}{2p^2(1p)^2}}p_{12}^Tp_{13}^Tp_{23}^T.`$ We note that this approximation has a maximum absolute error of $`O(10^3)`$ for the materials considered here. As such it is an order of magnitude better than a previously suggested approximation $`p_3(r_{12},r_{13},r_{23})p_2(u)p_2(v)/p`$ where $`u`$ and $`v`$ are the smallest, and next to smallest, values of $`r_{12},r_{13}`$ and $`r_{23}`$. ## IV Determination of $`\zeta _1`$ Actual calculations of the microstructural parameter $`\zeta _1`$ (eqn. (8)) have, to date, been for four classes of materials. Cellular materials , empirically measured physical composites , periodic arrays of spheres and materials with spherical inclusions. In the latter class the cases studied include: identical overlapping spheres (the IOS model), identical hard spheres , and poly-dispersed spheres (many of these results are summarised in ref. ). We now describe aspects of the computation of $`\zeta _1`$ for several spectra of the Gaussian random interface (GRI) model. It can be shown that $`\zeta _1=\frac{1}{2}`$ for $`p=\frac{1}{2}`$ (see appendix B) and that $`\zeta _1=1\zeta _2`$ where $`\zeta _2`$ is the microstructure parameter associated with phase 2. As $`\zeta _1`$ is dimensionless it must depend only on the ratios of the length scales associated with the spatial variables in the (dimensionless) correlation functions. That this should be so also follows from a simple dimensional analysis of the equations governing the electric field (no physical length scale is present). Henceforth we scale all spatial variables against a characteristic decay length without loss of generality. We consider three types of media generated from Gaussian random fields. The field-field correlation functions and their corresponding spectra are: Model I: $`g(r)`$ $`=`$ $`e^r{\displaystyle \frac{\mathrm{sin}\nu r}{\nu r}}`$ (29) $`\rho (k)`$ $`=`$ $`\pi ^2\left((1\nu ^2+k^2)^2+4\nu ^2\right)^1.`$ (30) Here $`\nu =2\pi l_1/l_2`$ with $`l_1`$ the decay length of the field and $`l_2`$ the characteristic domain size. When $`r0`$ the correlation functions arising from this model are similar to those considered in refs. ($`\nu =0`$) and ($`\nu >0`$). Note that this model has an infinite surface to volume ratio since $`<k^2>`$ diverges, however for computational realizations of the model the ratio is finite (see section V) and the model is well defined. However it is interesting to study $`\zeta _1`$ for this model to investigate the effect of interfacial roughness on effective properties. Model II: $`g(r)`$ $`=`$ $`e^{r^2}`$ (31) $`\rho (k)`$ $`=`$ $`{\displaystyle \frac{e^{\frac{1}{4}k^2}}{(4\pi )^{\frac{3}{2}}}}`$ (32) Model III: $`g(r)`$ $`=`$ $`{\displaystyle \frac{3\left(\mathrm{sin}\mu r\mu r\mathrm{cos}\mu r\mathrm{sin}rr\mathrm{cos}r\right)}{r^3(\mu ^31)}}`$ (33) $`\rho (k)`$ $`=`$ $`{\displaystyle \frac{3}{4\pi (\mu ^31)}}\left(H(\mu )H(1)\right)`$ (34) where $`\mu =k_1/k_0`$. Note that $`\rho (k)\delta (k)`$ as $`\mu 1`$ and the simple model used by Berk is recovered. To perform the integration (8) we use cylindrical co-ordinates (which damp the singularity at the origin), interchange the order of integration and exploit the $`rs`$ symmetry to give, $$\zeta _1=\frac{9}{2pq}_1^1\left(_0^{\mathrm{}}_0^{\frac{\pi }{4}}I(w,\varphi ,u)𝑑w𝑑\varphi \right)P_2(u)𝑑u$$ (35) with $$I(w,\varphi ,u)=\frac{pp_3(r,s,t)p_2(r)p_2(s)}{p^3w\mathrm{sin}\varphi \mathrm{cos}\varphi }.$$ (36) To elucidate the nature of the singularity in the integrand at $`w=\varphi =0`$ we consider $`w`$, $`\varphi `$ (and hence $`r`$, $`s`$ and $`t`$) to be small and assume the form $`g(r)1ar^n`$ (where $`n=1,2`$ in accord with models I-III). Now for $`p=\frac{1}{2}`$ the numerator of the expression for $`I`$ is given by, $$\frac{\mathrm{sin}^1(g(t))}{8\pi }\frac{\mathrm{sin}^1(g(r))\mathrm{sin}^1(g(s))}{4\pi ^2}\sqrt{2aw^n\varphi ^n},$$ (37) where we have used the results $`\mathrm{sin}^1(g(r))\frac{\pi }{2}\sqrt{2ar^n}`$, $`r=w\mathrm{cos}\varphi `$, $`s=w\mathrm{sin}\varphi `$ and $`t=w\sqrt{1\mathrm{sin}2\varphi u}`$. Therefore $`I(w,\varphi ,u)(w\varphi )^{\frac{n}{2}1}`$ with $`p=\frac{1}{2}`$ and numerical analysis shows this scaling also holds for $`p\frac{1}{2}`$. An integration rule which takes the singularity into account is employed. Note that for finite surface area to volume ratios $`n=2`$. The number of abscissae in each of three integration ranges was increased until the third significant figure in the estimation of $`\zeta _1`$ remained constant. The integration method was tested on the known correlation functions for the IOS model . The results are in exact agreement (to the reported 3rd significant figure) with those of Torquato and Stell and agree to the second significant figure with the results of Berryman . The calculation of the correlation functions in the integrand is done using a combination of iterative quadrature rules, the asymptotic results presented in appendix A and the non-rigorous approximation for the truncated three point function (28). The latter is used whenever $`1g<10^5`$ since in this case the functions $`F_{ij}`$ exhibit large derivatives and the quadrature rules converge too slowly. The accuracy sought in the application of each approximation is $`O(10^6)`$. The results for each of the models is presented in table I and plotted in figure 1 along side the results for the IOS model. Several comments on the qualitative relationship between $`\zeta _1`$ for the GRI model and prior calculations can be made. Consider the expansion, $$\zeta _1=\underset{i=0}{\overset{\mathrm{}}{}}e_ip^i.$$ (38) For a general class of materials with spherical inclusions it has been shown that $`e_0=0`$ . If the inverse of these materials is considered (so that the correlation functions refer to the material surrounding the inclusion) or ellipsoidal inclusions are considered then $`e_0>0`$ (see appendix C). For the case of symmetric cell materials $`e_0=M[0,1]`$ where $`M=0`$ for spherical cells and $`M=1`$ for plate like cells . Another interesting feature of $`\zeta _1`$ is that it is observed to be linear with $`p`$ over a wide range . Indeed for the symmetric cell model $`e_1=(12M)`$ and $`e_2,e_3\mathrm{}=0`$. By inspection of figure 1 we see that $`e_0>0`$ in qualitative agreement with the results for non-spherical inclusions which will occur in the GRI models. Also note that $`\zeta _1`$ is very similar to the results for the symmetric cell model for some $`0<M<1`$ over a wide range of $`p`$. This discussion demonstrates the success of simple models of random media to capture the qualitative features of $`\zeta _1`$ for the amorphous materials considered here. Calculations of the related microstructure parameter $`\eta _1`$ which arises in bounds on the elastic properties of random composites are reported in appendix E. ## V Generation of Fields For computational purposes we consider a $`T`$-periodic Gaussian random field with a maximum wavenumber $`K=2\pi N/T`$, $$y_{\text{ }K\text{ }}(𝐫)=\underset{l=N}{\overset{N}{}}\underset{m=N}{\overset{N}{}}\underset{n=N}{\overset{N}{}}c_{lmn}e^{i𝐤_{lmn}.𝐫}$$ (39) where, $$𝐤_{lmn}=\frac{2\pi }{T}(l𝐢+m𝐣+n𝐤).$$ (40) For $`y_{\text{ }K\text{ }}`$ real we require $`c_{l,m,n}=\overline{c}_{l,m,n}`$ and as $`<y_{\text{ }K\text{ }}>=0`$ we set $`c_{0,0,0}=0`$. For reasons which become clear below we also take $`c_{lmn}=0`$ for $`k_{lmn}=\|𝐤_{lmn}\|K`$. With $`c_{lmn}=a_{lmn}+ib_{lmn}`$, $`a_{lmn}`$ and $`b_{lmn}`$ are random independent variables (subject to the conditions on $`c_{lmn}`$) with Gaussian distributions such that $`<a_{lmn}>=<b_{lmn}>=0`$ and $$<a_{lmn}^2>=<b_{lmn}^2>=\frac{1}{2}\rho _{\text{ }K\text{ }}(k_{lmn})\left(\frac{2\pi }{T}\right)^3$$ (41) with $`\rho _{\text{ }K\text{ }}(k)`$ the spectral density. The field-field correlation function $`g_{\text{ }K\text{ }}`$ is given by $`g_{\text{ }K\text{ }}(r_{12})`$ $`=`$ $`<y_{\text{ }K\text{ }}(𝐫_1)\overline{y_{\text{ }K\text{ }}(𝐫_2)}>`$ (42) $`=`$ $`{\displaystyle \underset{N}{\overset{N}{}}}{\displaystyle \underset{N}{\overset{N}{}}}{\displaystyle \underset{N}{\overset{N}{}}}<c_{lmn}\overline{c_{lmn}}>e^{i𝐤_{lmn}.(𝐫_1𝐫_2)}`$ (43) $``$ $`{\displaystyle _0^K}4\pi k^2\rho _{\text{ }K\text{ }}(k){\displaystyle \frac{\mathrm{sin}k\|𝐫_1𝐫_2\|}{k\|𝐫_1𝐫_2\|}}𝑑k.`$ (44) The last integral is obtained by taking $`N`$ and $`T`$ large, using equation (41) and recognising that the summation is the approximation of a triple integral. Following a transformation to spherical co-ordinates we integrate over the angular variables to obtain (44). Since $`g_{\text{ }K\text{ }}(0)=1`$ we define, $$\rho _{\text{ }K\text{ }}(k)=P^1\rho (k),P=_0^K4\pi k^2\rho (k)𝑑k$$ (45) where the $`\rho (k)`$ are defined in section III. If in addition we take $`K\mathrm{}`$ then the conventional correlation function (and spectral density) are recovered. The Fourier expansion (39) is evaluated using a FFT. Consider materials derived from the field $`y_{\text{ }K\text{ }}`$ (eqn. (39)) as discussed in section III. Cross sections of the media for four different variants of the models are plotted in figure 2 (a)-(d). The large scale structure of the interface is determined by the terms in the expansion with small $`k`$ while the small scale structure (ripples on the surface) are determined by the terms where $`k`$ is large. A physical material will naturally contain a finite cutoff wavelength either imposed by the molecular size or by the manifestation of surface tension at the phase boundaries. This wavelength will then dictate the grid resolution necessary to properly resolve the structure for simulational purposes. Conversely for this study the discretisation $`\mathrm{\Delta }x=T/M`$, where $`M`$ is the number of grid points in the $`x`$ direction, will restrict the choice of $`K`$. Thus $`\mathrm{\Delta }x\lambda _{min}=2\pi /K`$ or $`NM`$. Another constraint is that $`Tl_d`$ where $`l_d`$ is an effective decay length of the field defined by $`g(l_d)=e^1`$. This must be so to ensure that the edges of the sample are uncorrelated to properly simulate an infinite random medium. For the approximation involved in determining $`g_{\text{ }K\text{ }}`$ (eqn. (44)) to be accurate would require $`T2\pi `$, however in this study this constraint was found to be of less importance and is ignored. The computational parameters used for each of the four variants of the models are given in table II. Cross-sections of each of the materials are plotted in figures 2 (e)-(h) and the interface for models II and III are plotted in figures 3 and 4. To properly account for the effect of $`K`$ for the computational models $`\zeta _1`$ must be recalculated using $`g_{\text{ }K\text{ }}(r)`$. This is done using a look up table generated by numerical integration of (44) and an asymptotic expansion for large $`r`$ (where the integral is costly to evaluate), $$g_{\text{ }K\text{ }}(r)=\frac{1}{P}g(r)4\pi K\rho _{\text{ }K\text{ }}(K)\mathrm{cos}Kr\frac{1}{r^2}+O\left(\frac{1}{r^3}\right).$$ (46) The function $`g_{\text{ }K\text{ }}(r)`$ for model I is plotted along with direct measurements of $`<y(𝐫_1)y(𝐫_2)>`$ in figure 5. The agreement is seen to very good. The values of $`\zeta _1`$ for the computational models are presented in table III. The effect of $`K`$ on $`\zeta _1`$ can be seen most clearly for Model I where the surface area to volume ratio becomes arbitrarily large as $`K`$ increases. The interfacial smoothing effect of imposing a finite cut-off ($`K`$) is shown in figures 2 (a) ($`K=18`$) and (e) ($`K=8`$). This has the effect of decreasing (increasing) the magnitude of $`\zeta _1`$ for $`p<\frac{1}{2}`$ ($`p>\frac{1}{2}`$) (compare tables I and III). For small (and large) volume fractions this can be qualitatively explained by the fact that the inclusion phase of the model will be much rougher (and hence less ‘spherical’) for $`K=\mathrm{}`$ than for $`K=8`$ (See appendix C). $`\zeta _1`$ for model II is unchanged for finite $`K`$ to the accuracy calculated here. ## VI Solution of the Laplace equation There are a variety of different methods of simulating the effective conductivity of an inhomogeneous medium. These include direct solution of the partial differential equations governing the potential $`\varphi `$ , Brownian motion algorithms, Fourier methods and other techniques. We solve Laplace’s equation with the charge conservation boundary conditions discussed in section II in a cube of side length $`T`$ using $`M^3`$ nodes. A potential of $`\varphi _1`$ and $`\varphi _0`$ are applied on the faces $`z=0`$ and $`z=T`$ and periodic boundary conditions are imposed on the four faces parallel to the direction of the current to model a sample of infinite extent in the $`x`$ and $`y`$ directions. A finite difference scheme is used to approximate the field and is solved using a conjugate gradient method. In appendix D we discuss the efficient implementation of the algorithm on a parallel computer. The $`z`$ components of the current and the field are then used in equation (4) to determine $`(\sigma _e)_M`$. The convergence criterion of the CG solver is decreased until the third significant figure of $`(\sigma _e)_M`$ remained constant. To estimate the continuum value of $`\sigma _e`$ we assume that $`\sigma _e(\sigma _e)_M+a_1M^1`$ and fit a line (using least squares) to several values of $`(\sigma _e)_M`$ vs. $`M^1`$. The intercept of this line with the axis $`M^1=0`$ then provides $`\sigma _e`$. Before proceeding to the random media we simulate the effective conductivity of a regular array of spheres of conductivity $`\sigma _1=10`$ in a matrix of conductivity $`\sigma _2=1`$. Exact results for this model have been calculated by McKenzie et al. and the model has been used by previous authors to test the accuracy of their algorithms. For computational purposes the array contains four spheres in the $`z`$ direction (using six spheres changes $`\sigma _e`$ by less than $`1\%`$). The values of $`(\sigma _e)_M`$ for increasing concentration are plotted along with the lines of best fit used to estimate $`\sigma _e`$ in figure 6. The graph demonstrates the necessity of extrapolating the data to $`M^10`$. The results for $`\sigma _e`$ are presented in table IV. The error is less than $`1\%`$ for $`p0.4`$ but increases to around $`3\%`$ at $`p=0.5`$ near the percolation threshold $`p_c0.52`$. For the random media it was found that computing $`(\sigma _e)_M`$ at more than two values of $`M`$ did not significantly alter the estimation of $`\sigma _e`$. For the random media we must consider how to assign the conductivity of a bond lying between two nodes which lie in different phases. Let $`y_{i,j}`$ and $`\sigma _{i,j}`$ be the respective values of the field and the conductivity at two such neighbouring nodes. There are three obvious ways of determining $`\sigma _{ij}`$. Defining $`a=(\alpha y_i)/(y_jy_i)`$ we can choose $`\sigma _{ij}=a\sigma _j+(1a)\sigma _i`$ (as if the portions of the volume element associated with the bond are like conductors in parallel), $`\sigma _{ij}=(a/\sigma _j+(1a)/\sigma _i)^1`$ (as if in series) or $`\sigma _{ij}=\sigma _1`$ or $`\sigma _2`$ as $`(y_i+y_j)/2>\alpha `$ or $`(y_i+y_j)/2<\alpha `$ (a simple field average). In figure 7 we show the effect of using these rules for two samples of material I generated with $`N=4,16`$. For a given discretisation ($`M`$) a large difference in $`\sigma _e`$ occurs depending on the rule employed. However extrapolation to $`M^1=0`$ demonstrates remarkably well that the choice is immaterial. As the simple averaging rule provides the least error for given $`M`$ it will be used. Finally, for a given volume fraction we use the bisection method to calculate the value of the level cut parameter $`\alpha ^{}`$ such that $`<\sigma >=p\sigma _1+q\sigma _2`$ (where $`<>`$ refers to bond averaging). This substantially reduces the statistical fluctuations in $`\sigma _e`$ compared to using the theoretical value of $`\alpha `$ determined from equation (13). ## VII Simulation Results We simulate the effective conductivity for the four GRI models listed in table II for a range compositions. As we are dealing with finite sized samples we report $`\sigma _e`$ as the average over a number of different realizations of the materials. Error bars, which represent 95% confidence limits on the results, are equal to twice the standard error. The samples are examined at three different contrast values; $`\sigma _{1,2}=10,1`$, $`\sigma _{1,2}=50,1`$ and $`\sigma _{1,2}=1,0`$. Previous authors have considered media with $`\sigma _{1,2}=\mathrm{},1`$ but this is not possible using the methods discussed here. The results for the case $`\sigma _{1,2}=10,1`$ are presented in table V and plotted in figure 8 for $`p[0.1,0.9]`$. To obtain the data five samples of each model with discretisations using $`M=64`$ and $`M=96`$ were considered. The results show little variation in $`\sigma _e`$ (for fixed $`p`$) amongst the four different media. The maximum relative difference of $`4.2\%`$ occurs at $`p=0.6`$ between $`\sigma _e`$ for model I ($`\nu =0`$, $`K=8`$) and model III ($`\mu =1.5`$). As the differences are relatively small we restrict further attention to the latter two materials. The bounds calculated from equations (6) and (9) using $`\zeta _1`$ from tables I and III are presented in table VI and plotted along with the simulation data in figures 9, for $`p[0.2,0.8]`$, and figure 10 for $`p[0.86,0.96]`$. The latter figure illustrates very clearly that the upper bound discriminates between model I and III. For the case $`\sigma _{1,2}=50,1`$ the results for model I ($`\nu =0`$, $`K=8`$) and model III ($`\mu =1.5`$) are reported in table VII and plotted along with relevant bounds in figure 11. Again $`5`$ samples of the media were considered with $`M=64`$ and $`M=96`$. Figure 11 shows very pronounced differences between the IOS model and GRI models. The results for the case $`\sigma _{1,2}=1,0`$ are given in table VIII and plotted along with the upper bound (the lower bound vanishes) in figure 12. Five samples at discretisations $`M=48`$ and $`M=64`$ were considered. There are several qualitative trends in the data which can be commented on. Note that for first two contrasts considered $`\sigma _e`$ of the GRI models is greater than that of the IOS model over the entire range of $`p`$. At low volume fractions this can be attributed to the fact that the inclusions of the GRI models are qualitatively less spherical than those of the IOS model (see appendix C). This can be clearly seen in figures 13 and 14 where the inclusions are plotted for each of the GRI models at $`p=0.07`$ (the IOS model will containing predominantly spherical inclusions at this volume fraction). At high volume fractions the situation is reversed; the matrix phase of the IOS model is extremely ramified and hence $`\sigma _e`$ is lower. Similarly near $`p=0,1`$ the small differences in $`\sigma _e`$ for Models I and III can be explained by the fact that the latter model has more spherical inclusions (compare figures 13 & 14). This behavior is consistent with the relative variations in $`\zeta _1`$ for each of the three models as discussed in appendix C. For mid-range $`p`$ the differences between the IOS and GRI models correspond to the fact that the more highly conducting regions of the latter are generally better connected than those of the IOS model. Again this difference can be anticipated from the relative behavior of the parameter $`\zeta _1`$ for the two classes of models. However this is not necessarily always so as can be seen by comparing the respective values of $`\zeta _1`$ (tables I & III) and $`\sigma _e`$ (table V) for Models I and III at $`p=0.4`$. In this case $`\sigma _e^\mathrm{I}<\sigma _e^{\mathrm{III}}`$ but $`\zeta _1^\mathrm{I}>\zeta _1^{\mathrm{III}}`$. Note that the simulation data for the IOS model in the case $`\sigma _{1,2}=1,0`$ were obtained for insulating spheres in a matrix of unit conductivity. Therefore the microstructure of the conducting phase is generally better connected than the GRI models and the arguments pertaining to the qualitative variations in $`\sigma _e`$ used above are reversed. In all of the above cases the data lies between the appropriate bounds. Furthermore the upper bound is seen to provide a good estimate for $`\sigma _e`$ for the GRI models for $`p0.7`$. This has been observed previously for both high and low $`p`$ (the latter case has not been considered in detail here). This fact provides evidence of the near optimality of bounds in these regions and a good confirmation of the techniques used in this paper (see figure 10 especially). We have calculated the percolation threshold for Model I and III using the algorithm of Skal et al. . In contrast to previous results we found $`p_c`$ to exhibit finite size effects and to depend on the spectrum. Averaging $`p_c`$ over 10 fields at $`M=20,32,48,64`$ and extrapolating to $`M=\mathrm{}`$ we found $`p_c0.07`$ ($`\alpha _c1.47`$) for Model I and $`p_c0.13`$ for Model III ($`\alpha _c1.13`$). In the absence of theoretical percolation results for the GRI model it is interesting to discuss the threshold in terms of the transition of the structures from elliptic to hyperbolic as $`p`$ increases. The average Gaussian curvature ($`K_G`$) for the GRI models is given by , $$<K_G>=\frac{1}{6}<k^2>(\alpha ^21).$$ (47) Therefore the nature of the interface is dominated by inclusions of positive curvature (eg. ellipsoids) for $`\|\alpha \|>1`$ and is predominantly hyperbolic (eg. bicontinuos) for $`\|\alpha \|<1`$. The fact that $`\alpha _c>1`$ indicates that connecting structures persist below the elliptic/hyperbolic transition as would be expected since $`<K_G>`$ is an average quantity. Finally it is interesting to discuss the surprisingly small variation in $`\sigma _e`$ (and $`\zeta _1`$) amongst the GRI materials. As can be seen in figures 2 (a)-(d) these materials appear to be very different when viewed at the same scale. However the major qualitative differences are related to the effective decay length, and when the materials are viewed at a scale proportional to this length they appear to be remarkably similar as shown in figures 2 (e)-(h). Equivalently if the length parameters are retained in eqns. (29)-(34) they can be tuned to achieve the latter group of figures at the same scale (without effecting $`\sigma _e`$). The remaining smaller qualitative differences account for the variation in $`\sigma _e`$ observed. ## VIII Conclusion We have investigated the effective conductivity of a two phase random composite material using bounding techniques and direct simulation. Our calculations of $`\zeta _1`$ increase the classes of composites to which the bounds can be applied, while the simulation data can be used to assess predictive theories and be compared with higher order bounds. The results also pertain to a variety of other effective properties of amorphous composites as discussed in section II. The bounds encompass all of the simulational data and the upper bound yields a reasonable estimate of $`\sigma _e`$ for $`p0.7`$ ($`\sigma _1>\sigma _2`$). Reasonably large differences in $`\sigma _e`$ and $`\zeta _1`$ are observed between the amorphous GRI models and the IOS model. This highlights the importance of incorporating microstructure effects in the calculation of the effective properties of composite materials. Conversely there is relatively little variation in $`\sigma _e`$ amongst the GRI models as the major qualitative microstructural differences between the models are related to an effective decay length upon which $`\sigma _e`$ is necessarily independent. We expect that other properties of such composites where no intrinsic field length scale is present (eg. the elastic bulk and shear moduli) will show similar behavior. It is clear that the Gaussian random interface model discussed here can serve as a useful ‘model’ amorphous medium in the study of the effective properties of random composites. Furthermore the bounds and simulations can be related to physical composites by experimentally relating such systems to one of the theoretically known models. This could be done by comparing the spectra obtained from small angle scattering studies with that obtained from the 2-point correlation functions of each model. Or, more simply, by comparing images of the models with electron micrographs. Although such schemes are only approximate, our results indicate that the fine microstructural details are relatively unimportant in both the calculation of $`\zeta _1`$ and the simulation of $`\sigma _e`$. We note that the GRI model can be extended to the case of membranes and foams and that higher order correlation functions can be calculated for use in more precise bounds. Random walker algorithms can be utilised to investigate the often studied scaling properties of $`\sigma _e`$ near the percolation threshold. ## ACKNOWLEDGMENTS The authors thank Stjepan Marčelja for suggesting the problem and M. Knackstedt, K. McGrath, P. Pieruschka, D. Singleton and X. Zhang for helpful discussions. The simulations reported here were carried out on the Thinking Machines CM-5 at the Australian National University Supercomputer Facility. A.R. is supported by an Australian postgraduate research award. ## A Asymptotic forms of $`p_2`$ The two point function (eqn. (14)) can be expanded in powers of $`\alpha `$ to yield, $$p_2(g)=\frac{1}{2\pi }\mathrm{arcsin}g+\frac{e^{\frac{1}{2}\alpha ^2}}{2\pi }\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^n\alpha ^{2n}a_n}{2^nn!}+p^2$$ (A1) with, $$a_n=\frac{2}{2n1}\left(1\left(\frac{1g}{1+g}\right)^{n\frac{1}{2}}\right)a_{n1}$$ (A2) and $`a_0=0`$. This expansion converges rapidly for small $`\|\alpha \|`$. For the case $`\alpha 0`$ we have (using successive integration by parts), $$p\frac{e^{\frac{1}{2}\alpha ^2}}{\sqrt{2\pi }\alpha }\left(1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^n1.3\mathrm{}(2n1)}{\alpha ^{2n}}\right).$$ (A3) Special care must be taken to determine a practical expansion for $`p_2`$. A simple application of Watson’s Lemma yields a solution which does not possess the correct limiting behavior as $`g0`$ and gives two different expansions for the cases $`g=1`$ and $`g<1`$. The first problem is dealt with by appropriately partitioning the integral, while the second necessitates a further transformation of variable, followed by an expansion in scaled parabolic cylinder functions (see for example ref. ). Thus we write $$p_2(g)=\frac{1}{2\pi }\left(_1^g_1^0\right)\mathrm{exp}\left(\frac{\alpha ^2}{1+t}\right)\frac{dt}{\sqrt{1t^2}}+p^2.$$ (A4) In the second of these integrals we make the substitution $`v=1/(1+t)1`$ to give $$\frac{e^{\alpha ^2}}{2^{\frac{3}{2}}\pi }_0^{\mathrm{}}\frac{e^{\alpha ^2v}}{(v+1)(v+\frac{1}{2})^{\frac{1}{2}}}𝑑v$$ (A5) which can be expanded using Watson’s Lemma, $$\frac{e^{\alpha ^2}}{2\pi }\left(\frac{1}{\alpha ^2}\frac{2}{\alpha ^4}+\frac{7}{\alpha ^6}+O\left(\frac{1}{\alpha ^8}\right)\right).$$ (A6) This is just the expansion of $`p^2`$ which can be cancelled from equation (A4). The remaining integral is put in a standard form by the substitution $`v=1/(1+t)1/(1+g)`$, $$p_2(g)=\frac{e^{\frac{\alpha ^2}{1+g}}}{2^{\frac{3}{2}}\pi }_0^{\mathrm{}}\frac{e^{\alpha ^2v}}{(v+\frac{1}{1+g})(v+\frac{1}{2}\frac{1g}{1+g})^{\frac{1}{2}}}𝑑v.$$ (A7) Note that the nature of the singularity of the integrand changes order as $`g1`$ (this makes it impossible to generate an expansion for the full range of $`g`$ using Watson’s Lemma). To develop a uniform expansion for large $`\alpha `$ valid near $`g=1`$ we make the further substitution, $$v=\frac{1}{2}u^2+\delta _gu,\delta _g=\sqrt{\frac{1g}{1+g}},$$ (A8) to give $$\frac{e^{\frac{\alpha ^2}{1+g}}}{2\pi }_0^{\mathrm{}}\frac{e^{\alpha ^2(\frac{1}{2}u^2+\delta _gu)}}{\left(\frac{1}{2}u^2+\delta _gu+\frac{1}{1+g}\right)}𝑑u.$$ (A9) In the usual way the non-exponential component of the integrand can be expanded in powers of $`u`$ and integrated term by term to give, $`p_2(g)`$ $``$ $`{\displaystyle \frac{e^{\frac{\alpha ^2}{1+g}}}{2\pi }}({\displaystyle \frac{1+g}{\alpha }}T_1\left(\delta _g\alpha \right){\displaystyle \frac{(1+g)^2}{\alpha ^2}}\delta _gT_2\left(\delta _g\alpha \right)`$ (A10) $`+`$ $`{\displaystyle \frac{(12g)(1+g)^2}{2\alpha ^3}}T_3\left(\delta _g\alpha \right))`$ (A11) where, $$T_n(z)=_0^{\mathrm{}}e^{zs\frac{1}{2}s^2}s^{n1}𝑑s=\mathrm{\Gamma }(n)e^{\frac{1}{4}z^2}\mathrm{U}(n\frac{1}{2},z).$$ (A12) $`\mathrm{U}(a,z)`$ is a parabolic cylinder function. Two simple checks can be made on this expansion. For $`g=1`$ we have $`T_n(0)=2^{n/21}\mathrm{\Gamma }(\frac{n}{2})`$ which, when substituted in (A11), gives the expansion of $`p`$. For $`g<1`$ we again employ Watson’s Lemma to determine the asymptotic expansion, $$T_n(z)=\underset{j=0}{\overset{\mathrm{}}{}}\frac{(1)^j\mathrm{\Gamma }(n+2j)}{j!2^jz^{n+2j}}.$$ (A13) Now taking $`g=0`$ in (A11) and using this expansion gives the asymptotic form of $`p^2`$ as it should. The expansions for the case $`\alpha 0`$ ($`p1`$) are simply, $`p=1\mathrm{AE}_1`$ and $`p_2(g)=2p1+\mathrm{AE}_2`$ where $`\mathrm{AE}_1`$ and $`\mathrm{AE}_2`$ are the asymptotic expansions given by (A3) and (A11) respectively with $`\alpha `$ replaced by $`\|\alpha \|`$. Note that Berk has derived a formal series representation of $`p_2`$ valid for all $`\alpha `$, however the convergence of the series is slow for $`g1`$ and not guaranteed at $`g=1`$. ## B Proof of $`\zeta _1=\frac{1}{2}`$ for $`p=\frac{1}{2}`$. In Brown’s formulation the parameter $`\zeta _1`$ arises as the limit as $`ϵ0`$ of the integral, $$\zeta _1=\frac{9}{2pq}_ϵ^{\mathrm{}}\frac{dr}{r}_ϵ^{\mathrm{}}\frac{ds}{s}_1^1𝑑uP_2(u)f(r,s,t)$$ (B1) where $`f`$ is the term in brackets of eqn. (8). Now by taking $`p=\frac{1}{2}`$ in eqn. (22) we have $`f=p_2^T(t)/22p_2^T(r)p_2^T(s)`$. Note that the integral of the second term vanishes since it does not depend on $`u`$ and $`_1^1P_2(u)𝑑u=0`$. Therefore after making the substitution $`t^2=r^2+s^22sru`$ eqn. (B1) becomes, $$\zeta _1=9_ϵ^{\mathrm{}}𝑑r_ϵ^r𝑑s_{rs}^{r+s}p_2^T(t)h(r,s,t)𝑑t$$ (B2) where $`h(r,s,t)=t(sr)^4(\frac{3}{4}(t^2s^2r^2)^2r^2t^2)`$. Interchanging the order of integration results in, $`\zeta _1`$ $`=`$ $`9\left({\displaystyle _0^{2ϵ}}+{\displaystyle _{2ϵ}^{\mathrm{}}}\right)p_2^T(t){\displaystyle _ϵ^{\mathrm{}}}𝑑s{\displaystyle _s^{t+s}}h𝑑r`$ (B3) $``$ $`9{\displaystyle _{2ϵ}^{\mathrm{}}}p_2^T(t){\displaystyle _ϵ^{t/2}}𝑑s{\displaystyle _s^{ts}}h𝑑r,`$ (B4) and carrying out the straight forward integrations over $`r`$ and $`s`$ leads to, $$\zeta _1=9_0^{2ϵ}p_2^T(t)\left(\frac{2t^2}{3ϵ^3}\frac{t^3}{2ϵ^4}+\frac{t^5}{24ϵ^6}\right)𝑑t.$$ (B5) Now by taking $`ϵ0`$ and using the fact that $`p_2^T(0)=\frac{1}{4}`$ gives $`\zeta _1=\frac{1}{2}`$. ## C Relationship between $`\sigma _e`$, $`\zeta _1`$ and shape at low $`p`$ Consider the small concentration approximation to $`\sigma _e`$ for the case of randomly distributed and oriented spheroids with axial ratios $`A,A,12A`$, $$\sigma _e=\sigma _2+\frac{1}{3}p(\sigma _1\sigma _2)z(\sigma _1,\sigma _2,A)+O(p^2)$$ (C1) where, $$z=\frac{2\sigma _2}{\sigma _2+A(\sigma _1\sigma _2)}+\frac{\sigma _2}{\sigma _2+(12A)(\sigma _1\sigma _2)}.$$ (C2) With $`A[0,1/2]`$ it can be easily shown that $`z`$ has a unique minimum at $`A=1/3`$ (spherical inclusions) and is monotonically increasing as $`\|A1/3\|`$ increases. Therefore with $`\sigma _1>\sigma _2`$, $`\sigma _e`$ will be higher the lower the sphericity of the inclusions (and conversely for $`\sigma _2>\sigma _1`$). The same argument should qualitatively hold for arbitrary shapes. To see how this relates to $`\zeta _1`$ we match the terms of the expansions of eqn. (C1) and eqn. (6) to order $`p`$ and $`(\sigma _1\sigma _2)^3`$ which gives $`\zeta _1=(13A)^2+O(p)`$. Thus $`\zeta _1`$ will be higher for less spherically shaped inclusions. ## D Implementation of the finite difference scheme The finite difference scheme (see for example ref. ) for the equations and boundary conditions discussed in section II leads to a system of simultaneous equations for the value of the potential at each of the interior nodes (including those on lateral faces if we define $`\varphi `$ to be periodic in the $`x`$ and $`y`$ directions). For each such node $`u`$ we have $$\underset{vnn}{}\sigma _{uv}(\varphi _u\varphi _v)=0$$ (D1) where $`nn`$ is the set of nearest neigbours of node $`u`$ and $`\sigma _{uv}`$ is the conductivity of the bond lying between nodes $`u`$ and $`v`$. Conventionally these equations are cast as a matrix equation with $`\varphi `$ and the boundary conditions as 1 dimensional matrices (vectors) . On a parallel computer it saves coding and implementation time to retain the potential $`\varphi _{i,j,k}`$ as a 3D matrix. Define $`A`$ as the operator which performs the operations defined on the left hand side of (D1) for interior nodes and $`A\varphi \varphi `$ on the nodes where Dirichlet conditions are to be applied. Also define $`b`$ as a 3D matrix containing the boundary conditions on the field ($`b_{i,j,1}=\varphi _1,b_{i,j,M}=\varphi _0`$ for all $`i,j`$ and $`b=0`$ elsewhere). Then solving the system of equations for $`\varphi `$ is equivalent to minimising $`A\varphi b_2`$, which can be done using a conjugate gradient method which handles vectors of general dimension. ## E Calculation of the microstructure parameter $`\eta _1`$ Three point bounds have also been derived for the elastic bulk and shear moduli which can be expressed in terms of $`\zeta _1`$ and an additional parameter, $`\eta _1`$ $`=`$ $`{\displaystyle \frac{5}{21}}\zeta _1+{\displaystyle \frac{150}{7pq}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dr}{r}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}{\displaystyle _1^1}duP_4(u)\times `$ (E2) $`\left(p_3(r,s,t){\displaystyle \frac{p_2(r)p_2(s)}{p}}\right).`$ We have calculated $`\eta _1`$ for several different GRI models and tabulated (see table IX) it along with data for the IOS model. Qualitatively the results are similar to those for $`\zeta _1`$ and we expect the differences in the effective shear and bulk moduli amongst the GRI models to be similar to those observed for the conductivity case.
warning/0004/hep-th0004176.html	ar5iv	text	# 1. Introduction ## 1. Introduction In this paper we discuss a large class of gauge theory models where the gauge fields are coupled to $`N_\mathrm{F}`$ Dirac fermions with different coupling constants (= charges) $`e_{j,\sigma }^\alpha `$ where the index $`\sigma =+,`$ distinguishes the chiral components of the fermions, $`j=1,2,\mathrm{},N_\mathrm{F}`$ is a fermion flavor index, and $`\alpha =1,2,\mathrm{}N_\mathrm{P}`$ labels different gauge fields. As we shall see, the existence of large gauge transformations implies that the charges have to be quantized i.e. $`e_{j,\sigma }^\alpha =n_{j,\sigma }^\alpha e^\alpha `$ (no summation) with $`n_{j,\sigma }^\alpha `$ integers. We say that such a model is chirally symmetric if for all $`\alpha `$ and $`j`$, $`n_{j,+}^\alpha =n_{\pi (j),}^\alpha `$ for some permutation $`\pi `$. It is worth noting that the standard Schwinger model is obtained as the special case $`N_\mathrm{P}=N_\mathrm{F}=1`$ and $`n_+=n_{}`$, and the chiral Schwinger model corresponds to $`N_\mathrm{P}=N_\mathrm{F}=1`$ and $`n_+=1`$ and $`n_{}=0`$ (for review on previous work on these and similar models see ). The class of models with only one photon field ($`N_\mathrm{P}=1`$) and coupling constants $`en_{j,+}`$ and $`en_{j,}`$ to the right- and left handed chiral components of the fermions, respectively ($`j=1,\mathrm{},N_\mathrm{F}`$, and the $`n_{j,\sigma }`$ are integers) was previously proposed and studied in . For these models it is known that the gauge anomaly cancels whenever $$\underset{j=1}{\overset{N_\mathrm{F}}{}}\left(n_{j,+}^2n_{j,}^2\right)=0.$$ (1) The simplest non-trivial (i.e. chirally asymmetric) solution is for $`N_\mathrm{F}=2`$, $`n_{j,+}=(3,4)`$ and $`n_{j,}=(0,5)`$, and therefore one sometimes refers to this model as the 3-4-5 model . Similarly one can find non-trivial examples for all Pythagorean Triple i.e. integer solutions of $`a^2+b^2=c^2`$. For the general class of models with $`N_\mathrm{F}`$ fermions and $`N_\mathrm{P}N_\mathrm{F}`$ photons which we study we find the following conditions for gauge anomalies to be absent, $$\underset{j=1}{\overset{N_\mathrm{F}}{}}\left(n_{j,+}^\alpha n_{j,+}^\beta n_{j,}^\alpha n_{j,}^\beta \right)=0\alpha ,\beta =1,\mathrm{},N_\mathrm{P}.$$ (2) A simple non-trivial example is for $`N_\mathrm{F}=N_\mathrm{P}=2`$, $`n_{j,+}^1=(3,4)`$, $`n_{j,}^1=(0,5)`$, $`n_{j,+}^2=(1,2)`$, and $`n_{j,}^2=(2,1)`$. To find solutions of these conditions for given $`N_\mathrm{P}`$ and as small $`N_\mathrm{F}`$ as possible seems to be an interesting generalization of the problem of finding Pythagorean Triples. The study of the 3-4-5 model in Ref. was in the path integral formalism. Here we use a Hamiltonian approach in the spirit of Ref. . This approach allows a rigorous construction of these model in the Hamiltonian framework using the quasi-free representation of boson- and fermion field algebras , and we also can solve the model using standard bosonization techniques . Our general class of models (i.e. arbitrary $`N_\mathrm{F}`$ and $`N_\mathrm{P}`$) is quite complicated, and it is somewhat surprising that it is possible to find the solution in such an explicit manner as we do in this paper. It is also intriguing to see the importance of the no-anomaly conditions in Eq. (2) at several different, seemingly unrelated points of our construction and solution. To simplify notation we explain our methods and computations in detail for the simplest case $`N_\mathrm{P}=1`$, and we are careful to do things such that the generalization to the case $`N_\mathrm{P}>1`$ is easy. The plan of this paper is as follows. We first concentrate on the models with one photon field and an arbitrary number $`N_\mathrm{F}`$ of fermion fields. After a formal definition we summarize the rigorous construction of these model in our framework. Anomalies are a consequence of Schwinger terms which result from the normal ordering necessary to construct the fermion currents as well-defined operators on the Hilbert space of the model. Especially we compute the gauge anomalies from the commutators of the implementers of gauge transformations (= Gauss’ law generators), and Eq. (1) is obtained as condition for a vanishing gauge anomaly. If and only if this latter condition holds a simple construction and solution of this model is obtained. As mentioned, these results for $`N_\mathrm{P}=1`$ are presented such that the generalization to our general class of models with $`N_\mathrm{P}`$ photons and $`N_\mathrm{F}`$ fermions is easy. The presentation of our results for this latter case are given in a rather short final paragraph. ## 2. Notation Throughout this paper we consider Abelian gauge theories with massless fermions on two dimensional spacetime which is a cylinder, $`R\times S^1`$, parametrized by coordinates $`(x^\mu )=(x^0,x^1)`$ where $`x^0=tR`$ (= time) and $`x^1=x[L/2,L/2]`$ (= spatial coordinate; $`0<L<\mathrm{}`$ is the spatial length). Our metric tensor is $`diag(1,1)`$. We use the Einstein summation convention only for spacetime indices $`\mu ,\nu =0,1`$ but not for flavor- or spin indices. ## 3. Formal definition of the model We now define in some detail the simplest non-trivial example for a chiral model without gauge anomalies. As described above, this model contains $`N_\mathrm{F}`$ flavors of Dirac fermion fields $`\psi _{j,\sigma }`$, $`\overline{\psi }_{j,\sigma }`$ and one Abelian gauge field $`A_\mu `$; here and in the following, $`\sigma ,\sigma ^{}=+,`$ are spin indices, and indices $`j,j^{}=1,\mathrm{},N_\mathrm{F}`$ distinguish the different fermion flavors. The model is formally defined by the Lagrangian $$=\frac{1}{4}F_{\mu \nu }F^{\mu \nu }+\underset{j=1}{\overset{N_\mathrm{F}}{}}\overline{\psi }_j\gamma ^\mu [\mathrm{i}_\mu +e(n_{j,+}P_++n_{j,}P_{})A_\mu ]\psi _j$$ (3) where $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$, and $$P_\pm =\frac{1}{2}(\mathrm{𝟏}\pm \gamma _3)$$ (4) are chiral projections; here $`\gamma ^\mu (\gamma ^\mu )_{\sigma ,\sigma ^{}}`$ are Dirac matrices which we take as $`\gamma ^0=\sigma _1`$, $`\gamma ^1=\mathrm{i}\sigma _2`$, and $`\gamma _3=\gamma ^0\gamma ^1=\sigma _3`$ ($`\sigma _{1,2,3}`$ are the Pauli spin matrices as usual), and the real parameters $`e_{j,\sigma }=en_{j,\sigma }`$ are coupling constants. The gauge group for this model is $`𝒢=C^{\mathrm{}}(R\times S^1;\mathrm{U}(1))`$ (= smooth $`\mathrm{U}(1)`$-valued functions on spacetime), and the Lagrangian Eq. (3) is obviously invariant under the following gauge transformations, $`\psi _j`$ $``$ $`(\text{e}^{\mathrm{i}n_{j,+}\chi }P_++\text{e}^{\mathrm{i}n_{j,}\chi }P_{})\psi _j`$ $`\overline{\psi }_j`$ $``$ $`\overline{\psi }_j(\text{e}^{\mathrm{i}n_{j,+}\chi }P_{}+\text{e}^{\mathrm{i}n_{j,}\chi }P_+)`$ $`A_\mu `$ $``$ $`A_\mu {\displaystyle \frac{1}{e}}_\mu \chi `$ (5) for all $`\text{e}^{\mathrm{i}\chi }𝒢`$. Note that the existence of the large gauge transformation $`\text{e}^{\mathrm{i}\chi (x,t)}=\text{e}^{\mathrm{i}x/L}`$ forces us to require that the $`n_{j,\sigma }`$ are integers (otherwise the large gauge transformations cannot be implemented in our model): The charges of the different fermion flavors have to be quantized. To motivate our construction of the model in the Hamiltonian framework below we recall some formulas from the formal canonical procedure; see e.g. . From the action Eq. (3) one computes the canonical momenta for the various fields of the model and obtains the following canonical (anti-) commutator relations<sup>1</sup><sup>1</sup>1$`\{a,b\}=ab+ba`$ and $`[a,b]=abba`$<sup>,</sup><sup>2</sup><sup>2</sup>2here and in the following we set $`t=0`$ and make explicit the dependence on the spatial coordinate only $`\{\psi _{j,\sigma }(x),\psi _{j^{},\sigma ^{}}^{}(y)\}`$ $`=`$ $`\delta _{\sigma ,\sigma ^{}}\delta _{j,j^{}}\delta (xy)`$ $`[E(x),A_1(y)]`$ $`=`$ $`\mathrm{i}\delta (xy)`$ (6) etc. where $`\psi ^{}=\overline{\psi }\gamma ^0`$ and $`E=F_{01}`$. Moreover, with the notation $`H_0`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N_\mathrm{F}}{}}}{\displaystyle _{L/2}^{L/2}}dx\psi _j^{}(x)\gamma _3(\mathrm{i}_1)\psi _j(x)`$ $`J(x)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N_\mathrm{F}}{}}}{\displaystyle \underset{\sigma =\pm }{}}n_{j,\sigma }\psi _j^{}(x)\gamma _3P_\sigma \psi _j(x)`$ $`\rho (x)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N_\mathrm{F}}{}}}{\displaystyle \underset{\sigma =\pm }{}}n_{j,\sigma }\psi _j^{}(x)P_\sigma \psi _j(x)`$ (7) the resulting Hamiltonian for the model can be written as $$H=H_0+_{L/2}^{L/2}dx\left(\frac{1}{2}E(x)^2+eA_1(x)J(x)\right)$$ (8) and has to be supplemented by the constraint (= Gauss’ law) $$G(x)=_1E(x)+e\rho (x)0.$$ (9) We note in passing that we also have two Noether currents, a vector current $`(J_\mathrm{V}^\mu )`$ and an axial current $`(J_\mathrm{A}^\mu )`$. These currents are given by $`J_\mathrm{V}^0=\rho ,J_\mathrm{V}^1=J`$ $`J_\mathrm{A}^0=J,J_\mathrm{A}^1=\rho `$ (10) and formally (i.e. prior to quantization) obey continuity equations, $`_\mu J_\mathrm{V}^\mu =_\mu J_\mathrm{A}^\mu =0`$. We also introduce more general kinds of fermion currents $$\rho _{j,\sigma }(x)=\psi _{j,\sigma }^{}(x)\psi _{j,\sigma }(x)$$ (11) which all are observables of interest for our model. ## 4. Construction of the model We now outline a rigorous construction of this model using the representation theory of loop groups . This construction amounts to representing the (Fourier modes) of the field operators $`\psi _{j,\sigma }^{()}`$, $`A_1`$ and $`E`$, and the observable algebra of the model by closed operators on a Hilbert space $``$ such that the Hamiltonian $`H`$ is represented by a self-adjoint operator on $``$. In this construction it is crucial to establish gauge- and Lorentz invariance. As we shall see, this will lead us to the condition in Eq. (1). The essential physical requirement determining the construction of the model and implying a non-trivial vacuum structure is positivity of the Hamiltonian $`H`$ on the physical Hilbert space. As is well-known, this forces one to use a non-trivial representation of the field operators of the model. The essential simplification in (1+1) (and not possible in higher) dimensions is that one can use a quasi-free representation for the fermion field operators corresponding to “filling up the Dirac sea” associated with the free fermion Hamiltonian $`H_0`$, and for the photon operators one can use a naive boson representation. We now describe this representation in more detail. In the following it is convenient to work in Fourier space. We first introduce some useful notation which is such that in all equations the limit $`L\mathrm{}`$ is obvious. The Fourier space for even (periodic) and odd functions is $`\mathrm{\Lambda }^{}\left\{p={\displaystyle \frac{2\pi }{L}}n\right\|nZ\}\text{ and }\mathrm{\Lambda }_{odd}^{}\left\{q={\displaystyle \frac{2\pi }{L}}\left(n+{\displaystyle \frac{1}{2}}\right)\right\|nZ\},`$ respectively. For functions $`\widehat{f}`$ on Fourier space we write $`\widehat{}_\mathrm{\Lambda }^{}\widehat{\mathrm{d}}p\widehat{f}(p)_{p\mathrm{\Lambda }^{}}\frac{2\pi }{L}\widehat{f}(p)`$ and similarly for $`\mathrm{\Lambda }_{odd}^{}`$, and the corresponding $`\delta `$-function satisfying $`\widehat{}_\mathrm{\Lambda }^{}\widehat{\mathrm{d}}p^{}\widehat{\delta }(pp^{})\widehat{f}(p^{})=\widehat{f}(p)`$ is $`\widehat{\delta }(pp^{})\frac{L}{2\pi }\delta _{p,p^{}}`$. Our conventions for the Fourier transformed operators are, $$\widehat{\psi }_{j,\sigma }^{()}(q)=_{L/2}^{L/2}\frac{\mathrm{d}x}{\sqrt{2\pi }}\psi _{j,\sigma }^{()}(x)\text{e}^{\stackrel{(+)}{}\mathrm{i}qx}(q\mathrm{\Lambda }_{odd}^{})$$ (12) (i.e. we use anti-periodic boundary conditions for the fermions), $$\widehat{A}_1(p)=_{L/2}^{L/2}\frac{\mathrm{d}x}{2\pi }A_1(x)\text{e}^{\mathrm{i}px}(p\mathrm{\Lambda }^{})$$ (13) and in the other cases $$\widehat{Y}(p)=_{L/2}^{L/2}dxY(x)\text{e}^{\mathrm{i}px}(p\mathrm{\Lambda }^{})\text{ for }Y=E,\rho ,J,\rho _{j,\sigma }\text{ etc.}$$ (14) Therefore the non-trivial canonical (anti-) commutator relations of the field operators become $`[\widehat{A}_1(p),\widehat{E}(p^{})]=\mathrm{i}\widehat{\delta }(p+p^{})`$ and $`\{\widehat{\psi }_{j,\sigma }(q),\widehat{\psi }_{j^{},\sigma ^{}}^{}(q^{})\}`$ $`=`$ $`\delta _{\sigma ,\sigma ^{}}\delta _{j,j^{}}\widehat{\delta }(qq^{})`$ $`\{\widehat{\psi }_{j,\sigma }^{()}(q),\widehat{\psi }_{j^{},\sigma ^{}}^{()}(q^{})\}`$ $`=`$ $`0p,p^{},j,j^{}.`$ (15) The model will be constructed in the full Hilbert space $`=_{\mathrm{Photon}}_{\mathrm{Fermion}}`$. For $`_{\mathrm{Photon}}`$ we take the Fock space generated by boson field operators $`b^{}(p)`$ and $`b^{}(p)=b(p)^{}`$, $`p\mathrm{\Lambda }^{}`$, ($``$ is the Hilbert space adjoint) obeying the commutator relations $$[b(p),b^{}(p^{})]=\widehat{\delta }(pp^{}),[b(p),b(p^{})]=0p,p^{}$$ (16) and a normalized state $`\mathrm{\Omega }_\mathrm{P}_{\mathrm{Photon}}`$ such that $$b(p)\mathrm{\Omega }_\mathrm{P}=0p.$$ (17) We then set $$\widehat{A}_1(p)=\frac{1}{s}\left(b(p)+b^{}(p)\right)\widehat{E}(p)=\frac{\mathrm{i}s}{2}\left(b(p)b^{}(p)\right)$$ (18) with a parameter $`s`$ to be determined later. We recall that these requirements fix the Hilbert space $`_{\mathrm{Photon}}`$ completely. In this setting we now construct bilinears in the photon field using normal ordering $`:\mathrm{}:`$ with respect to the state $`\mathrm{\Omega }_\mathrm{P}`$, e.g. $`:b(p)b^{}(p^{}):=b^{}(p^{})b(p)`$ for all $`p,p^{}`$. For $`_{\mathrm{Fermion}}`$ we take the Fermion Fock space generated by operators $`\widehat{\psi }_{j,\sigma }(q)`$ and $`\widehat{\psi }_{j,\sigma }^{}(q)=\widehat{\psi }_{j,\sigma }(q)^{}`$ obeying the relations Eq. (4. Construction of the model) and a normalized state $`\mathrm{\Omega }_\mathrm{F}_{\mathrm{Fermion}}`$ such that $`\widehat{\psi }_{+,j}(q)\mathrm{\Omega }_\mathrm{F}=\widehat{\psi }_{,j}^{}(q)\mathrm{\Omega }_\mathrm{F}=0q>0,j.`$ (19) We note that the state $`\mathrm{\Omega }_\mathrm{F}`$ characterizing this representation of the fermion field algebra can be interpreted as Dirac sea associated with the free fermion Hamiltonian $`H_0`$. The presence of this Dirac sea requires normal-ordering $`:\mathrm{}:`$ of the Fermion bilinears such as $`H_0=_j\widehat{}_{\mathrm{\Lambda }_{odd}^{}}\widehat{\mathrm{d}}q:q\widehat{\psi }_j^{}(q)\gamma _3\widehat{\psi }_j(q):`$ and $`\widehat{\rho }_{j,\sigma }(p)`$. This modifies the naive commutator relations of these operators which follow from the canonical anti-commutator relations Eq. (4. Construction of the model) as Schwinger terms show up, see e.g. . The relevant commutators for us are, $`[\widehat{\rho }_{j,\sigma }(p),\widehat{\rho }_{j^{},\sigma ^{}}(p^{})]`$ $`=`$ $`\sigma \delta _{\sigma ,\sigma ^{}}\delta _{j,j^{}}p\widehat{\delta }(p+p^{})`$ $`[H_0,\widehat{\rho }_{j,\sigma }(p)]`$ $`=`$ $`\sigma p\widehat{\rho }_{j,\sigma }(p).`$ (20) We also note that $$\widehat{\rho }_{j,+}(p)\mathrm{\Omega }_\mathrm{F}=\widehat{\rho }_{j,}(p)\mathrm{\Omega }_\mathrm{F}=0p>0,j$$ (21) which together with (4. Construction of the model) shows that the $`\widehat{\rho }_{j,+}(p)`$ (resp. $`\widehat{\rho }_{j,}(p)`$) give a highest (resp. lowest) weight representation of the Heisenberg algebra. The (Fourier transformed) Gauss’ law operators are now well-defined on $``$, $$\widehat{G}(p)=\mathrm{i}p\widehat{E}(p)+e\widehat{\rho }(p)$$ (22) with $`\widehat{\rho }(p)=_{j,\sigma }n_{j,\sigma }\widehat{\rho }_{j,\sigma }(p)`$. To determine a (possible) gauge anomaly we compute the commutators of these operators and obtain (we use Eq. (4. Construction of the model)) $$[\widehat{G}(p),\widehat{G}(p^{})]=e^2p\widehat{\delta }(p+p^{})\underset{j=1}{\overset{N_\mathrm{F}}{}}\underset{\sigma =\pm }{}\sigma n_{j,\sigma }^2.$$ (23) We see that no gauge anomaly occurs if and only if the condition Eq. (1) holds. In the following we assume that this is the case. We now can give a precise meaning to the formal Hamiltonian $`H`$ in Eq. (8) as follows, $$H=H_0+\widehat{}_\mathrm{\Lambda }^{}\widehat{\mathrm{d}}p:\left(\frac{1}{4\pi }\widehat{E}(p)\widehat{E}(p)+e\widehat{A}_1(p)\widehat{J}(p)+\pi M^2\widehat{A}_1(p)\widehat{A}_1(p)\right):.$$ (24) Similarly as in the Schwinger model we have added a photon mass term (= last term on the r.h.s. of Eq. (24)) to restore gauge invariance which otherwise would be spoiled by Schwinger terms. Note that our normalization is such that this term formally equals $`\frac{1}{2}M^2_{L/2}^{L/2}dxA_1(x)^2`$, i.e. $`M`$ can be interpreted as photon mass. To determine the parameter $`M^2`$ we compute $`[\widehat{G}(p),H]`$ and demand this to be always zero. By straightforward computation (we use Eq. (4. Construction of the model)), $$M^2=\frac{e^2}{2\pi }\underset{j=1}{\overset{N_\mathrm{F}}{}}\underset{\sigma =\pm }{}n_{j,\sigma }^2.$$ (25) We finally have to fix the parameter $`s`$ in Eq. (18). We observe that $`H`$ in Eq. (24) has the form $`H=H_{0,\mathrm{F}}+H_{0,\mathrm{P}}+H_{\mathrm{int}}`$ where $`H_{0,\mathrm{F}}=H_0`$ and $$H_{0,\mathrm{P}}=\widehat{}_\mathrm{\Lambda }^{}\widehat{\mathrm{d}}p:\left(\frac{1}{4\pi }\widehat{E}(p)\widehat{E}(p)+\pi M^2\widehat{A}_1(p)\widehat{A}_1(p)\right):$$ (26) are the free fermion- and Photon Hamiltonians. We fix $`s`$ by requiring that $`H_{0,\mathrm{P}}`$ has the form $`\widehat{}_\mathrm{\Lambda }^{}\widehat{\mathrm{d}}p\omega (p)b^{}(p)b(p)`$. We thus obtain $`s=2\sqrt{\pi \|M\|}`$ and $`\omega (p)=\|M\|`$. Then $`H_{0,\mathrm{F}}+H_{0,\mathrm{P}}`$ obviously is defined as self-adjoint, positive operators on $``$. Our results in this paper imply that also the interacting Hamiltonian $`H`$ is well-defined on $``$: the operator $`H`$ is self-adjoint and bounded from below on $``$. ## 5. Bosonization and gauge fixing I As in the standard Schwinger model we can rewrite the Hamiltonians of our models as boson Hamiltonians by using the so-called Kronig identity $`\widehat{H}_0=\frac{1}{2}\widehat{}_\mathrm{\Lambda }^{}\widehat{\mathrm{d}}p:_{j,\sigma }\widehat{\rho }_{j,\sigma }(p)\widehat{\rho }_{j,\sigma }(p):`$ where $`:\mathrm{}:`$ denotes normal ordering of the fermion currents (see for more details). We obtain $$H=\widehat{}_\mathrm{\Lambda }^{}\widehat{\mathrm{d}}q:\left(\frac{1}{4\pi }\widehat{E}(p)\widehat{E}(p)+\frac{1}{2}\underset{j=1}{\overset{N_\mathrm{F}}{}}\underset{\sigma =\pm }{}\stackrel{~}{\rho }_{j,\sigma }(p)\stackrel{~}{\rho }_{j,\sigma }(p)\right):$$ (27) where we have defined $$\stackrel{~}{\rho }_{j,\sigma }(p)=\widehat{\rho }_{j,\sigma }(p)+\sigma en_{j,\sigma }\widehat{A}_1(p).$$ (28) Note that $`[\widehat{G}(p),\stackrel{~}{\rho }_{j,\sigma }(p^{})]=0`$ i.e. the operators $`\stackrel{~}{\rho }_{j,\sigma }(p)`$ all are gauge invariant and thus observables for our model. The Hilbert space $``$ which we have obtained still contains gauge equivalent states i.e. states related to each other by static gauge transformations. To obtain the physical Hilbert space $`_{\mathrm{phys}}`$ for our models we have to go through a procedure called fixing the gauge in the physics literature. We will do it in two steps. The first step will be done in this paragraph and eliminates the ‘small’ gauge transformations (i.e. the $`\text{e}^{\mathrm{i}\chi }𝒢`$ with $`\chi `$ smooth and periodic functions on space $`S^1`$). This step is easy since the resulting Hilbert space $`_{\mathrm{phys}}^{}`$ can be identified as the sub-Hilbert space of $``$ which is annihilated by all the Gauss law operators. In a second step, which will be described in Paragraph 7 further below, we will account for large gauge transformations, too. For the first step we recall that the only gauge invariant degree of freedom of the Photon field at fixed time is the holonomy $`_{L/2}^{L/2}dxA_1(x)`$. Due to the absence of a gauge anomaly we can therefore impose the gauge condition $$\widehat{A}_1(p)=\delta _{p,0}Y$$ (29) and solve the Gauss’ law $`\widehat{G}(p)0`$ (cf. Eq. (22)) as $$\widehat{E}(p)=\frac{e\widehat{\rho }(p)}{\mathrm{i}p}\text{for }p0.$$ (30) This determines all components of $`\widehat{E}`$ except the one conjugate to $`Y`$: $$\widehat{E}(0)=\frac{L}{2\pi \mathrm{i}}\frac{}{Y}.$$ (31) After that we are left with the ($`p=0`$)-component of Gauss’ law, viz. $$eQ_\mathrm{V}0,Q_\mathrm{V}=\widehat{\rho }(0)=\underset{j=1}{\overset{N_\mathrm{F}}{}}\underset{\sigma =\pm }{}n_{j,\sigma }\widehat{\rho }_{j,\sigma }(0).$$ (32) We thus can identify the Hilbert space of states invariant under all small gauge transformations as $`_{\mathrm{phys}}^{}=L^2(R,\mathrm{d}Y)_{\mathrm{Fermion}}^{}`$ where $`_{\mathrm{Fermion}}^{}`$ is the zero charge sector of the fermionic Fock space. Moreover, by inserting Eqs. (29)–(31), each gauge invariant operator on $``$ becomes an operator on $`_{\mathrm{phys}}^{}`$. Especially the Hamiltonian becomes $`H=\frac{2\pi }{L}_{p0}:h_p:`$ where $$h_0=\frac{1}{4\pi }\left(\frac{L}{2\pi }\right)^2\frac{^2}{Y^2}+\frac{1}{2}\underset{j=1}{\overset{N_\mathrm{F}}{}}\underset{\sigma =\pm }{}\left(\widehat{\rho }_{j,\sigma }(0)+\sigma n_{j,\sigma }eY\right)^2$$ (33) and $`:h_p:=:({\displaystyle \frac{e^2}{2\pi p^2}}\widehat{\rho }(p)\widehat{\rho }(p)+{\displaystyle \underset{j=1}{\overset{N_\mathrm{F}}{}}}{\displaystyle \underset{\sigma =\pm }{}}\widehat{\rho }_{j,\sigma }(p)\widehat{\rho }_{j,\sigma }(p)):.`$ (34) This completes the first step of the gauge fixing procedure. As mentioned we have not yet accounted for the existence of the large gauge transformations. We will come back to this in Paragraph 7 below. ## 6. Solution of the model The operators $`:h_p:`$ obviously all commute with each other. We thus can diagonalize the Hamiltonian $`H`$ by diagonalizing each term $`:h_p:`$ separately. We first consider $`:h_p:`$ with $`p>0`$. We introduce the boson operators $$c_j(p)=\{\begin{array}{cc}\frac{1}{\sqrt{\|p\|}}\widehat{\rho }_{j,+}(p)& \text{ for }p>0\\ \frac{1}{\sqrt{\|p\|}}\widehat{\rho }_{j,}(p)& \text{ for }p<0\end{array}$$ (35) and $`c_j^{}(p)=c_j(p)^{}=c_j(p)`$, so that $`[c_j(p),c_j^{}^{}(p^{})]=\frac{L}{2\pi }\delta _{j,j^{}}\widehat{\delta }(pp^{})`$ and $`c_j(p)\mathrm{\Omega }=0`$ for all $`j,j^{}`$ and $`p,p^{}0`$. We also find it convenient to define $$c_j(p)=c_j(p),c_{N_\mathrm{F}+j}(p)=c_j^{}(p)j=1,\mathrm{}N_\mathrm{F},$$ (36) to fix $`p>0`$, and suppress the argument $`p`$ in the following. Then $`[c_j,c_j^{}^{}]=\frac{L}{2\pi }\delta _{j,j^{}}q_j`$ where $`q_j=1`$ and $`q_{N_\mathrm{F}+j}=1`$ for $`j=1,\mathrm{},N_\mathrm{F}`$. We then can write $`:h_p:`$ in matrix notation as $`:𝐜^{}\mathrm{𝐡𝐜}:`$ where we defined $`𝐜^{}=(c_1^{},\mathrm{},c_{2N_\mathrm{F}}^{})`$ etc. and $`𝐡`$ is the $`2N_\mathrm{F}\times 2N_\mathrm{F}`$ matrix $$𝐡=p\mathrm{𝟏}+\frac{e^2}{2\pi p}𝐧𝐧^T$$ (37) where $`𝐧^T=(n_1,\mathrm{},n_{2N_\mathrm{F}})`$ with $`n_j=n_{j,+}`$ and $`n_{N_\mathrm{F}+j}=n_{j,}`$ for $`j=1,\mathrm{},N_\mathrm{F}`$ ($`\mathrm{𝟏}`$ is the $`2N_\mathrm{F}\times 2N_\mathrm{F}`$ unit matrix here, and we use a standard tensor notation; the <sup>T</sup> denotes the matrix transpose). Defining $`h_p=𝐜^{}\mathrm{𝐡𝐜}`$ we get $`:h_p:=h_p<\mathrm{\Omega },h_p\mathrm{\Omega }>`$ with $$<\mathrm{\Omega },h_p\mathrm{\Omega }>=\frac{L}{2\pi }\underset{j=1}{\overset{N_\mathrm{F}}{}}\left(p+\frac{e^2}{2\pi p}n_{j,}^2\right)=\frac{L}{2\pi }\left(N_\mathrm{F}\|p\|+\frac{M^2}{2p}\right)$$ (38) (we used Eqs. (25) and (1)). We now can diagonalize $`h_p`$ by finding a boson Bogoliubov transformation $`𝐂=\mathrm{𝐔𝐜}`$ ($`𝐔`$ some $`2N_\mathrm{F}\times 2N_\mathrm{F}`$ matrix) so that the operators $`C_j`$ obey the same relations as the $`c_j`$ and are such that $`h_p=𝐂^{}\mathrm{𝐃𝐂}=_{j=1}^{2N_\mathrm{F}}\omega _jC_j^{}C_j`$. It is easy to see that these two conditions are equivalent to $$\mathrm{𝐔𝐪𝐔}^{}=𝐪\text{ and }(𝐔^1)^{}\mathrm{𝐡𝐔}^1=𝐃$$ (39) where $`𝐪=diag(q_1,\mathrm{},q_{2N_\mathrm{F}})`$ and $`𝐃=diag(\omega _1,\mathrm{},\omega _{2N_\mathrm{F}})`$ ( and <sup>-1</sup> is the matrix adjungation and matrix inverse, respectively). To solve this somewhat unconventional diagonalization problem we note that Eq. (39) implies $$𝐃^2=(\mathrm{𝐪𝐃})^2=𝐔(\mathrm{𝐪𝐡})^2𝐔^1,$$ (40) and $`𝐔`$ is determined from this equation up to transformations $`𝐔\mathrm{𝐕𝐔}`$ where $`𝐕`$ commutes with $`𝐃^2`$. We shall see that Eq. (40) corresponds to a standard diagonalization problem: the matrix $`(\mathrm{𝐪𝐡})^2`$ is self-adjoint and thus can be diagonalized by a unitary matrix $`\stackrel{~}{𝐔}`$. We thus can solve the problem in Eq. (39) by first determining $`\stackrel{~}{𝐔}`$ and then making the ansatz $`𝐔=𝐕\stackrel{~}{𝐔}`$ with $`𝐕`$ commuting with $`𝐃^2`$. From Eq. (39) we then get the following condition, $$𝐕\stackrel{~}{𝐔}\mathrm{𝐪𝐡}\stackrel{~}{𝐔}^{}𝐕^1=\mathrm{𝐪𝐃}$$ (41) which again is a standard diagonalization problem and will allow us to determine $`𝐕`$. We now compute $`𝐃^2`$ using Eq. (40). We write $`\mathrm{𝐪𝐡}=p𝐪+\frac{M^2}{p}𝐞_{N_\mathrm{F}+1}𝐞_1^T`$ where we define $`𝐞_1=𝐧/\|𝐧\|`$ ($`\|𝐧\|=\sqrt{𝐧^T𝐧}`$) and $`𝐞_{N_\mathrm{F}+1}:=\mathrm{𝐪𝐞}_1`$ (we used $`\|𝐧\|=\|\mathrm{𝐪𝐧}\|`$ and $`\frac{e^2}{2\pi }\|𝐧\|^2=M^2`$). It is now crucial to observe that the condition Eq. (1) is equivalent to $`𝐧^T\mathrm{𝐪𝐧}=0`$ i.e. if there is no gauge anomaly the two vectors $`𝐞_1`$ and $`𝐞_{N_\mathrm{F}+1}`$ are orthonormal. Moreover, in this case we can extend these vectors to a complete orthonormal real basis $`\{𝐞_j\}_{j=1}^{2N_\mathrm{F}}`$ in $`R^{2N_\mathrm{F}}`$ so that $`\mathrm{𝐪𝐞}_j=𝐞_{N_\mathrm{F}+j}`$ for $`j=1,\mathrm{},N_\mathrm{F}`$. We thus obtain $$(\mathrm{𝐪𝐡})^2=p^2\mathrm{𝟏}+M^2\left(𝐞_1𝐞_1^T+𝐞_{N_\mathrm{F}+1}𝐞_{N_\mathrm{F}+1}^T\right)$$ from which we can immediately read off $`\stackrel{~}{𝐔}`$ and the matrix elements of $`𝐃^2`$: denoting as $`𝐄_j`$ the standard basis in $`R^{2N_\mathrm{F}}`$ (i.e. $`(𝐄_j)_j^{}=\delta _{j,j^{}}`$) we have $$\stackrel{~}{𝐔}=\underset{j=1}{\overset{2N_\mathrm{F}}{}}𝐄_j𝐞_j^T,$$ (42) and $`\omega _j^2=p^2+M^2`$ and for $`j=1,N_\mathrm{F}+1`$ and $`\omega _j^2=p^2`$ otherwise. We then compute $$\stackrel{~}{𝐔}\mathrm{𝐪𝐡}\stackrel{~}{𝐔}^{}=\underset{j=1}{\overset{N_\mathrm{F}}{}}\left(\frac{M_j^2}{p}𝐄_{N_\mathrm{F}+j}𝐄_j^T+p𝐄_{N_\mathrm{F}+j}𝐄_j^T+p𝐄_j𝐄_{N_\mathrm{F}+j}^T\right)$$ (43) with $`M_1=M`$ and $`M_{j1}=0`$, which shows that we can determine $`𝐕`$ by diagonalizing $`2\times 2`$ matrices of the form $`\left(\begin{array}{cc}0& p\\ p+\frac{M_j^2}{p}& 0\end{array}\right)`$. We find $$𝐕=\underset{j=1}{\overset{2N_\mathrm{F}}{}}𝐄_j𝐅_j^T$$ (44) where $$𝐅_{j,N_\mathrm{F}+j}=\frac{p𝐄_j\pm \sqrt{p^2+M_j^2}𝐄_{N_\mathrm{F}+j}}{\sqrt{2p^2+M_j^2}}.$$ (45) Thus $$\omega _j=\omega _{N_\mathrm{F}+j}=\sqrt{p^2+M_j^2}\text{ for }j=1,\mathrm{},N_\mathrm{F}.$$ (46) We thus obtain $$h_p=\underset{j=1}{\overset{N_\mathrm{F}}{}}\sqrt{p^2+M_j^2}\left(C_j(p)^{}C_j(p)+C_j(p)C_j(p)^{}\right)$$ (47) where we used $`C_{N_\mathrm{F}+j}(p)=C_j^{}(p)`$. This completes the diagonalization of the operators $`:h_p:`$. We are left to diagonalize $`h_0`$. We note that this is the part of the Hamiltonian which contains the operators $`Y`$ and $`\widehat{\rho }_{j,\sigma }(0)`$ which are not invariant under the large gauge transformation $`\text{e}^{\mathrm{i}x/L}`$ but transform as follows, $`Y`$ $``$ $`Y{\displaystyle \frac{1}{e}}`$ $`\widehat{\rho }_{j,\sigma }(0)`$ $``$ $`\widehat{\rho }_{j,\sigma }(0)+\sigma n_{j,\sigma }.`$ (48) It is therefore not immediately obvious that $`h_0`$ is indeed invariant under the large gauge transformations. To make this invariance manifest we note that the operator $`Q_\mathrm{A}=_{j,\sigma }\sigma n_{j,\sigma }\widehat{\rho }_{j,\sigma }(0)`$ changes under the transformations Eq. (6. Solution of the model) as $`Q_\mathrm{A}Q_\mathrm{A}+\|𝐧\|^2`$, and therefore $$\stackrel{~}{Y}:=Y+\frac{1}{e\|𝐧\|^2}Q_\mathrm{A}$$ (49) is indeed invariant. By straightforward computations we obtain $$h_0=\frac{\|M\|}{2}\left(C_1(0)^{}C_1(0)+C_1(0)C_1(0)^{}\right)+𝒞$$ (50) where $`C_1^{()}(0)=\stackrel{()}{+}\frac{L}{4\pi }(\pi \|M\|)^{1/2}\frac{}{\stackrel{~}{Y}}+(\pi \|M\|)^{1/2}\stackrel{~}{Y}`$ and $$𝒞=\frac{1}{2}\left(\underset{j=1}{\overset{N_\mathrm{F}}{}}\underset{\sigma =\pm }{}\widehat{\rho }_{j,\sigma }(0)^2\frac{1}{\|𝐧\|^2}Q_\mathrm{A}^2\right)$$ (51) all are invariant under the transformations Eq. (6. Solution of the model). We also have $`[C_1(0),C_1^{}(0)]=\frac{L}{2\pi }`$ and thus see that $`h_0`$ is essentially a harmonic oscillator Hamiltonian. We can combine our results above in a compact form as follows, $$H=\underset{j=1}{\overset{N_\mathrm{F}}{}}\widehat{}_\mathrm{\Lambda }^{}\widehat{\mathrm{d}}p\sqrt{p^2+M_j^2}C_j^{}(p)C_j(p)+\frac{2\pi }{L}𝒞+_0,$$ (52) with the constant $$_0=\underset{j=1}{\overset{N_\mathrm{P}}{}}\left(\frac{\|M_j\|}{2}+\frac{L}{2\pi }\widehat{}_{p>0}\widehat{\mathrm{d}}p\left(\sqrt{p^2+M_j^2}(p+\frac{M_j^2}{2p})\right)\right)$$ (53) (note that actually only the term with $`j=1`$ contributes to $`_0`$). This shows explicitely that the Hamiltonian $`H`$ is self-adjoint and bounded from below, and since $`𝒞`$ is non-negative (we will show this further below) the constant $`_0`$ is equal to the the ground state energy. Moreover, we also see explicitly that our model has a relativistic spectrum: the physical degrees of freedom correspond to one massive- and $`N_\mathrm{F}1`$ massless boson fields. We now briefly describe how to construct a groundstate for our model. We note that with the explicit formulas given above it is straightforward to construct a unitary operator $`𝒰(p)`$ in terms of the operators $`c_j^{()}(p)`$ such that $`C_j(p)=𝒰(p)c_j(p)𝒰(p)^{}`$. Moreover, one can check that $`𝒰=_{p>0}𝒰(p)`$ defines a unitary operator on $`_{\mathrm{Fermion}}^{}`$. It is then easy to see that the state $`\phi _0(Y)𝒰\mathrm{\Omega }_\mathrm{F}_{\mathrm{phys}}^{}`$ with $`\phi _0(Y)\mathrm{exp}(\pi \|M\|\frac{2\pi }{L}Y^2)`$ is annihilated by all operators $`C_j(p)`$ and $`𝒞`$ and thus a ground state of $`H`$. We thus have found one groundstate for our model. However, this state is highly degenerate, and it is actually not invariant under the large gauge transformations Eq. (6. Solution of the model). In the next paragraph we will discuss this issue in more detail. We finally mention that in a complete solution of the model one also needs to find the Green’s functions i.e. vacuum expectation values of gauge invariant combinations of the fermion- and photon field operators. All these Green’s functions can be computed explicitly by using the so-called boson-fermion correspondence which allows to write the fermion field operators in terms of the fermion currents. These computations are similar to the ones for the usual Schwinger model (see e.g. ) but beyond the scope of the present paper. ## 7. Gauge fixing II: Vacuum structure and all that We now describe the structure of the Hilbert space of our model and then perform the second step of the gauge fixing procedure described already above. We first recall the well-known structure of the fermion Fock space $`_{\mathrm{Fermion}}`$: this space can be generated by the fermion current $`\widehat{\rho }_{j,\sigma }(p)`$, $`p0`$, together with unitary operators $`R_{j,\sigma }`$ which obey the relations $$R_{j,\sigma }^1\widehat{\rho }_{j^{},\sigma ^{}}(p)R_{j,\sigma }=\widehat{\rho }_{j^{},\sigma ^{}}(p)+\delta _{p,0}\delta _{j,j^{}}\sigma \delta _{\sigma ,\sigma ^{}}$$ (54) and which interpolate between different sectors labeled by the eigenvalues of the charge operators. Thus for all $`𝐦=(m_{1,+},\mathrm{},m_{N_\mathrm{F},+},m_{1,},\mathrm{},m_{N_\mathrm{F},})Z^{2N_\mathrm{F}}`$, the operators $$𝐑^𝐦=R_{1,+}^{m_{1,+}}\mathrm{}R_{N_\mathrm{F},+}^{m_{N_\mathrm{F},+}}R_{1,}^{m_{1,}}\mathrm{}R_{N_\mathrm{F},}^{m_{N_\mathrm{F},}}$$ (55) commute with all $`h_{p>0}`$ and $`𝒰`$, and if $`\mathrm{\Psi }`$ is a vector in $`_{\mathrm{Fermion}}`$ with $`\widehat{\rho }_{j,\sigma }(0)\mathrm{\Psi }=0`$ $`j,\sigma `$ then $`\widehat{\rho }_{j,\sigma }(0)𝐑^𝐦\mathrm{\Psi }=m_{j,\sigma }𝐑^𝐦\mathrm{\Psi }`$, which implies that $`𝐑^𝐦\mathrm{\Psi }`$ is an eigenvector of $`Q_\mathrm{V}`$ and $`Q_\mathrm{A}`$ with eigenvalues $`𝐧^T𝐦`$ and $`𝐧^T\mathrm{𝐪𝐦}`$, respectively. This implies that if such a vector $`\mathrm{\Psi }`$ is also an eigenstate of all $`h_{p>0}`$ and $`\phi L^2(R)`$ an eigenstate of $`h_0`$ then all states $$\phi (\stackrel{~}{Y})𝐑^𝐦\mathrm{\Psi }=\phi \left(Y+\frac{𝐧\mathrm{𝐪𝐦}}{e\|𝐧\|^2}\right)𝐑^𝐦\mathrm{\Psi }$$ with $`𝐧^T𝐦=0`$ (= charge zero condition) are eigenstates of $`H`$ with eigenvalues which are of the form $`+c(𝐦)`$ and differ only by the contribution from $`𝒞`$ Eq. (51), $`c(𝐦)=\frac{2\pi }{L\|𝐧\|^2}(\|𝐦\|^2\|𝐧\|^2(𝐧^T\mathrm{𝐪𝐦})^2)`$. Note that the latter is always non-negative due to the Cauchy-Schwartz inequality. We thus see that all these states with $`𝐦=k\mathrm{𝐪𝐧}`$ ($`k`$ integer) are degenerate, and especially all states $`\phi _0\left(Y+\frac{k}{e}\right)𝐑^{k\mathrm{𝐪𝐧}}𝒰\mathrm{\Omega }`$ are groundstates for our models. This degeneracy actually is explained by the invariance under large gauge transformation Eq. (6. Solution of the model): This transformation acts on states in $`_{\mathrm{phys}}^{}`$ as $`\phi (Y)\mathrm{\Psi }\phi \left(Y+\frac{1}{e}\right)𝐑^{\mathrm{𝐪𝐧}}\mathrm{\Psi }.`$ We now come to the second step of our gauge fixing procedure. Our discussion above implies that the states $$[\mathrm{\Psi },\phi ]^\theta (Y)=\frac{1}{\sqrt{2\pi }}\underset{kZ}{}\text{e}^{\mathrm{i}k\theta }\phi \left(Y+\frac{k}{e}\right)𝐑^{k\mathrm{𝐪𝐧}}\mathrm{\Psi }$$ (56) ($`\theta `$ real) have simple transformation properties under large gauge transformation Eq. (6. Solution of the model), $`[\mathrm{\Psi },\phi ]^\theta \text{e}^{\mathrm{i}\theta }[\mathrm{\Psi },\phi ]^\theta `$. We thus can define $`_{\mathrm{phys}}`$ as the vector space spanned by all $`[\mathrm{\Psi },\phi ]^0(Y)`$. However, this does not yet make $`_{\mathrm{phys}}`$ into a Hilbert space: a simple computation shows that $$<[\mathrm{\Psi }_1,\phi _1]^\theta ,[\mathrm{\Psi }_2,\phi _2]^\theta ^{}>_{_{\mathrm{phys}}}=\delta (\theta \theta ^{})<[\mathrm{\Psi }_1,\phi _1],[\mathrm{\Psi }_2,\phi _2]>_{\mathrm{phys}}$$ where $`<[\mathrm{\Psi }_1,\phi _1],[\mathrm{\Psi }_2,\phi _2]>_{\mathrm{phys}}=<\mathrm{\Psi }_1,\mathrm{\Psi }_2>_{_{\mathrm{Fermion}}}<\phi _1,\phi _2>_{L^2(R)}`$ is independent of $`\theta `$ and $`\theta ^{}`$. We see that the vectors $`[\mathrm{\Psi },\phi ]^\theta (Y)`$ are actually not contained in $`_{\mathrm{phys}}^{}`$ (they do not have finite norm). The remedy of this problem is a simple multiplicative regularization i.e. ‘dropping the infinite constant $`\delta (0)`$’. This is equivalent to using $`<,>_{\mathrm{phys}}`$ as inner product in $`_{\mathrm{phys}}`$ which is well-defined. This completes the construction of the model. ## 8. Generalization to an arbitrary number of photons We now describe how the above results generalize to a large class of models with $`N_\mathrm{P}`$ different photon fields $`A_\mu ^\alpha `$ where $`\alpha =1,\mathrm{},N_\mathrm{P}`$. These models are (formally) given by , $$=\frac{1}{4}\underset{\alpha =1}{\overset{N_\mathrm{P}}{}}F_{\mu \nu }^\alpha (F^\alpha )^{\mu \nu }+\underset{j=1}{\overset{N_\mathrm{F}}{}}\overline{\psi }_j\gamma ^\mu [\mathrm{i}_\mu +\underset{\alpha =1}{\overset{N_\mathrm{P}}{}}e^\alpha (n_{j,+}^\alpha P_++n_{j,}^\alpha P_{})A_\mu ^\alpha ]\psi _j$$ (57) where $`F_{\mu \nu }^\alpha =_\mu A_\nu ^\alpha _\nu A_\mu ^\alpha `$ and the charge units $`e^\alpha `$ corresponding to the different gauge fields can be different. This model obviously is invariant under transformations belonging to the gauge group $`𝒢=C^{\mathrm{}}(R\times S^1;\mathrm{U}(1)^{N_\mathrm{P}})`$, $`A_\mu ^\alpha A_\mu ^\alpha \frac{1}{e^\alpha }_\mu \chi ^\alpha `$ etc., and as before the existence of large gauge transformations requires all the $`n_{j,\sigma }^\alpha `$ to be integers. The canonical procedure and the construction of the models with $`N_\mathrm{P}=1`$ generalize with minor changes: now we have $`N_\mathrm{P}`$ copies of the photon fields and correspondingly $`N_\mathrm{P}`$ copies of the Gauss’ law operators etc. To be more specific: We now have $`[\widehat{A}_1^\alpha (p),\widehat{E}^\beta (p^{})]=\mathrm{i}\delta ^{\alpha ,\beta }\widehat{\delta }(p+p^{})`$, a Hamiltonian $$H=H_0+\widehat{}_\mathrm{\Lambda }^{}\widehat{\mathrm{d}}p\underset{\alpha =1}{\overset{N_\mathrm{P}}{}}:\left(\frac{1}{4\pi }\widehat{E}^\alpha (p)\widehat{E}^\alpha (p)+e^\alpha \widehat{A}_1^\alpha (p)\widehat{J}^\alpha (p)+\pi \widehat{A}_1^\alpha (p)\underset{\beta =1}{\overset{N_\mathrm{P}}{}}\mathrm{{\rm Y}}^{\alpha \beta }\widehat{A}_1^\beta (p)\right):,$$ (58) and Gauss’ law operators $`\widehat{G}^\alpha (p)=ip\widehat{E}^\alpha (p)+\widehat{\rho }^\alpha (p)`$ where $$\widehat{\rho }^\alpha (p)=\underset{j,\sigma }{}n_{j,\sigma }^\alpha \widehat{\rho }_{j,\sigma }(p),\widehat{J}^\alpha (p)=\underset{j,\sigma }{}\sigma n_{j,\sigma }^\alpha \widehat{\rho }_{j,\sigma }(p)$$ (59) with $`\widehat{\rho }_{j,\sigma }(p)`$ as before. The model has no gauge anomalies if and only if all commutators $`[\widehat{G}^\alpha (p),\widehat{G}^\beta (p^{})]`$ vanish, and similarly as for $`N_\mathrm{P}=1`$ we obtain the conditions in Eq. (2). Note that the mass term we have to add to the naive Hamiltonian depends on a $`N_\mathrm{P}\times N_\mathrm{P}`$-matrix $`\mathrm{{\rm Y}}^{\alpha \beta }`$ which is determined such that $`H`$ commutes with all $`\widehat{G}^\alpha (p)`$. We thus obtain $$\mathrm{{\rm Y}}^{\alpha \beta }=\frac{e^\alpha e^\beta }{2\pi }\underset{j=1}{\overset{N_\mathrm{F}}{}}\underset{\sigma =\pm }{}n_{j,\sigma }^\alpha n_{j,\sigma }^\beta .$$ (60) We note already here that $`𝚼=(\mathrm{{\rm Y}}^{\alpha \beta })_{\alpha ,\beta =1}^{N_\mathrm{P}}`$ is a self-adjoint, real, non-negative $`N_\mathrm{P}\times N_\mathrm{P}`$ matrix, and therefore we can write $$\mathrm{{\rm Y}}^{\alpha \beta }=\underset{\gamma =1}{\overset{N_\mathrm{P}}{}}M_\gamma ^2(a_\gamma )^\alpha (a_\gamma )^\beta $$ (61) where $`(a_\gamma )^\alpha `$ are the components of the orthonormal eigenvectors of $`𝚼`$ and $`M_\gamma ^2`$ the corresponding eigenvalues. For later convenience we assume that $`\mathrm{rank}(𝚼)=N_\mathrm{P}`$ i.e. that all the $`M_\gamma ^2`$ are non-zero.<sup>3</sup><sup>3</sup>3We believe that this assumption could be easily dropped. The generalization of the representation Eq. (18) of the photon fields etc. generalizes in a straightforward manner (we only mention that the generalization of the free Photon Hamiltonian now becomes $`H_{0,\mathrm{P}}=_\alpha \widehat{}_\mathrm{\Lambda }^{}\widehat{\mathrm{d}}p\|M_\alpha \|b_\alpha ^{}(p)b_\alpha (p)`$ where $`M_\alpha ^2`$ are the eigenvalues of the matrix $`𝚼`$). Moreover, also the generalization of Eq. (27) is obvious where the gauge invariant currents now are $`\stackrel{~}{\rho }_{j,\sigma }(p)=\widehat{\rho }_{j,\sigma }(p)+\sigma _{\alpha =1}^{N_\mathrm{P}}e^\alpha n_{j,\sigma }^\alpha \widehat{A}_1^\alpha (p)`$ (note that due to Eq. (2) these latter currents indeed commute with all Gauss law operators). We now come to the solution of the models: the gauge fixing condition generalizing the one in Eq. (29) obviously is $`\widehat{A}_1^\alpha (p)=\delta _{p,0}Y^\alpha `$. Imposing that condition, the Hilbert space of states invariant under all small gauge transformations becomes $`_{\mathrm{phys}}^{}=L^2(R^{N_\mathrm{P}},\mathrm{d}Y^1\mathrm{}\mathrm{d}Y^{N_\mathrm{P}})_{\mathrm{Fermion}}^{}`$ where the zero charge sector in the fermion Fock space now is defined such that $$Q_\mathrm{V}^\alpha =\underset{j=1}{\overset{N_\mathrm{F}}{}}\underset{\sigma =\pm }{}n_{j,\sigma }^\alpha \widehat{\rho }_{j,\sigma }(0)=0\alpha =1,\mathrm{},N_\mathrm{P}$$ (62) on $`_{\mathrm{Fermion}}^{}`$ (again, these latter conditions come from the ($`p=0`$)-components of Gauss’ law). Moreover, the Hamiltonian can again be written as $`H=\frac{2\pi }{L}_{p0}:h_p:`$. For $`p>0`$ we can use the matrix notation introduced in Paragraph 6 and write $`h_p=𝐜^{}\mathrm{𝐡𝐜}`$ where $`𝐡=p\mathrm{𝟏}+\frac{1}{p}𝐘`$ with $$𝐘=\underset{\alpha =1}{\overset{N_\mathrm{P}}{}}\frac{e^\alpha e^\alpha }{2\pi }𝐧^\alpha (𝐧^\alpha )^T$$ (63) and $`(𝐧^\alpha )^T=(n_{1,+}^\alpha ,\mathrm{},n_{N_\mathrm{F},+}^\alpha ,n_{1,}^\alpha ,\mathrm{},n_{N_\mathrm{F},}^\alpha )`$ etc. as before. We now show how to diagonalize these Hamiltonians, following the method explained in Paragraph 6: we compute $`(\mathrm{𝐪𝐡})^2=p^2\mathrm{𝟏}+𝐘+\mathrm{𝐪𝐘𝐪}`$ where we used $`\mathrm{𝐘𝐪𝐘}=0`$ which follows from Eq. (2). We now observe that $`𝐘`$ is a self-adjoint, real, non-negative matrix, and it therefore can be written as $$𝐘=\underset{j=1}{\overset{N_\mathrm{F}}{}}M_j^2𝐞_j𝐞_j^T$$ (64) where the $`𝐞_j`$ are orthonormal and $`M_1^2M_2^2\mathrm{}M_{N_\mathrm{F}}^2`$. The $`M_j^2`$ and $`𝐞_j`$ can computed by diagonalizing the $`N_\mathrm{F}\times N_\mathrm{F}`$ matrix $`𝐘`$. We also observe that the non-zero eigenvalues $`M_j^2`$ of the matrix $`𝐘`$ are identical with the non-zero eigenvalues of the matrix $`𝚼=(\mathrm{{\rm Y}}^{\alpha \beta })_{\alpha ,\beta =1}^{N_\mathrm{P}}`$ defined in Eq. (60). Thus only the $`M_j^2`$ with $`j=1,\mathrm{},\mathrm{rank}(𝚼)=N_\mathrm{P}`$ are non-zero. (To see this, note that if $`(a_j)^\alpha `$ are the components of an eigenvector of $`𝚼`$ with non-zero eigenvalue $`M_j^2`$, then $`𝐟_j=_\alpha e^\alpha (a_j)^\alpha 𝐧^\alpha `$ is an eigenvector of the matrix $`𝐘`$ with the same eigenvalue $`M_j^2`$. One can also check easily that the $`𝐟_j`$ span a space of dimension equal to the rank of the matrix $`𝚼`$ which we assumed to be equal to $`N_\mathrm{P}`$.) Defining $`𝐞_{N_\mathrm{F}+j}=\mathrm{𝐪𝐞}_j`$ we obtain a complete orthonormal basis in $`R^{2N_\mathrm{F}}`$ (orthogonality of the $`𝐞_{jN_\mathrm{F}}`$ and $`𝐞_{j>N_\mathrm{F}}`$ again follows from Eq. (2)), and we can write $$(\mathrm{𝐪𝐡})^2=p^2\mathrm{𝟏}+\underset{j=1}{\overset{N_\mathrm{P}}{}}M_j^2\left(𝐞_j𝐞_j^T+𝐞_{N_\mathrm{F}+j}𝐞_{N_\mathrm{F}+j}^T\right).$$ Thus $`\stackrel{~}{𝐔}`$ Eq. (42) diagonalizes the matrix $`(\mathrm{𝐪𝐡})^2`$, and the eigenvalues of this latter matrix are $`\omega _j^2=p^2+M_j^2`$ where $`M_{N_\mathrm{F}+j}^2=M_j^2`$. It is then easy to check that all the Eqs. (43)–(46) remain true also in the present case, and one finally obtains a representation of $`h_p`$ as in Eq. (47). We now turn to $`h_0`$. After some computations we obtain, $$h_0=\frac{1}{4\pi }\left(\frac{L}{2\pi }\right)^2\underset{\alpha =1}{\overset{N_\mathrm{P}}{}}\frac{^2}{(\stackrel{~}{Y}^\alpha )^2}+\frac{1}{2}\underset{\alpha ,\beta =1}{\overset{N_\mathrm{P}}{}}\mathrm{{\rm Y}}^{\alpha \beta }\stackrel{~}{Y}^\alpha \stackrel{~}{Y}^\beta +𝒞$$ (65) where $`\stackrel{~}{Y}^\alpha =Y^\alpha +_\beta (\mathrm{{\rm Y}}^1)^{\alpha \beta }Q_\mathrm{A}^\beta `$, $`Q_\mathrm{A}=_{j,\sigma }\sigma n_{j,\sigma }^\alpha \widehat{\rho }_{j,\sigma }(0)`$, and $$𝒞=\frac{1}{2}\left(\underset{j=1}{\overset{N_\mathrm{F}}{}}\underset{\sigma =\pm }{}\widehat{\rho }_{j,\sigma }(0)^2\underset{\alpha ,\beta =1}{\overset{N_\mathrm{P}}{}}Q_\mathrm{A}^\alpha (\mathrm{{\rm Y}}^1)^{\alpha \beta }Q_\mathrm{A}^\beta \right).$$ (66) Note that $`\stackrel{~}{Y}^\alpha `$, $`𝒞`$ and $`h_0`$ all are invariant under the large gauge transformations $`Y^\alpha `$ $``$ $`Y^\alpha {\displaystyle \frac{k^\alpha }{e^\alpha }}`$ $`\widehat{\rho }_{j,\sigma }(0)`$ $``$ $`\widehat{\rho }_{j,\sigma }(0)+{\displaystyle \underset{\alpha =1}{\overset{N_\mathrm{P}}{}}}\sigma n_{j,\sigma }^\alpha k^\alpha `$ (67) where $`𝐤=(k^1,\mathrm{},k^{N_\mathrm{P}})Z^{N_\mathrm{P}}`$. Introducing $`Z_j=_\alpha (a_j)^\alpha \stackrel{~}{Y}^\alpha `$ where $`(a_j)^\alpha `$ the components of the eigenvectors of the matrix $`𝚼`$ (cf. Eq. (61)) and $`C_j^{()}(0)=\stackrel{()}{+}\frac{L}{4\pi }(\pi \|M_j\|)^{1/2}\frac{}{Z_j}+(\pi \|M_j\|)^{1/2}Z_j`$ we can write $$h_0=\underset{j=1}{\overset{N_\mathrm{P}}{}}\frac{\|M_j\|}{2}\left(C_j(0)^{}C_j(0)+C_j(0)C_j(0)^{}\right)+𝒞.$$ (68) Again we can combine our results and write $`H`$ as in Eqs. (52)–(53), and we see explicitly that our model has a relativistic spectrum: we have $`N_\mathrm{P}`$ massive and $`(N_\mathrm{F}N_\mathrm{P})`$ massless bosons. It is straightforward to extend the construction of the unitary operator $`𝒰`$ on $`_{\mathrm{Fermion}}^{}`$ diagonalizing all the $`h_p`$ and then check that $`\phi _0(Y^1,\mathrm{},Y^{N_\mathrm{P}})𝒰\mathrm{\Omega }_\mathrm{F}_{\mathrm{phys}}^{}`$ with $`\phi _0=\mathrm{exp}(\frac{2\pi ^2}{L}_{\alpha ,\beta ,\gamma }\|M_\gamma \|(a_\gamma )^\alpha (a_\gamma )^\beta Y^\alpha Y^\beta )`$ (cf. Eq. (61)) is one groundstate of the model. One then can check that for any $`\phi L^2(R^{N_\mathrm{P}})`$ which is an eigenstate of $`h_0`$ and any $`\psi _{\mathrm{Fermion}}^{}`$ which is a common eigenvector of all $`h_p`$, the state $$\phi (\stackrel{~}{Y}^1,\mathrm{},\stackrel{~}{Y}^{N_\mathrm{P}})𝐑^𝐦\mathrm{\Psi }=\phi (Y^1+\frac{k^1}{e^1},\mathrm{},Y^{N_\mathrm{P}}+\frac{k^{N_\mathrm{P}}}{e^{N_\mathrm{P}}})𝐑^{_\alpha k^\alpha \mathrm{𝐪𝐧}^\alpha }\mathrm{\Psi },$$ with $`(𝐧^\alpha )^T𝐦=1`$ for all $`\alpha `$, is an eigenstate of $`H`$ with an eigenvalue of the form $`E+c(𝐦)`$ where $$c(𝐦)=\frac{2\pi }{L}\left(𝐦^2\underset{j=1}{\overset{N_\mathrm{P}}{}}(𝐞_j^T\mathrm{𝐪𝐦})^2\right)$$ with $`𝐞_j`$ the orthonormal eigenvectors of the matrix $`𝐘`$ in Eq. (63). With that one can check that the eigenstate of $`H`$ which also are invariant under the large gauge transformations are $`[\mathrm{\Psi },\phi ]^{(\theta ^1,\mathrm{},\theta ^{N_\mathrm{P}})}(Y)={\displaystyle \frac{1}{(2\pi )^{N_\mathrm{P}/2}}}{\displaystyle \underset{𝐤Z^{N_\mathrm{P}}}{}}\text{e}^{\mathrm{i}k^1\theta ^1}\mathrm{}\text{e}^{\mathrm{i}k^{N_\mathrm{P}}\theta ^{N_\mathrm{P}}}\times `$ $`\times \phi (Y^1+{\displaystyle \frac{k^1}{e^1}},\mathrm{},Y^{N_\mathrm{P}}+{\displaystyle \frac{k^{N_\mathrm{P}}}{e^{N_\mathrm{P}}}})𝐑^{k^1\mathrm{𝐪𝐧}^1}\mathrm{}𝐑^{k^{N_\mathrm{P}}\mathrm{𝐪𝐧}^{N_\mathrm{P}}}\mathrm{\Psi }`$ (69) where now we have $`N_\mathrm{P}`$ real $`\theta `$ parameters. Similarly as for $`N_\mathrm{P}=1`$ the physical Hilbert space $`_{\mathrm{phys}}`$ of the model is spanned by the states $`[\mathrm{\Psi },\phi ]^{(0,\mathrm{},0)}(Y)`$, and one needs to renormalize the inner product of these states, i.e. ‘drop the infinite constant $`\delta (0)^{N_\mathrm{P}}`$’, to get a proper inner product on $`_{\mathrm{phys}}`$. Acknowledgement: We thank M. Luescher for interesting discussions which prompted this work.
warning/0004/nucl-th0004030.html	ar5iv	text	# 1 Introduction ## 1 Introduction The scalar resonances produced in the scattering of pseudoscalar mesons have been a permanent source of discussion.They are advocated as ordinary $`q\overline{q}`$ states , $`q^2\overline{q}^2`$ states , $`K\overline{K}`$ molecules , glueballs and/or hybrids . The properties of these resonances, like their mass, width, partial decay widths and their influence in different reactions are closely tied to their nature, although different interpretations are sometimes possible. The properties of the resonances are modified in the presence of a nuclear medium. They get a selfenergy which changes the position of the peak, the decay width and in addition introduces new decay channels. The predictions for the modification of these properties are also closely related to the hypothesis made on the nature of the resonance. On the other hand, the knowledge of the renormalization of the resonance properties is an essential tool in order to interpret correctly experiments where pairs of pseudoscalar mesons are produced in nuclei close to the energy regions where the resonances appear. This would be the case for the production of two pions which can lead to the $`\rho `$ resonance in L=1 and isospin I=1, or the $`\sigma `$ and the $`f_0(980)`$ in L=0 and I=0, or the production of $`\pi \eta `$ which can lead to the $`a_0(980)`$ resonance in L=0 and I=1. In particular, the renormalization of the $`\rho `$ meson properties in nuclei has witnessed a spectacular effort both theoretical and experimental after the suggestion that there should be a universal scaling of the masses . Although not supported experimentally<sup>1</sup><sup>1</sup>1The pion selfenergy is repulsive instead of attractive, as suggested for the $`\rho `$ in on the basis of the value for the nucleon effective mass in the nucleus, which actually has a different meaning than the one attributed to the $`\rho `$ effective mass in that work., and also challenged theoretically by more recent calculations, the hypothesis certainly had a stimulating effect which has lead to a series of thorough and detailed studies which have set the issue on firmer grounds . Most works would conclude that the $`\rho `$ mass does not appreciably change in the medium, but the width is substantially increased. This latter change alone might be sufficient to explain the spectra of dileptons produced in heavy ion collisions , although further theoretical and experimental work is being done. Another front where there has been much progress is the renormalization of the $`\pi \pi `$ scattering amplitude in a nuclear medium and its possible relationship to the enhancement of the $`\pi \pi `$ invariant mass distribution close to threshold, seen in the experiments of pion induced two pion production in nuclei . It was suggested in that the $`\pi \pi `$ interaction could develop a singularity just below the two pion threshold which would correspond to a kind of Cooper pair. Since this shows up in the scalar and isoscalar channel it was also interpreted as a drop of the sigma mass in the nuclear medium. More refined calculations which have payed special attention to the chiral constraints of the $`\pi \pi `$ amplitude appreciably weaken the renormalization of the amplitude, although there is still an appreciable enhancement of the imaginary part close to threshold . Although there were initial hopes that this alone could explain the experimental features of the $`(\pi ,\pi \pi )`$ reaction in nuclei, more detailed calculations have shown that some additional mechanisms may be required . The former examples show two fronts where the medium properties of mesons are thoroughly investigated. In the present paper we want to pay attention to the modification of the meson scalar resonances $`f_0(980)`$ and $`a_0(980)`$ inside the nuclear medium. The $`f_0(980)`$ resonance has the same quantum numbers as the $`\sigma `$, only its mass is much larger, and conversely the width is much smaller, of the order of 40-100 MeV (versus $`m_\sigma =450MeV`$ and a width of around 450 MeV according to the PDG book ). The $`a_0(980)`$ in the I=1 channel has about the same mass as the $`f_0(980)`$ and a width of the order of 50-100 MeV. In spite of its relevance as a source of information on the nature of these resonances and the interest that it should have in the analysis of the two meson production in nuclei, the surprising fact is that there are neither theoretical nor experimental studies on this issue, which contrast with the large efforts devoted to the study of the $`\rho `$ meson properties in an energy region very close to the one where these resonances appear. The reasons for this might be simply technical. On the one hand the confusion about the nature of these states was a deterrent. On the other hand, the study of the properties of these resonances in nuclei, which couple both of them largely to the $`K\overline{K}`$ system, had their own problems since the question of the $`K^{}`$ selfenergy in a nuclear medium was itself unclear . Fortunately things have changed in both fronts recently, to the point that one can count on reasonable models with which to tackle the problem. One of the areas which has witnessed an important progress in recent years is the meson meson interaction by means of Chiral Perturbation Theory ($`\chi PT`$), which is supposed to be the effective theory of QCD at low energies . The theory has proved rich in applications to strong, weak and electromagnetic processes in which pairs of mesons appear at small energies . Yet, implicit to the perturbative nature of the theory is the fact that it does not generate poles in the scattering amplitudes and hence is unsuited to study the energy regions where meson resonances appear. In this respect, there have been recent advances which have shown the usefulness of $`\chi PT`$ as a means to constrain non perturbative unitary methods. For instance, by means of the inverse amplitude method (IAM) and the chiral Lagrangians, one can obtain the $`\sigma `$ and $`\rho `$ mesons in the $`\pi \pi `$ scattering and the $`K^{}`$ resonance in $`K\pi `$ scattering. The generation of the $`f_0(980)`$ and $`a_0(980)`$ resonances, however, required the extension of the IAM to coupled channels, which was done in where these two mesons were also generated. In addition, all meson meson properties up to about 1.2 GeV were reproduced using simply the standard $`O(p^2)`$ and $`O(p^4)`$ chiral Lagrangians as input. Advances have also been made using the hypothesis of resonance saturation of Ref. , which states that the information of the $`O(p^4)`$ chiral Lagrangian is tied to the exchange of resonances which survive in the large $`N_c`$ limit. By allowing such genuine (preexisting to the unitarization, or multiple scattering of the mesons) resonances and using the $`O(p^2)`$ chiral Lagrangian in addition, together with a proper unitarization scheme based on the N/D method, it was shown in that a good description of all the meson meson information up about 1.5 GeV could be accomplished. In this respect it is interesting to observe that the use of the lowest order $`\chi PT`$ Lagrangian, properly unitarized by means of the Bethe Salpeter equation, together with an appropriate cut off to regularize the loops, is able to reproduce all the information of the meson meson scattering in the scalar sector up to about 1.2 GeV and gives the same results as the more general methods reported above. This is possible due to the large weight of the lowest order Lagrangian in the scalar sector, in contrast to what happens in the vector sector where higher orders play an essential role. The meson baryon interaction, from the point of view of $`\chi PT`$, has also witnessed much progress . On the other hand the meson baryon interaction has been the ground for application of the chiral unitary techniques, using the Lippmann Schwinger summation or the Bethe Salpeter integral equation and more recently the IAM in or the N/D method . The work of initiated a fruitful line which has allowed to link the $`K^{}N`$ interaction with its coupled channels at low energies, plus the properties of the $`\mathrm{\Lambda }(1405)`$ resonance, with the information of the chiral Lagrangians. The work of added more channels in the coupled channels set and proved that it was possible to get a good reproduction of the low energy data in terms of the lowest order chiral Lagrangian and a suitable cut off for the loops. This progress in the elementary reaction also stimulated work directed to obtain the selfenergy of the kaons in nuclei. It was soon realized that the consideration of Pauli blocking in the intermediate nucleon states lead to a shift towards higher energies of the $`\mathrm{\Lambda }(1405)`$ resonance and, as a consequence, to an appreciable attraction on the kaon . However, the large kaon attraction which was found suggested that the problem had to be solved selfconsistently. This was done in , where it was found that the selfconsistent treatment rendered the resonance to the same free position, although it appeared with a substantially larger width. This problem has been further pursued in where all elements of were considered but the renormalization of the pions in the intermediate states was also taken into account. With all these ingredients a $`K^{}`$ nucleus optical potential was found which has proved consistent with the information of kaonic atoms . Thus, the situation at present is that one has the information and the means to generate the scalar mesons from the chiral Lagrangians and also the interaction of the kaons with the nucleus, the two ingredients that blocked progress in the topic of the scalar meson renormalization in nuclei, and this is the problem that we shall tackle in this paper. ## 2 The meson meson interaction in the nuclear medium Here, we briefly review the formalism for the meson meson interaction which was developed in . We start by taking the states $`\pi \pi `$ and $`K\overline{K}`$ ,which we label 1 and 2. In the I=0 channel they are given by $$\begin{array}{c}\|K\overline{K}>=\frac{1}{\sqrt{2}}\|K^+(\stackrel{}{q})K^{}(\stackrel{}{q})+K^0(\stackrel{}{q})\overline{K}^0(\stackrel{}{q})>\hfill \\ \|\pi \pi >=\frac{1}{\sqrt{6}}\|\pi ^+(\stackrel{}{q})\pi ^{}(\stackrel{}{q})+\pi ^{}(\stackrel{}{q})\pi ^+(\stackrel{}{q})+\pi ^0(\stackrel{}{q})\pi ^0(\stackrel{}{q})>,\hfill \end{array}$$ (1) and the $`K\overline{K}`$, $`\pi \eta `$ which we label 1 and 2, in the isospin I=1 channel, $$\begin{array}{c}\|K\overline{K}>=\frac{1}{\sqrt{2}}\|K^+(\stackrel{}{q})K^{}(\stackrel{}{q})K^0(\stackrel{}{q})\overline{K}^0(\stackrel{}{q})>\hfill \\ \|\pi \eta >=\|\pi ^0(\stackrel{}{q})\eta (\stackrel{}{q})>,\hfill \end{array}$$ (2) where $`\stackrel{}{q}`$ is the momentum of the particles in the CM of the pair. We follow the convention $`\|\pi ^+>=\|1,1>`$ and $`\|K^{}>=\|\frac{1}{2},\frac{1}{2}>`$ isospin states. In the I=0 case we neglect the $`\eta \eta `$ channel. Its contribution has been assessed in and it has only relevance at energies beyond 1.2 GeV, hence, as done in , we shall omit this channel here too. The Bethe Salpeter (BS) equation is given by $$T=VGT.$$ (3) Eq. 3 is meant as a coupled channel equation with two channels in each of the isospin states. It is an integral equation, meaning that the term $`VGT`$ involves one loop integral where both $`V`$ and $`T`$ would appear off shell. However, it was shown in that, for the free case, these amplitudes could be factorized on shell out of the integral and the off shell part was absorbed by a renormalization of the coupling constants. Thus, the equation becomes purely algebraic and is quite easy to solve. In Eq. 3, the formal product of $`VGT`$ inside the loop integral becomes then the product of $`V`$,$`G`$ and $`T`$, with $`V`$ and $`T`$ the on shell amplitudes, and the function $`G`$ is given by the diagonal matrix $$G_{ii}=i\frac{d^4q}{(2\pi )^4}\frac{1}{q^2m_{1i}^2+iϵ}\frac{1}{(Pq)^2m_{2i}^2+iϵ}$$ (4) where $`P`$ is the total fourmomentum of the meson-meson system. The BS equation sums the series of diagrams depicted in Fig. 1. In the nuclear medium, one has to add the diagrams depicted in Fig. 2, which stem from the interaction of the pions with the medium, through $`ph`$ and $`\mathrm{\Delta }h`$ excitation. We neglect here the small s-wave pion-nucleon interaction. The interesting finding of and was that the contact terms involving the $`ph`$ ($`\mathrm{\Delta }h`$) excitations of diagrams (b)(c)(d) canceled exactly the off shell contribution from the meson meson vertices in the term of Fig. 2(a). Hence technically one only had to evaluate the diagrams of the free type (Fig. 1) and those of Fig. 2(a) (plus higher order iterations ) in order to evaluate the pion pion scattering amplitude in the nuclear medium. In , only the pions were renormalized inside the medium since the work was only concerned with the meson-meson interaction at low energies where the kaons do not play much of a role. However, in order to address now the question of the $`f_0(980)`$ and $`a_0(980)`$ resonances we shall have to take into account the kaon selfenergy in the medium. Since we are close to the two kaon threshold, the s-wave part of the kaon selfenergy is the relevant ingredient. For completeness, we shall also consider the p-wave selfenergy due to $`\mathrm{\Sigma }h`$ or $`\mathrm{\Lambda }h`$ excitation, as in , although it does not play an important role in the present problem. The loop involving pions is done as in . This provides the $`\pi \pi `$ matrix element of the $`G`$ function in the medium. As for the $`K\overline{K}`$ intermediate states we have an asymmetric situation. In particular, the selfenergy of the $`\overline{K}`$ requires special care. We will use the results of the detailed model of Ref. . On the other hand, the $`K`$ selfenergy can be accounted for in a much simpler way since there are no resonances with strangeness $`S=1`$ and the $`KN`$ interaction is quite smooth. Altogether, $`t\rho `$ gives a very reasonable approximation to the $`K`$ selfenergy. By taking results from or we can write $$\mathrm{\Pi }=\frac{1}{2}(t_{Kp}+t_{Kn})\rho 0.13m_K^2\frac{\rho }{\rho _0}$$ (5) where $`t`$ is the elastic $`K`$-nucleon amplitude, $`\rho `$ is the nuclear density and $`\rho _0`$ is the normal nuclear density. The new meson-meson amplitude in the medium is now given by means of the modified BS equation $$\stackrel{~}{T}=V\stackrel{~}{G}\stackrel{~}{T},$$ (6) where the $`K\overline{K}`$ matrix element of $`\stackrel{~}{G}`$ is given by $$\stackrel{~}{G}_{K\overline{K}}=i\frac{d^4q}{(2\pi )^4}\frac{1}{q^2m_K^2\mathrm{\Pi }_K(q^0,q,\rho )}\frac{1}{(Pq)^2m_K^2\mathrm{\Pi }_{\overline{K}}(P^0q^0,q,\rho )}$$ (7) The $`\overline{K}`$ propagator can be written in the Lehmann representation, $$\frac{1}{(Pq)^2m_K^2\mathrm{\Pi }_{\overline{K}}(P^0q^0,q,\rho )}=_0^{\mathrm{}}𝑑\omega 2\omega \frac{S_{\overline{K}}(\omega ,q,\rho )}{(P^0q^0)^2\omega ^2+iϵ}$$ (8) where $`S_{\overline{K}}(\omega ,q,\rho )`$ is the spectral function of the $`\overline{K}`$ in the medium, as described in . The $`q^0`$ and angular integrations can be done analytically and then we find $$\stackrel{~}{G}_{K\overline{K}}=\frac{1}{2\pi ^2}_0^{\mathrm{}}𝑑\omega _0^{q_{max}}𝑑qq^2S_{\overline{K}}(\omega ,q,\rho )\frac{\omega +\stackrel{~}{\omega }(q)}{\stackrel{~}{\omega }(q)(s(\omega +\stackrel{~}{\omega }(q))^2)+iϵ}$$ (9) where $`\stackrel{~}{\omega }(q)`$ is given by $$\stackrel{~}{\omega }(q)=\sqrt{\stackrel{}{q}^2+m_K^2+\mathrm{\Pi }_K}.$$ (10) For the $`\pi \eta `$ intermediate channel we proceed in the same way. The $`\pi `$ propagator is also written in terms of the Lehmann representation, as done in and the $`\eta `$ propagator explicitly in terms of the eta selfenergy, as we have done above for the $`K`$. We also take the $`t\rho `$ approximation for the $`\eta `$ selfenergy. However, the $`t`$ matrix is not so well known in this case. We take the following amplitudes from the fit of $$t^1=\frac{M}{4\pi \sqrt{s}}(\frac{1}{a}+r_0q_\eta ^2+s_0q_\eta ^4iq_\eta ),$$ (11) where $`M`$ is the nucleon mass, $`\sqrt{s}`$ is the center of mass energy of the $`\eta `$-nucleon system, and $`q_\eta `$ is the momentum of the $`\eta `$ meson in the same system, assuming $`\eta `$ and nucleon to be on shell, $`a=0.75+0.27ifm`$, $`r_0=1.500.24ifm`$ and $`s_0=0.100.01ifm^3`$. We use these results from $`\eta N`$ threshold up to a value of $`\sqrt{s}`$ 200 MeV above it. Outside the region of validity of this fit, we have assumed a constant selfenergy equal to the one of the closest extreme of the parametrized region. The uncertainties associated to this assumption are small since taking just a cero $`\eta `$ selfenergy outside that range of energies changes the results by less than 5%. ## 3 Other medium corrections ### 3.1 Tadpole terms In this section we consider some additional many body corrections. We begin by the contribution of the diagram depicted in Fig. 3. which has been considered in using the linear sigma model, with sigmas and pions as elementary fields. The Lagrangian gives rise to a vertex with three sigmas as shown in the figure which produces a many body correction, assuming a certain coupling of the sigma to the nucleons, which in is borrowed from the Bonn phenomenological boson exchange models of the $`NN`$ interaction . On the other hand, the $`\chi PT`$ Lagrangian involves only pseudoscalar meson fields. The sigma can be generated dynamically through the rescattering of the pions. The closest analog to the diagram of Fig. 3 is given by the series of Fig. 4 where the sigma is generated to the left and right of the nucleon-hole loop by means of the iteration of the chiral Lagrangian. In order to evaluate the contribution of these diagrams we use the chiral Lagrangians involving the octet of baryons and the octet of pseudoscalar mesons. The terms needed come from the covariant derivative terms of the Lagrangian. After some trivial algebra, we find that for the scalar channel these terms are proportional to $`\overline{p}\gamma ^\mu p\overline{n}\gamma ^\mu n`$ and therefore would vanish in symmetric nuclear matter after summation over protons and neutrons. Considering higher order corrections, there is another possible way to generate an analogous structure as it is shown in Fig. 5 Here one also needs the coupling of the sigma, generated through pion-pion interaction, with the nucleons. The scalar isoscalar exchange in the $`NN`$ interaction has been addressed in only at the perturbative level. In it has been revisited taking into account all the meson meson rescattering which generates the sigma. In this latter work one finds an attraction (not generated in the perturbative approach) at intermediate distances, but weaker than in the Bonn model because an appreciable repulsion sets up already at distances like 0.7 fm, which grows fast at small distances. However, the relevant magnitude would be the strength of the potential in momentum space at zero momentum, which is what is met in Fig. 5 and there the strength of the potential of is about one order of magnitude smaller than for the Bonn potential. Hence, from this source we also get a negligible contribution to the modification of the $`\pi \pi `$ scattering in the nuclear medium. ### 3.2 Roper-hole excitation Finally, we will consider the excitation of resonances from the occupied nucleon states by the pair of mesons. This requires, in our case, a resonance which decays into a nucleon and two pions (two mesons in general) in s-wave and with I=0. There is such a candidate at the energies where we are concerned which is the $`N^{}(1440)`$ Roper resonance. The branching ratio for this decay is small but it plays a very important role in the $`\pi N\pi \pi N`$ reaction . It has also been shown that this mechanism is the dominant in the $`NNNN\pi \pi `$ reaction close to threshold . In the present problem it could contribute via the mechanism depicted in Fig. 6. The evaluation of this mechanism is straightforward in the $`SU(2)`$ sector. In the first place we take the Lagrangian suggested in which reads $$_{N^{}N\pi \pi }=c_1^{}\overline{\psi }_N^{}\chi _+\psi _N\frac{c_2^{}}{M^{}^2}(_\mu _\nu \overline{\psi }_N^{})u^\mu u^\nu \psi _N+h.c.$$ (12) where $$\chi _+=m_\pi ^2(2\frac{\stackrel{}{\varphi }^2}{f^2}+\mathrm{});u_\mu =\frac{1}{f}\stackrel{}{\tau }_\mu \stackrel{}{\varphi }+\mathrm{}$$ (13) and $`\stackrel{}{\varphi }`$ is the pion field. The value of the $`c_1^{}`$ and $`c_2^{}`$ constants is not fully determined by the partial decay width of the $`N^{}`$, which only tells us that these values satisfy the equation of an ellipse. One needs to find extra constraints and in it was found that the best set that lead to agreement with the $`\pi ^{}p\pi ^+\pi ^{}p`$ reaction was $`c_1^{}=7.27GeV^1`$ and $`c_2^{}`$=0. This was corroborated in by looking at all the different isospin channels. With the $`c_1^{}`$ term alone we obtain the Lagrangian, after we expand up to two pion fields $$_{N^{}N\pi \pi }=\frac{1}{f^2}c_1^{}m_\pi ^2(\pi ^0\pi ^0+2\pi ^+\pi ^{})(\overline{p}^{}p+\overline{n}^{}n)+h.c.$$ (14) There is no experimental information on the $`N^{}NK\overline{K}`$ coupling. For simplicity we will assume the simple generalization of the previous Lagrangian $$=\frac{1}{2}c_1^{}Tr(\overline{B}B)Tr(\chi _+),$$ (15) The baryonic matrix $`B`$ and the mesonic matrix $`\chi _+`$ can be found in Ref. . Expanding this Lagrangian up to terms with two meson fields we obtain and keeping only those terms relevant to our calculation we get $$=\frac{c_1^{}}{f^2}\{m_\pi ^2(\pi ^0\pi ^0+2\pi ^+\pi ^{})+m_K^2(K^0\overline{K}^0+K^+K^{})\}(\overline{p}^{}p+\overline{n}^{}n)+h.c.$$ (16) We can see that the strength of the coupling to $`K\overline{K}`$ is $`\frac{m_K^2}{m_\pi ^2}`$ times that of the charged pions. Other simple generalizations of the $`SU(2)`$ Lagrangian, i.e. $`Tr(\overline{B}\chi _+B)`$, also provide that scaling, although with some different numeric factors and charge combinations. With the choice of Eq. 16, the contribution of the mechanism of Fig.6 to the meson meson scalar isoscalar interaction is then given by $$V_{ij(N^{}N)}^{}=(c_1^{})^2\frac{4m_i^2m_j^2}{f^4}𝒰(p^0,\stackrel{}{p}=0,\rho )$$ (17) where $`i,j`$ stand for the physical meson states, $`m_i,m_j`$ are the masses of the initial and final mesons and $`𝒰(p^0,\stackrel{}{p}=0,\rho )`$ is the complex Lindhard function for the excitation of resonances , taken from , which in the limit of $`\stackrel{}{p}`$=0, which one has for a meson pair in their center of mass frame, has the simple expression $$𝒰(p^0,\stackrel{}{p}=0,\rho )=\frac{\rho }{p^0m_N^{}+m_N+i\frac{\mathrm{\Gamma }_N^{}}{2}}.$$ (18) This leads to an additional contribution to the $`\pi \pi `$ ”potential” used in the Bethe Salpeter equation which is given by the matrix elements for I=0 (for I=1 one gets zero contribution) $`V_{11}^{}`$ $`=`$ $`2(c_1^{})^2{\displaystyle \frac{4m_K^4}{f^4}}𝒰(p^0,\stackrel{}{p}=0,\rho )`$ (19) $`V_{12}^{}`$ $`=`$ $`\sqrt{3}(c_1^{})^2{\displaystyle \frac{4m_\pi ^2m_K^4}{f^4}}𝒰(p^0,\stackrel{}{p}=0,\rho )`$ (20) $`V_{22}^{}`$ $`=`$ $`{\displaystyle \frac{3}{2}}(c_1^{})^2{\displaystyle \frac{4m_\pi ^4}{f^4}}𝒰(p^0,\stackrel{}{p}=0,\rho )`$ (21) Now one must be cautious to use the empirical value of the $`c_1^{}`$ parameter when the potential of the former equations is put into the kernel of the BS equation. Indeed, at the level of terms linear in the density one will generate the diagrams of the figure This means that we are generating the iteration of the free mesons to the right and the left of the $`N^{}h`$ excitation. As a consequence the empirical coupling should be related to a bare one by $$c_1^{}=c_{1B}^{}(1+G_{\pi \pi }T_{\pi \pi ,\pi \pi }^{I=0}+\sqrt{\frac{4}{3}}\frac{m_K^2}{m_\pi ^2}G_{KK}T_{K\overline{K},\pi \pi }^{I=0})$$ (22) Thus the calculations must be done using the bare coupling, since the free pion interaction renormalizes it to the effective values demanded by the $`N^{}`$ decay into two pions in s-wave. ## 4 Results and discussion In the first place let us discuss the results in the I=1 channel. In figs. 8,9 we show the real and imaginary parts of the $`K\overline{K}K\overline{K}`$ and $`\pi \eta \pi \eta `$ amplitudes for different values of the nuclear density at energies around the $`a_0(980)`$ meson. An inspection to the imaginary part of Fig. 8 seems to indicate that the peak of its magnitude, corresponding to the apparent mass of the $`a_0(980)`$ resonance, moves to low energies as the density increases, producing a shift of about $`50`$ MeV at $`\rho =\rho _0`$. This shift is also visible in the real part by looking at the point where the real part changes sign. The apparent width measured from the imaginary part of the amplitude becomes bigger as the density increases and for $`\rho =\rho _0`$ becomes as large as 200 MeV from an apparent free width of around 90 MeV. This behaviour is, however, not reproduced in the $`\pi \eta \pi \eta `$ amplitude. The first difference one may observe is the presence of a larger background, even in free space, which makes it not resemble a Breit Wigner resonance so much. For this channel, the $`a_0`$ resonance does not move much with the density and the width becomes very large already at small densities, to the point that at densities of the order of one half $`\rho _0`$ the resonant shape is practically lost. Given the fact that the $`a_0(980)`$ meson is usually observed in mass distributions of $`\pi \eta `$ in the final state of some reaction, the relevant magnitude entering the cross section of these reactions close to the resonance is the modulus squared of the $`\pi \eta \pi \eta `$ amplitude plotted in Fig. 10. What we observe there is, indeed, that the resonance melts very fast as the density increases and at densities of the order of $`\rho _0/2`$ there is practically no resonant trace left. The fast disappearance of this relatively narrow resonance in nuclei is probably one of the most striking predictions for this channel. The I=0 channel has a richer structure as a function of the density. In Fig. 11 we show the amplitude $`\pi \pi \pi \pi `$ in a range of energies from 200 MeV to 1100 MeV. The figure shows results at low energies already discussed in . As one can see in the figure, there is an accumulation of strength in the imaginary part below threshold which was first pointed out in and has also been predicted in other approaches . The relationship of this increased strength to the enhanced two pion distribution in $`(\pi ,2\pi )`$ reactions in nuclei at small invariant masses has been discussed in , but according to the detailed calculation of it is not enough to reproduce the experimental data. The intermediate region of energies is quite interesting and no much attention has been given to it so far. There one can see a drastic decrease of the strength of the imaginary part as the density increases. This reduction could lead to appreciable changes in the two pion production reactions in nuclei, like the $`(\gamma ,2\pi )`$ reaction for which experiments are already becoming available . The region of the $`f_0(980)`$ resonance is also interesting. By looking both at the dip of the imaginary part of the amplitude, as well as to the position of the zero of the real part, we can see that the position of the resonance does not change when the density increases. We observe, however, a gradual melting of the dip of the imaginary part which comes as an interference between the background of the $`\sigma `$ meson contribution and the contribution of the $`f_0(980)`$ resonance. It is interesting to discuss what happens when we introduce the $`N^{}h`$ excitation, which is a novel ingredient with respect to the approach of . In Fig. 12 we show the results in which the SU(3) version of the coupling of the $`N^{}`$ to $`N\pi \pi `$ and $`NK\overline{K}`$ is used. We can see that the inclusion of this new ingredient barely modifies the results of the amplitude in all the range of energies shown. These changes amount to about a 10 per cent increase of the strength of the imaginary part of the amplitude in the region around 300 to 400 MeV. We do not show it here, but point out that using the SU(2) version of the $`N^{}`$ coupling to $`N\pi \pi `$, in which the $`N^{}`$ only couples to $`N\pi \pi `$, practically does not change the results with respect to those in which the $`N^{}`$ is allowed to couple to pions and kaons. The role of the $`N^{}h`$ excitation is more apparent in the $`K\overline{K}\pi \pi `$ and $`K\overline{K}K\overline{K}`$ amplitudes. In Fig. 13 we show the $`K\overline{K}\pi \pi `$ amplitude for different densities. This amplitude is better suited than the $`\pi \pi \pi \pi `$ in order to show the $`f_0(980)`$ resonance because it does not have a large background. Notice that in this inelastic channel the $`f_0(980)`$ shows up as a Breit Wigner contribution rotated 90 degrees. Thus, the roles of the real and imaginary parts of the amplitude are interchanged . We can see that as the density increases the position of the resonance barely moves. The width, however, grows with the density from a free value of around 30 MeV to about 100 MeV at $`\rho =\rho _0`$. If we include the $`N^{}h`$ contribution in the SU(2) formulation there are no appreciable changes with respect to those shown in the figure. The results are however quite different if we include the $`N^{}h`$ contribution in the SU(3) formulation. These results are shown in Fig. 14. We can see there that around 300 to 400 MeV a resonant like structure develops with the imaginary part showing a negative peak and the real part changing fast around a zero value. This is a reflection of the $`N^{}h`$ excitation which in this case is magnified because of the large coupling of the $`N^{}h`$ excitation to $`K\overline{K}`$, as we saw in section 3. As we discussed there, we found a large coupling, of the order of $`m_K^2/m_\pi ^2`$ that of the pion, based on a generalization to SU(3) of the pion coupling. We also saw that there were ambiguities, but any simple generalization led to a coupling of this order of magnitude. In spite of this huge coupling, we saw no visible effects in the $`\pi \pi \pi \pi `$ amplitude. Here it shows clearly in a large medium change of the $`\pi \pi K\overline{K}`$ amplitude. The effects in the $`K\overline{K}K\overline{K}`$ amplitude are even more pronounced. In Fig. 15 we show the $`K\overline{K}K\overline{K}`$ in the energy region around the $`f_0`$ resonance. The density effects around the pole can be appreciated better and one can see that even at $`\rho =\rho _0`$ the shape of the resonance is not lost, but the width increases to about 100 MeV at $`\rho _0`$. As we have said, there are some elements of uncertainty tied to the extrapolation of the $`N^{}`$ coupling to kaons, and it would be worth trying to find some observable consequences of the assumptions made. Since the larger effects are seen in a region where the kaons are far off shell, one can only hope to observe indirect effects on reactions were the kaons appear as intermediate states. The large values obtained in the amplitudes will certainly be softened by the small weight of a $`K\overline{K}`$ propagator where the two kaons are quite off shell, so one should not expect drastic changes. Yet, even moderate changes might be relevant in some processes like the $`(\pi ,2\pi )`$ reaction in nuclei, where it was shown in that there were large cancellations between terms to give a final result smaller than the contribution of individual terms, such that any small changes in one of them might alter the final balance. The finding of indirect evidence of this directly unobservable $`N^{}`$ coupling to N and kaons would be an important test of particle symmetries. In any case the interesting medium effects found here, independent of the still unknown couplings, would certainly call for devoted experiments from which we could learn more about the nature of the scalar resonances and the way the meson meson interaction is changed in a medium. Reactions like $`\gamma p\pi \pi p`$ have already been suggested as a means to observe the scalar resonances . Their extension using nuclear targets is certainly feasible and, together with other experiments, should be encouraged. ## 5 Conclusions In section 3 we addressed the question of new contributions to the $`\pi \pi `$ scattering in a nuclear medium beyond those already considered in other approaches. One of the terms considered in which a nucleon loop is attached to the four meson vertex was found to be zero for symmetric nuclear matter. Other possible mechanisms which would simulate a three sigma vertex coupling were also estimated to be much weaker than previously suggested. These results would further strengthen those obtained in , which would in turn mean that the experimental problem of the enhanced invariant $`\pi \pi `$ mass close to threshold would not be solved yet. The main topic of the present paper has been the discussion of the renormalization of the properties of the scalar meson resonances, concretely the $`f_0(980)`$ and the $`a_0(980)`$ resonances, in the nuclear medium. The renormalization required the use of the kaon selfenergy for which we have used a recent one deduced from chiral Lagrangians and which is consistent with the information of kaonic atoms. We have systematically tried to use the chiral unitary formalism in the different aspects of the problem, be the generation of the resonances through the meson meson interaction given by the chiral meson Lagrangians, or the meson baryon interaction, which for the most delicate case, the one of the $`K^{}`$, is also obtained by means of a nonperturbative chiral approach. The results obtained are interesting, we do not observe an appreciable change of the position of either resonance. However, the widths are substantially changed. In the case of the $`f_0(980)`$ resonance the width passes from 30 MeV in the free case to about 100 MeV at normal nuclear matter. In the case of the $`a_0(980)`$ resonance the width grows so fast with density that even at $`\rho _0/2`$ there is practically no trace of the resonance. The next step should be the search for these effects in nuclear experiments which can help shed new light on the nature of these resonances and the behaviour of kaons in nuclear matter. Acknowledgments: We would like to thank A. Ramos for providing us with her codes to calculate kaonic spectral functions. Useful discussions with N. Grion, T. Hatsuda, G. Chanfray, R. Rapp, and J. Wambach are also acknowledged. This work is partly supported by DGICYT contract no. PB 96-0753 and by the EEC-TMR Program, EURODAPHNE, Contract No. ERBFMRX-CT98-0169.
warning/0004/cond-mat0004445.html	ar5iv	text	# Mode Coupling Approach for spin–facilitated kinetic Ising models ## I Introduction In spite of advances in the description of liquids near the glass transition using different approaches the phenomenon is generally not completely understood. Supercooled fluids reveal often a non-Arrhenius behavior of relaxation times and a characteristic stretched exponential decay of correlation functions. But a long range order is not developed in contrast to conventional phase transitions. The dynamics of the glass transition can be described by an increasing cooperativity of local processes with decreasing temperature. The cooperativity leads to the well known slowing down of the relaxation times which can be illustrated by a strongly curved trajectory in the Arrhenius plot (logarithm of the relaxation time $`\tau `$ versus the inverse temperature $`T^1`$). One possible fit of this curve is given by the Williams-Landel-Ferry (WLF) relation, i.e. $`\mathrm{ln}\tau \left(TT_0\right)^1`$ with a finite Vogel temperature $`T_0`$. Mode coupling theories (MCT) predict a completely ergodic decay of correlation functions above a critical temperature $`T_c`$ and an incomplete decay below $`T_c`$. The incomplete decay is usually identified as the relatively fast $`\beta `$–process caused by processes related to the breaking up of the local cages. The long time regime of the ergodic decay above $`T_c`$ is denoted as $`\alpha `$–process. The relaxation time of the $`\alpha `$–process increases rapidly with decreasing temperature and below $`T_c`$ only the $`\beta `$–process remains effective, i.e. at $`T_c`$ the system undergoes a sharp phase transition to a state with partially frozen (density) fluctuations. Note that $`T_c`$ is in the range between the melting temperature $`T_m`$ and the glass temperature $`T_g`$, e.g. $`T_m>T_c>T_g`$. It is really the $`\alpha `$–process exists also below $`T_c`$. This process leads to a very slow decay of the apparently frozen structures, i.e. the nonergodic structures obtained from the MCT are approximately stable only for a finite time interval. This slow decay shows the typical properties which correspond usually to the dynamics of the main glass transition (WLF like behavior of the relaxation time, stretched exponential decay of the correlation function). This effects can be partially described in terms of an extended mode coupling theory introducing additional hopping processes. There exists also various alternative descriptions which explain the cooperative motion of the particles inside a supercooled liquid below $`T_c`$. One of these possibilities is the $`n`$–spin facilitated kinetic Ising model , originally introduced by Fredrickson and Andersen. The base idea of this model consists in a coarse graining of space and time scales and simultaneously a reduction of the degrees of freedom. In detail that means: 1. Coarse graining of spatial scales: The supercooled liquid is separated into cells, so that each cell contains a sufficiently large number of particles which realize a representative number of molecular motions, i.e. the many body system is considered of a virtual lattice with the unit size $`l`$. This lattice has no influence on the underlying dynamics of the supercooled liquid. 2. Reduction of the degrees of freedom: Each cell will be characterized by only one trivial degree of freedom, i.e. the cell structure enables us to attach to each cell an observable $`\sigma _j`$ which characterizes the actual state of particles inside the cell $`j`$. The usual realization is given by the local density $`\rho _j`$ (particles per cell) with $`\sigma _j=1`$ if $`\rho _j>\overline{\rho }`$ and $`\sigma _j=1`$ if $`\rho _j<\overline{\rho }`$ where $`\overline{\rho }`$ is the averaged density of the system. This mapping implies consequently different mobilities of the particles inside such a cell, i.e. $`\sigma _j=1`$ corresponds to a more immobile (or solid like) state and $`\sigma _j=1`$ to a more mobile (or liquid like) state of the cell $`j`$. The set of all observables $`\sigma =\left\{\sigma _j\right\}`$ forms a configuration. The evolution of the statistical probability distribution function $`P(\sigma ,t)`$ can be described by a generalized master equation using a projection of the real dynamics onto the dynamics of $`\sigma `$: $`{\displaystyle \frac{P(\sigma ,t)}{t}}={\displaystyle \underset{\sigma ^{}}{}}L^{}(\sigma ,\sigma ^{})P(\sigma ,t)+{\displaystyle \underset{\sigma ^{}}{}}{\displaystyle \underset{0}{\overset{t}{}}}K(\sigma ,\sigma ^{},tt^{})P(\sigma ,t^{})𝑑t^{}`$ 3. Coarse graining of the time scale: The last step bases on the assumption that possible memory terms $`K(\sigma ,\sigma ^{},t)`$ of the generalized master equation is mainly determined by fast molecular processes while the slow dynamics is mainly reflected by the temporally local contributions $`\underset{\sigma ^{}}{}L^{}(\sigma ,\sigma ^{})P(\sigma ,t)`$. Of course, the validity of this assumption depends strongly on the choice of the remaining degrees of freedom, and in many cases it is very hard (or impossible from the actual point of view) to give a satisfactory explantation of this assumption. However, if this separation of the dynamics is justified, an elementary time scale larger than the time scale of the fast molecular processes can be introduced. Therefore, the memory will be reduced to simple temporally local terms, i.e. $`K(\sigma ,\sigma ^{},tt^{})=\delta \left(tt^{}\right)_0^{\mathrm{}}K(\sigma ,\sigma ^{},\tau )𝑑\tau `$. One obtains an evolution equation which is equivalent to the mathematical representation of a usual master equation. $$\frac{P(\sigma ,t)}{t}=\underset{\sigma ^{}}{}L(\sigma ,\sigma ^{})P(\sigma ,t)$$ (1) The dynamical matrix $`L(\sigma ,\sigma ^{})=L^{}(\sigma ,\sigma ^{})+_0^{\mathrm{}}K(\sigma ,\sigma ^{},\tau )𝑑\tau `$ is determined by the above discussed formal procedure. Unfortunately, a direct calculation is mostly very complicated, so that one should use reasonable assumptions about the mathematical structure of $`L`$. To make the time evolution of the glass configurations more transparent we use the argumentation following the idea of Fredrickson and Andersen , i.e. we suppose that the basic dynamics is a simple (Glauber) process $`\sigma _j=+1\sigma _j=1`$ controlled by the thermodynamical Gibb’s measure and by self induced topological restrictions. In particular, an elementary flip at a given cell is allowed only if the number of the nearest neighbored mobile cells ($`\sigma _j=+1`$) is equal or larger than a restriction number $`n`$ with $`0<n<z_c`$ ($`z_c`$: coordination number of the lattice). Elementary flip processes and geometrical restrictions lead to the cooperative rearrangement of the underlying system and therefore to mesoscopical models describing a supercooled liquid below $`T_c`$. These models are denoted as $`n`$–spin facilitated Ising model on a $`d`$–dimensional lattice SFM$`[n,d]`$. The self–adapting environments influence in particular the long time behavior of the spin-spin correlation functions and therefore of the corresponding density-density correlation functions. The SFM$`[n,d]`$ was studied numerically (SFM$`[2,2]`$) and recently also analytically (SFM$`[1,1]`$). From this point of view it will be an interesting task to derive a set of equations related to the SFM$`[n,d]`$ which are similar to the well known Mori–Zwanzig equations and which can be used as a reasonable basis for a further treatment, e.g. for a continuous fraction analysis or for a mode coupling approach. The aim of the present paper is the derivation of such evolution equations and their analysis in terms of a mode coupling approach. We restrict our investigation to the analysis of the SFM$`[2,d]`$ but a generalization to another class of spin facilitated kinetic Ising models is always possible. The starting point is the mapping of the master equations of the SFM$`[2,d]`$ to evolution equations in a Fock–space representation. Using a projection formalism one obtains evolution equations for a set of relevant observables and consequently for the corresponding correlation functions. The paper ends in an analysis of these correlation functions in terms of the frequency matrices and memory terms. ## II Fock–space approach Following Doi , compare also , the probability distribution $`P(\sigma ,t)`$ can be related to a state vector $`\|F(t)`$ in a Fock-space according to $`P(\sigma ,t)=\sigma \|F(t)`$ and $`\|F(t)=_\sigma P(\sigma ,t)\|\sigma `$, respectively, with the base vectors $`\|\sigma `$. Using this representation, the Master equation (1) can be transformed to an equivalent equation in a Fock-space $$_t\|F(t)=\widehat{L}\|F(t)$$ (2) The dynamical matrix $`L(\sigma ,\sigma ^{})`$ of (1) is now mapped onto the operator $`\widehat{L}`$ given in a second quantized form with $`d`$ and $`d^{}`$ being the annihilation and creation operators, respectively, for flips processes. Usually $`\widehat{L}`$ is expressed in terms of creation and annihilation operators which satisfy Bose commutation rules . The SFM$`[n,d]`$ can be interpreted as a lattice gas ($`\sigma _i=0`$: empty cell, $`\sigma _i=1`$: occupied cell) considering the excluded volume effect, i.e. changes of the configuration $`\sigma `$ are possible only under the presence of the exclusion principle. To preserve the restriction of the occupation number in the underlying dynamical equations too, the commutation rules of the operators $`\widehat{d}`$ and $`\widehat{d}^{}`$ are chosen as those of Pauli-operators : $$[\widehat{d}_i,\widehat{d}_j^{}]=\delta _{i,j}(12\widehat{d}_i^{}\widehat{d}_i)[\widehat{d}_i,\widehat{d}_j]=[\widehat{d}_i^{},\widehat{d}_i^{}]=0\widehat{d}_i^2=(\widehat{d}_i^{})^2=0$$ (3) It should be remarked that the method can be extended to the case of higher restricted occupation numbers . As it was shown by Doi the average of a physical quantity $`B(\sigma )`$ is given by the average of the corresponding operator $`\widehat{B}(t)=_\sigma \|\sigma B(\sigma )\sigma \|`$ via $$\widehat{B}(t)=\underset{\sigma }{}P(\sigma ,t)B(\sigma )=s\left\|\widehat{B}\right\|F(t)$$ (4) using the reference state $`s\|=_\sigma \sigma \|`$. The normalization condition is manifested in the relation $`<s\|F(t)>=1`$. In the same way, arbitrary correlation functions can be expressed by $`\widehat{A}(t)\widehat{B}(t^{})={\displaystyle \underset{\sigma ,\sigma ^{}}{}}A(\sigma )P(\sigma ,t;\sigma ^{},t^{})B(\sigma ^{})=s\left\|\widehat{A}\mathrm{exp}\left\{\widehat{L}\left(tt^{}\right)\right\}\widehat{B}\right\|F(t^{})`$ From this point of view, it is possible to create the evolution equations of various averages and correlation functions. For example, using (2) and (4) one obtains the evolution equation for an arbitrary operator $`\widehat{B}`$ : $$_t\widehat{B}=s\left\|[\widehat{B},\widehat{L}]\right\|F(t)$$ (5) Here we have used the necessary relation $`<s\|\widehat{L}=0`$, which is an immediately consequence of the normalization condition. The evolution operator for the SFM$`[2,d]`$ can be written as $$\widehat{L}=\underset{i,j,k}{}\kappa _{i\|jk}\widehat{D}_j\widehat{D}_k\left[\beta (\widehat{d}_i\widehat{D}_i)+\lambda (\widehat{d}_i^{}(1\widehat{D}_i))\right]$$ (6) with the particle number operator $`\widehat{D}_i=\widehat{d}_i^{}\widehat{d}_i`$ and temperature dependent jumping rates $`\lambda `$ and $`\beta `$. $`\kappa _{i\|jk}`$ is a lattice function with $`\kappa _{i\|jk}=1`$ if $`jk`$ and $`j`$ and $`k`$ are neighbored to lattice cell $`i`$. Applying a simple activation dynamics one obtains for the jumping rates: $$\beta =\nu ^1(T)\text{ and}\lambda =\nu ^1(T)\mathrm{exp}(\epsilon /T)$$ (7) where $`\nu ^1(T)`$ is an elementary temperature dependent time scale ($`\epsilon `$ is the energy difference between the solid and liquid like state). Note that the stationary state corresponds to an average $`\overline{\sigma }_{\mathrm{eq}}=\widehat{D}_j=\lambda /(\lambda +\beta )`$ which can be obtained directly from (5). The knowledge of $`\widehat{L}`$ and the corresponding evolution equation (2) allows a reasonable analysis of the SFM$`[2,d]`$. The Fock space formalism has the decisive advantage of a simple construction principle for each evolution operator $`\widehat{L}`$ on the basis of creation and annihilation operators. Therefore, this method allows investigations of master equations for various evolution processes, e.g. aggregation, chemical reactions , nonlinear diffusion and just the spin facilitated kinetic Ising model. ## III Projection equations ### A relevant operators The dynamics of an arbitrary physical system can be described by a reasonable set of relevant operators. We use the normalized local deviations of the configuration from the thermodynamical average and the corresponding derivatives with respect to the time as suitable relevant observables for the investigation of the SFM$`[2,d]`$ $$\widehat{\eta }_i^{(0)}(t)=\widehat{\eta }_i(t)=\frac{\widehat{D}_i(t)\overline{\sigma }_{\mathrm{eq}}}{\sqrt{\overline{\sigma }_{\mathrm{eq}}(1\overline{\sigma }_{\mathrm{eq}})}}\mathrm{and}\widehat{\eta }_i^{(\beta )}(t)=\frac{^\beta \widehat{\eta }_i(t)}{t^\beta }=\widehat{\eta }_i(t)\widehat{L}^\beta $$ (8) ($`\beta =0,1,\mathrm{},g_{\mathrm{max}}`$; usually the upper borderline is a finite integer number, but $`g_{\mathrm{max}}\mathrm{}`$ is also possible). These covariant operators must be extended by the corresponding contravariant operators $$\stackrel{~}{\eta }_i^{(0)}(t)=\widehat{\eta }_i(t)\mathrm{and}\stackrel{~}{\eta }_i^{(\beta )}(t)=\widehat{L}^\beta \widehat{\eta }_i(t)$$ (9) Using (8) and (9) we construct the backward projection operator $`\widehat{P}`$: $$\widehat{P}=\underset{\alpha ,\beta ,i,j}{}\mathrm{}..\stackrel{~}{\eta }_i^{(\alpha )}g_{ij}^{\alpha \beta }\widehat{\eta }_j^{(\beta )}\mathrm{with}\underset{\alpha ,i}{}\widehat{\eta }_k^{(\gamma )}\stackrel{~}{\eta }_i^{(\alpha )}g_{ij}^{\alpha \beta }=\delta ^{\gamma \beta }\delta _{kj}$$ (10) (with $`\alpha `$, $`\beta `$, …$`[0,g_{\mathrm{max}}]`$). The projection operator leads to an identical mapping of the relevant operators onto itself, i.e. $`\widehat{\eta }_k^{(\gamma )}\widehat{P}=\widehat{\eta }_k^{(\gamma )}`$. Consequently, the orthogonal projection operator $`\widehat{Q}`$ is given by $`\widehat{Q}=1\widehat{P}`$ with $`\widehat{\eta }_k^{(\gamma )}\widehat{Q}=0`$. ### B Basis equations The evolution equation (2) leads to the formal solution $`\|F(t)=\mathrm{exp}\left\{\widehat{L}t\right\}\|F(0)`$. The dependence of $`\|F(t)`$ on the time can be transferred to an arbitrary operator analogous to the transformation of Schrödinger’s representation into the Heissenberg picture. Therefore one obtains time–dependent operators $$\widehat{B}(t)=\widehat{B}\mathrm{exp}\left\{\widehat{L}t\right\}$$ (11) The derivation of the evolution equations for the relevant observables starts from the formal time evolution of $`\widehat{\eta }_k^{(\gamma )}(t)`$. It follows from (11) $$\frac{\widehat{\eta }_k^{(\gamma )}(t)}{t}=\widehat{\eta }_k^{(\gamma )}(t)\widehat{L}$$ (12) This equation is the basis for the derivation of projection equations for the relevant observables in analogy to the well known Mori–Zwanzig equations for classical or quantum mechanical equations of motion. The application of $`1=\widehat{P}+\widehat{Q}`$ onto the operator $`\widehat{L}`$ leads to a formal splitting into a relevant and an irrelevant part (Note that $`\widehat{P}`$ realizes a projection onto the subspace $`L_{}`$ of relevant operators, whereas $`\widehat{Q}`$ projects onto the linearly independent subspace $`L_{}`$ of all other operators). Hence, $$\begin{array}{ccc}\frac{\widehat{\eta }_k^{\left(\gamma \right)}\left(t\right)}{t}\hfill & =\hfill & \widehat{\eta }_k^{(\gamma )}\mathrm{exp}\left\{\widehat{L}t\right\}\widehat{L}=\widehat{\eta }_k^{(\gamma )}\widehat{L}(\widehat{P}+\widehat{Q})\mathrm{exp}\left\{\widehat{L}t\right\}\hfill \\ & =\hfill & \underset{\beta ,j}{}\mathrm{\Omega }_{kj}^{(\gamma \beta )}\widehat{\eta }_j^{(\beta )}+\widehat{\eta }_k^{(\gamma )}\widehat{L}\widehat{Q}\mathrm{exp}\left\{\widehat{L}t\right\}\hfill \end{array}$$ (13) with the frequency matrix $$\mathrm{\Omega }_{kj}^{(\gamma \beta )}=\underset{\alpha ,i}{}\widehat{\eta }_k^{(\gamma )}\widehat{L}\stackrel{~}{\eta }_i^{(\alpha )}g_{ij}^{\alpha \beta }$$ (14) The second term of (13) can be rewritten by using an identical transformation of $`\mathrm{exp}\left\{\widehat{L}t\right\}`$ into an integral expression: $$\mathrm{exp}\left\{\widehat{L}t\right\}=_0^t𝑑t^{}\mathrm{exp}\left\{\widehat{L}\widehat{Q}(tt^{})\right\}\widehat{L}\widehat{P}\mathrm{exp}\left\{\widehat{L}t^{}\right\}+\mathrm{exp}\left\{\widehat{L}\widehat{Q}t\right\}$$ (15) This relation allows the derivation of rigorous projection equations similar to the usual Mori-Zwanzig-equations : $$\frac{\widehat{\eta }_k^{(\gamma )}(t)}{t}=\underset{\beta ,j}{}\mathrm{\Omega }_{kj}^{(\gamma \beta )}\widehat{\eta }_j^{(\beta )}+_0^t𝑑t^{}\underset{\beta ,j}{}K_{kj}^{(\gamma \beta )}(tt^{})\widehat{\eta }_j^{(\beta )}(t^{})+\widehat{f}_k^{(\gamma )}(t)$$ (16) with the residual forces $$\widehat{f}_k^{(\gamma )}(t)=\widehat{\eta }_k^{(\gamma )}(t)\widehat{L}\widehat{Q}\mathrm{exp}\left\{\widehat{L}\widehat{Q}t\right\}=\widehat{f}_k^{(\gamma )}\mathrm{exp}\left\{\widehat{L}\widehat{Q}t\right\}$$ (17) (with the properties $`\widehat{f}_k^{(\gamma )}(t)\widehat{Q}=\widehat{f}_k^{(\gamma )}(t)`$ and $`\widehat{f}_k^{(\gamma )}(t)\widehat{P}=0`$) and the memory matrix: $$K_{kj}^{(\gamma \beta )}(tt^{})=\underset{\alpha ,i}{}\widehat{\eta }_k^{(\gamma )}\widehat{L}\widehat{Q}\mathrm{exp}\left\{\widehat{L}\widehat{Q}(tt^{})\right\}\widehat{L}\stackrel{~}{\eta }_i^{(\alpha )}g_{ij}^{\alpha \beta }$$ (18) The comparison between (16) and the usual Mori-Zwanzig-equations shows a formal equivalence. Both equations contains frequency terms, memory and residual forces with a similar mathematical structure. But there is a fundamental difference which can be studied directly by inspecting the memory kernel. The memory of Mori-Zwanzig-equations can be written always as a correlation function of the residual forces. This relation can be interpreted as a representation of the fluctuation–dissipation theorem, and it is causally connected with the fact, that the Mori–Zwanzig equations are related to reversible classical or quantum mechanical equations. On the other hand, the memory (18) cannot be completely constructed from residual forces (17). The cause is the irreversible character of the underlying master equation. ### C Projection equations for a reduced set of relevant observables We choose $`g_{\mathrm{max}}=1`$ for the following investigations, i.e. the relevant observables are $`\widehat{\eta }_i^{(0)}(t)=\widehat{\eta }_i(t)`$ and $`\widehat{\eta }_i^{(1)}(t)=\widehat{\eta }_i(t)\widehat{L}`$. This settling corresponds slightly to mechanical systems, which are completely determined by spatial coordinates and velocities. The general system of equations (16) becomes $$\begin{array}{ccc}\frac{}{t}\widehat{\eta }_n^{(0)}(t)\hfill & =\hfill & \underset{j}{}\underset{\beta =0,1}{}\left[\mathrm{\Omega }_{nj}^{(0\beta )}\widehat{\eta }_j^{(\beta )}(t)+_0^t𝑑t^{}K_{nj}^{(0\beta )}(tt^{})\widehat{\eta }_j^{(\beta )}(t^{})\right]+\widehat{f}_n^{(0)}(t)\hfill \\ \frac{}{t}\widehat{\eta }_n^{(1)}(t)\hfill & =\hfill & \underset{j}{}\underset{\beta =0,1}{}\left[\mathrm{\Omega }_{nj}^{(1\beta )}\widehat{\eta }_j^{(\beta )}(t)+_0^t𝑑t^{}K_{nj}^{(1\beta )}(tt^{})\widehat{\eta }_j^{(\beta )}(t^{})\right]+\widehat{f}_n^{(1)}(t)\hfill \end{array}$$ (19) A simple analysis leads to the simplifications $`\mathrm{\Omega }_{nj}^{(0\beta )}=\delta _{nj}\delta ^{1\beta }`$, $`\widehat{f}_n^{(0)}(t)=0`$ and $`K_{nj}^{(0\beta )}(tt^{})=0`$. Therefore, the first equation will be reduced to the identity $`_t\widehat{\eta }_n^{(0)}(t)=_t\widehat{\eta }_n(t)=\widehat{\eta }_n^{(1)}(t)`$ and the second equation can be rewritten as $$\begin{array}{ccc}\frac{^2}{t^2}\widehat{\eta }_n(t)\hfill & =\hfill & \underset{j}{}\left[\mathrm{\Omega }_{nj}^{(10)}\widehat{\eta }_j(t)+\mathrm{\Omega }_{nj}^{(11)}\frac{}{t}\widehat{\eta }_j(t)\right]\hfill \\ & +\hfill & \underset{j}{}_0^t𝑑t^{}\left[K_{nj}^{(10)}(tt^{})\widehat{\eta }_j(t)+K_{nj}^{(11)}(tt^{})\frac{}{t^{}}\widehat{\eta }_j(t^{})\right]+\widehat{f}_n^{(1)}(t)\hfill \end{array}$$ (20) The result is a second order differential equation which reflects the complete dynamics of the relevant observables. ### D Projection equations for correlation functions An important rule for experimental and theoretical investigations of the glass transition plays the time–dependent equilibrium correlation functions of the relevant observables. Especially the SFM$`[n,d]`$ should be characterized by $$\mathrm{\Phi }_{nm}(t)=\widehat{\eta }_n(t)\widehat{\eta }_m(0)=s\left\|\widehat{\eta }_n\mathrm{exp}\left\{\widehat{L}t\right\}\widehat{\eta }_m\right\|F(0)$$ (21) These correlation function is equivalent to the normalized spin-spin correlation: $`\mathrm{\Phi }_{nm}(t)={\displaystyle \frac{\sigma _n(t)\sigma _m(0)\overline{\sigma }_{\mathrm{eq}}^2}{\overline{\sigma }_{\mathrm{eq}}(1\overline{\sigma }_{\mathrm{eq}})}}`$ which should be similar to the normalized density–density correlation of the underlying supercooled liquid, i.e. $`\mathrm{\Phi }_{nm}(t)\delta \rho (𝐫,t)\delta \rho (𝐫^{},0)`$, see also the above discussed mapping $`\rho \sigma `$. The evolution equation of $`\mathrm{\Phi }_{nm}(t)`$ follows from (20) by a right hand multiplication with $`\widehat{\eta }_m`$ and a subsequent determination of the average. The contributions of the residual forces $`\widehat{f}_n^{(1)}(t)`$ vanish identically. One obtains a homogeneous integro–differential equation: $$\begin{array}{ccc}\frac{^2\mathrm{\Phi }_{nm}}{t^2}\hfill & =\hfill & \underset{j}{}\left[\mathrm{\Omega }_{nj}^{(10)}\mathrm{\Phi }_{jm}+\mathrm{\Omega }_{nj}^{(11)}\frac{\mathrm{\Phi }_{jm}}{t}\right]\hfill \\ & +\hfill & _0^t𝑑t^{}\left[K_{nj}^{(10)}(tt^{})\mathrm{\Phi }_{jm}(t^{})+K_{nj}^{(11)}(tt^{})\frac{\mathrm{\Phi }_{jm}(t^{})}{t}\right]\hfill \end{array}$$ (22) This evolution equation of $`\mathrm{\Phi }_{nm}(t)`$ is a rigorous second order integro–differential equation, which will be analyzed now. To this aim it is necessary to determine the frequency and memory parts for the thermodynamical equilibrium. Here, correlation functions, memory and frequency matrices should be homogeneous and isotropic functions, e.g. $`\mathrm{\Phi }_{mn}(t)=\mathrm{\Phi }(\left\|𝐧𝐦\right\|,t)`$ or $`K_{mn}^{(\alpha \beta )}(t)=K^{(\alpha \beta )}(\left\|𝐧𝐦\right\|,t)`$. The homogenity is a direct consequence of the underlying translation invariance. On the other hand, isotropy can be expected for the asymptotic case of the continuous limit, i.e. for $`\left\|𝐧𝐦\right\|\mathrm{}`$. However, the isotropy is partially disturbed at finite distances as a consequence of the underlying lattice structure. The Fourier transformation of (22) can be obtained by using the representation $`\mathrm{\Phi }_{mn}(t)={\displaystyle \underset{𝐪}{}}\mathrm{\Phi }(𝐪,t)\mathrm{exp}\left\{i𝐪(𝐧𝐦)\right\}\mathrm{\Phi }(𝐪,t)={\displaystyle \frac{1}{N}}{\displaystyle \underset{n}{}}\mathrm{\Phi }(\|𝐧\|,t)\mathrm{exp}\left\{i\mathrm{𝐪𝐧}\right\}`$ of the correlation function whereas the frequency matrix and the memory can be written as: $$\mathrm{\Omega }^{(il)}(𝐪)=\frac{1}{N}\underset{n}{}\mathrm{\Omega }^{(il)}(\|𝐧\|)\mathrm{exp}\left\{i\mathrm{𝐪𝐧}\right\}$$ (23) and $$K^{(il)}(𝐪,tt^{})=\frac{1}{N}\underset{n}{}K^{(il)}(\|𝐧\|,tt^{})\mathrm{exp}\left\{i\mathrm{𝐪𝐧}\right\}$$ (24) respectively. $`𝐧`$ and $`𝐦`$ denote the lattice vector of size $`n`$ and $`m`$, $`𝐪`$ is a vector of the first Brillouin zone corresponding to the lattice. We use a cubic lattice for the following calculations but an application of another lattice type is always possible. The Mori-Zwanzig-equation (22) becomes now: $$\begin{array}{ccc}\ddot{\mathrm{\Phi }}(q,t)\hfill & =\hfill & N\left[\mathrm{\Omega }^{(10)}(q)\mathrm{\Phi }(q,t)+\mathrm{\Omega }_{nj}^{(11)}(q)\dot{\mathrm{\Phi }}(q,t)\right]\hfill \\ & +\hfill & N_0^t𝑑t^{}\left[K^{(10)}(q,tt^{})\mathrm{\Phi }(q,t^{})+K^{(11)}(q,tt^{})\dot{\mathrm{\Phi }}(q,t^{})\right]\hfill \end{array}$$ (25) Note that all quantities depend only on $`q=\left\|𝐪\right\|`$ (at least for the continuous limit) because of the isotropy. Finally, the Laplace transformation $$\mathrm{\Phi }(q,z)=_0^{\mathrm{}}𝑑t\mathrm{exp}\left\{zt\right\}\mathrm{\Phi }(q,t)$$ (26) leads to the algebraic equation $$\mathrm{\Phi }(q,z)=\frac{\mathrm{\Phi }_0(q)}{z+{\displaystyle \frac{N\mathrm{\Omega }^{(10)}(q)NK^{(10)}(q,z)zg_0(q)}{z+N\mathrm{\Omega }^{(11)}(q)NK^{(11)}(q,z)+g_0(q)}}}$$ (27) which considers the initial conditions $`\mathrm{\Phi }(q,0)=\mathrm{\Phi }_0(q)`$ and $`\dot{\mathrm{\Phi }}(q,0)=\dot{\mathrm{\Phi }}_0(q)`$. Furthermore, the quantity $`g_0(q)`$ denotes the ratio $`g_0(q)=\dot{\mathrm{\Phi }}_0(q)/\mathrm{\Phi }_0(q)`$. It should remarked that especially $`\mathrm{\Omega }^{(11)}(q)0`$ and $`K^{(10)}(q,z)0`$ are consequences of the irreversible master equations. On the other hand, the usual Mori–Zwanzig equations are founded on reversible Liouville operators which lead immediately to $`\mathrm{\Omega }^{(11)}(q)=0`$ and $`K^{(10)}(q,z)=0`$. ## IV Determination of frequency matrices The concrete determination of the frequency matrices is possible by using the concrete evolution operator $`\widehat{L}`$ (6). Note that the projection equations (22) are valid for an arbitrary physical system which can be described by master equations. The frequency matrices contain always the matrix $`𝐠`$, which can be calculated from (10), i.e. $`𝐠`$ is determined by the following system of linear equations: $$h_{ik}^{(\alpha \gamma )}=\widehat{\eta }_i^{(\alpha )}\stackrel{~}{\eta }_k^{(\gamma )}\mathrm{with}\underset{\gamma ,k}{}h_{ik}^{(\alpha \gamma )}g_{kj}^{\gamma \beta }=\delta ^{\alpha \beta }\delta _{ij}$$ (28) Using the definition $$\mathrm{\Gamma }_{ik}^\beta =\widehat{\eta }_i\widehat{L}^\beta \widehat{\eta }_k$$ (29) one obtains simple expressions for the matrix $`h`$ (i.e. $`h_{ik}^{(\alpha \gamma )}=\mathrm{\Gamma }_{ik}^{\alpha +\gamma }`$) as well as for the frequency matrix (14): $$\mathrm{\Omega }_{kj}^{(\gamma \beta )}=\underset{\alpha ,i}{}\mathrm{\Gamma }_{ki}^{\alpha +\gamma +1}g_{ij}^{\alpha \beta }$$ (30) The knowledge of $`\mathrm{\Gamma }_{ik}^\beta `$ ($`\beta =0\mathrm{}3`$) allows the exact determination of $`\mathrm{\Omega }^{(10)}(q)`$ and $`\mathrm{\Omega }_{nj}^{(11)}(q)`$. The values $`\mathrm{\Gamma }_{ik}^\beta `$ can be obtained straightforwardly by using (6) and the commutation relations (3). It follows: $$\mathrm{\Gamma }_{mn}^\alpha =\left(\frac{1}{\tau _0}\right)^\alpha \left(A^\alpha \delta _{mn}+B^\alpha \mathrm{\Theta }_{nm}+C^\alpha \chi _{mn}+D^\alpha \zeta _{nm}\right)$$ (31) The lattice functions $`\mathrm{\Theta }_{nm}`$, $`\chi _{mn}`$ and $`\zeta _{nm}`$ vanish, except for the following cases: $`\mathrm{\Theta }_{nm}=1`$ for $`\left\|𝐦𝐧\right\|=1`$, $`\chi _{mn}=1`$ for $`\left\|𝐦𝐧\right\|=\sqrt{2}`$ and $`\zeta _{nm}=1`$ for $`\left\|𝐦𝐧\right\|=2`$. The values $`A^\alpha `$, $`B^\alpha `$, $`C^\alpha `$ and $`D^\alpha `$ are listed in appendix A. The Fourier transformation is now a simple calculation. The approximation for small wave vectors (continuous limit) is by a special interest. The actual lattice structure becomes irrelevant on these sufficiently large spatial scales, i.e. the Fourier transformed $`\mathrm{\Gamma }_{mn}^\alpha `$ are isotropic values. One obtains up to the second order in $`q`$: $$\mathrm{\Gamma }^\alpha (q)=\frac{1}{N}\frac{1}{(\tau _0)^\alpha }\left(\mathrm{\Gamma }_0^\alpha \gamma ^\alpha q^2\right)$$ (32) The coefficients $`\mathrm{\Gamma }_0^\alpha `$ and $`\gamma ^\alpha `$ follows immediately from the values $`A^\alpha `$, $`B^\alpha `$, $`C^\alpha `$ and $`D^\alpha `$: $$\mathrm{\Gamma }_0^\alpha =A^\alpha +z_cB^\alpha +\frac{1}{2}z_c(z_c2)C^\alpha +z_cD^\alpha \mathrm{and}\gamma ^\alpha =B^\alpha +(z_c2)C^\alpha +4D^\alpha $$ (33) ($`z_c`$ is the coordination number of the $`d`$-dimensional lattice, i.e. $`z_c`$ is the number of nearest neighbors per lattice cell. Straightforwardly, the Fourier transformed matrix $`g^{\alpha \beta }(q)`$ can be written as $$g^{(\alpha \beta )}(q)=\frac{1}{WN}\left(\begin{array}{cc}\mathrm{\Gamma }_0^2\gamma ^2q^2& \tau _0\mathrm{\Gamma }_0^1\\ \tau _0\mathrm{\Gamma }_0^1& \tau _0^2\end{array}\right)$$ (34) with $`W=\mathrm{\Gamma }_0^2(\mathrm{\Gamma }_0^1)^2\gamma ^2q^2`$. Finally, the frequency matrices $`\mathrm{\Omega }^{(10)}(q)`$ and $`\mathrm{\Omega }^{(11)}(q)`$ are given by $$N\mathrm{\Omega }^{(10)}(q)=\frac{1}{W\tau _0^2}\left\{\mathrm{\Gamma }_0^1\mathrm{\Gamma }_0^3\left(\mathrm{\Gamma }_0^2\right)^2+\left(2\mathrm{\Gamma }_0^2\gamma ^2\mathrm{\Gamma }_0^1\gamma ^3\right)q^2\right\}$$ (35) and $$N\mathrm{\Omega }^{(11)}(q)=\frac{1}{W\tau _0}\left\{\mathrm{\Gamma }_0^3\mathrm{\Gamma }_0^1\mathrm{\Gamma }_0^2+\left(\mathrm{\Gamma }_0^1\gamma ^2\gamma ^3\right)q^2\right\}$$ (36) in the continuous limit. Note, that (34), (35) and (36) considers already that $`\gamma ^0=\gamma ^1=0`$ and $`\mathrm{\Gamma }_0^0=1`$. ## V Analysis of the relaxation behavior A rough understanding of the results so far is possible by an analysis of the relaxation behavior of the correlation function $`\mathrm{\Phi }(q,t)`$ neglecting the memory terms containing in (27). In this case the correlation function is reduced to a finite continued fraction. It can be expected that this case is related to the high temperature limit corresponding to a more or less exponential decay of the correlation function. Furthermore, this approximation should be reasonable for the description of the correlation function at short time scales. Note that because of $`K^{(10)}(q,t)`$const. and $`K^{(11)}(q,t)`$const. for $`t0`$, one obtains $`K^{(10)}(q,z)z^1`$ and $`K^{(11)}(q,z)z^1`$ for $`z\mathrm{}`$. Therefore, the memory terms can be neglected at sufficiently short time scales $`t0`$ or $`z\mathrm{}`$. The initial conditions $`\mathrm{\Phi }_0(q)`$ and $`\dot{\mathrm{\Phi }}_0(q)`$ of the correlation function $`\mathrm{\Phi }`$ are defined by equilibrium averages: $$\mathrm{\Phi }_{nm}(0)=\widehat{\eta }_n\widehat{\eta }_m=\mathrm{\Gamma }_{nm}^0\mathrm{and}\dot{\mathrm{\Phi }}_{nm}(0)=\widehat{\eta }_m\widehat{L}\widehat{\eta }_n=\mathrm{\Gamma }_{nm}^1$$ (37) and consequently $$g_0(q)=\frac{\dot{\mathrm{\Phi }}_0(q)}{\mathrm{\Phi }_0(q)}=N\mathrm{\Gamma }^1(q)$$ (38) The normalized correlation function $`\stackrel{~}{\mathrm{\Phi }}(q,z)=\mathrm{\Phi }(q,z)/\mathrm{\Phi }_0(q)`$ follows from (27) under consideration of (38) and under neglect of memory terms. A simple calculation leads to $$\stackrel{~}{\mathrm{\Phi }}(q,z)=\frac{z+N\mathrm{\Omega }^{(11)}(q)+N\mathrm{\Gamma }^1(q)}{z^2+zN\mathrm{\Omega }^{(11)}(q)+N\mathrm{\Omega }^{(10)}(q)}$$ (39) which can be written as $$\stackrel{~}{\mathrm{\Phi }}(q,z)=\frac{A_1}{zz_1}+\frac{A_2}{zz_2}$$ (40) with the poles $$z_{1/2}=\frac{1}{2}\left[N\mathrm{\Omega }^{(11)}(q)\sqrt{\left(N\mathrm{\Omega }^{(11)}(q)\right)^24N\mathrm{\Omega }^{(10)}(q)}\right].$$ (41) and the intensities $`A_1={\displaystyle \frac{N\mathrm{\Gamma }^1(q)z_2}{z_1z_2}}A_2={\displaystyle \frac{z_1N\mathrm{\Gamma }^1(q)}{z_1z_2}}`$ The present approximation of the SFM$`[2,d]`$ is characterized by two relaxation times $`\tau _R^1(q)=z_1^1`$ and $`\tau _R^2(q)=z_2^1`$. It can be verified by a simple calculation that both relaxation times shows only a weak dependence on $`q.`$ Furthermore, both relaxation times approach finite values for $`q0`$. Obviously, the spin facilitated kinetic Ising model shows no diffusion–like modes which behave as $`\tau q^2`$ for the limit $`q0`$. This finding agrees with investigations of the one dimensional spin facilitated kinetic Ising model . The SFM$`[1,1]`$ corresponds to diffusion processes combined with creation and annihilation processes of active states. A cell has an active state if this cell can change its state without any support by further flip processes of neighbored cells. Creation and annihilation processes dominates at sufficiently large scales, i.e. an inhomogeneity reaching over a sufficiently large distance will be reduced by local creation and annihilation processes of mobile cells until diffusion processes becomes effective. As above mentioned, we restrict our investigations to the borderline case of macroscopic scales, i.e. $`q0`$. Fig.1 shows the relaxation times $`\tau _R^1(q)`$ and $`\tau _R^2(q)`$ as a function of temperature. As expected, there is no significant difference between the relaxation times for $`q=0`$ and $`q0`$, respectively. Furthermore, the slow relaxation time $`\tau _R^1(q)`$ shows a weak non–Arrhenius behavior. One obtains $`\mathrm{ln}\tau _R^2\mathrm{ln}\nu (T)+o(q^2)`$ (see also eq.7) and $`\mathrm{ln}\tau _R^1\mathrm{ln}\nu (T)+2\epsilon /T+o(q^2)`$, respectively, at low temperatures. On the other hand, the high temperature regime is characterized by another temperature dependence: $`\mathrm{ln}\tau _R^{1,2}=\mathrm{ln}\nu (T)+u^{1,2}\epsilon /T+o(q^2)`$. The coefficients $`u^{1,2}`$ depend on the actual lattice structure and the spatial dimension. But a simple analysis shows that always $`u^1<2`$, i.e. the activation energy of the SFM$`[2,d]`$ increases with decreasing temperature. The existence of two relaxation times means not that the SFM$`[2,d]`$ is characterized by an $`\alpha `$– and a $`\beta `$–process. The superposition of both decays, $`\mathrm{exp}\left\{t/\tau _R^1\right\}`$ and $`\mathrm{exp}\left\{t/\tau _R^2\right\}`$, considering of the intensities $`A_1`$ and $`A_2`$ (see fig.2), shows a continuous decay of the correlation function $`\stackrel{~}{\mathrm{\Phi }}(q,t)`$, see fig.3. This behavior is in an agreement with numerical simulations and it corresponds also to the above discussed thought, that spin facilitated kinetic Ising models are possible candidates modelling the behavior of supercooled liquids below the critical temperature of the usual mode coupling theory. ## VI Determination of the memory matrices ### A Complete and orthogonal basis All operators acting on the Fock–space can be represented by a complete collection of orthogonal base operators. The determination of such a basis is possible under consideration of the underlying $`\widehat{d}_i,\widehat{d}_i^{}`$– (pseudo fermionic) algebra (3) of the SFM$`[2,d]`$. The base operators can be expressed as all possible products of the above introduced operators $`\widehat{\eta }_i`$ (8). A base operator is denoted as $`\widehat{B}_{𝐍_n}^{(n)}`$. (The index $`n`$ corresponds to the order of the product, $`𝐍_n`$ is an $`n`$–dimensional vector indicating the concerning lattice cells). The first groups of the basis are: $`\begin{array}{ccccc}\widehat{B}^{(0)}\hfill & =\hfill & 1\hfill & & \\ \widehat{B}_i^{(1)}\hfill & =\hfill & \widehat{\eta }_i\hfill & & \\ \widehat{B}_{ij}^{(2)}\hfill & =\hfill & \widehat{\eta }_i\widehat{\eta }_j\hfill & \mathrm{for}\hfill & i<j\hfill \\ \widehat{B}_{ijk}^{(3)}\hfill & =\hfill & \widehat{\eta }_i\widehat{\eta }_j\widehat{\eta }_k\hfill & \mathrm{for}\hfill & i<j<k\hfill \end{array}`$ Note that because of the commutation relation $`[\widehat{\eta }_i,\widehat{\eta }_j]=0`$ the components of $`𝐍_n`$ can be ordered. The case of two or more equivalent indices is excluded because $`\widehat{\eta }_i^2=(12\overline{\sigma }_{\mathrm{eq}})/\overline{\sigma }_{\mathrm{eq}}(1\overline{\sigma }_{\mathrm{eq}})\widehat{\eta }_i+1`$, i.e. quadratic or higher powers of each operator $`\widehat{\eta }_i`$ can be always reduced to a linear representation. The orthogonality means that: $$\widehat{B}_{𝐍_n}^{(n)}\widehat{B}_{𝐍_m}^{(m)}=\delta ^{nm}\delta _{𝐍_n,𝐍_m}$$ (42) This relation can be checked considering that all equilibrium averages of operators on various cells decay in a product of averages with respect to these cells, e.g. $`\widehat{\eta }_i^2\mathrm{}\widehat{\eta }_j\mathrm{}..\widehat{\eta }_k^2\mathrm{}.=\widehat{\eta }_i^2\mathrm{}\widehat{\eta }_j\mathrm{}..\widehat{\eta }_k^2\mathrm{}.`$. This important relation is valid for all SFM$`[n,d]`$ if the neighbor–neighbor interaction vanishes. Note that the Hamiltonian of the analyzed class of spin facilitated kinetic Ising models is given by $`\widehat{H}=_i\epsilon \widehat{D}_i(t)_i\epsilon \sqrt{\overline{\sigma }_{\mathrm{eq}}(1\overline{\sigma }_{\mathrm{eq}})}\widehat{\eta }_i(t)`$. From this point of view, the relation (42) follows immediately because of $`\widehat{\eta }_i=0`$ and $`\widehat{\eta }_i^2=1`$. Thus, the basis $`\stackrel{~}{B}=\left\{\widehat{B}_{𝐍_n}^{(n)}\right\}`$ is orthogonal. The completeness of $`\stackrel{~}{B}=\left\{\widehat{B}_{𝐍_n}^{(n)}\right\}`$ is to be understood in relation to the reference state $`s\|`$, i.e. the following equation is fulfilled for an arbitrary operator $`\widehat{X}`$: $$s\|\widehat{X}=\underset{n}{}\underset{𝐍_n}{}\widehat{X}\widehat{B}_{𝐍_n}^{(n)}s\|\widehat{B}_{𝐍_n}^{(n)}$$ (43) The mathematical proof of this property is given in appendix B. ### B Decomposition of the memory terms Eq.18 is a reasonable starting point for an analysis of the memory terms. The consideration of (8) and (9) leads to: $$K_{kj}^{(\gamma \beta )}(t)=\underset{\alpha ,i}{}\widehat{\eta }_k\widehat{L}^{\gamma +1}\widehat{Q}\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\widehat{Q}\widehat{L}^{\alpha +1}\eta _ig_{ij}^{\alpha \beta }$$ (44) (Note that $`\widehat{Q}^2=\widehat{Q}`$). The application of the completeness relation (43) onto (44) yields: $$K_{kj}^{(\gamma \beta )}(t)=\underset{\alpha ,i}{}\underset{n,m}{}\underset{𝐍_n,𝐍_m}{}H_{k,𝐍_n}^{\gamma (n)}\widehat{B}_{𝐍_n}^{(n)}\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\widehat{B}_{𝐍_m}^{(m)}\stackrel{~}{H}_{𝐍_m.i}^{(m)\alpha }g_{ij}^{\alpha \beta }$$ (45) with $$H_{k,𝐍_n}^{\gamma (n)}=\widehat{\eta }_k\widehat{L}^{\gamma +1}\widehat{Q}\widehat{B}_{𝐍_n}^{(n)}\mathrm{and}\stackrel{~}{H}_{𝐍_n.k}^{(n)\gamma }=\widehat{B}_{𝐍_n}^{(n)}\widehat{Q}\widehat{L}^{\gamma +1}\eta _k$$ (46) These coefficients can be determined by simple algebraic calculations. One obtains immediately that both, $`H_{k,𝐍_n}^{\gamma (n)}`$ and $`\stackrel{~}{H}_{𝐍_n.i}^{(n)\alpha }`$ vanish identically for $`n=0,1`$. On the other hand, the coefficients (46) vanish also identically for $`n>5`$ because $`\gamma 1`$. Hence, the memory terms can be constructed by using a finite number of functions $`\widehat{B}_{𝐍_n}^{(n)}\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\widehat{B}_{𝐍_m}^{(m)}`$. These functions will be transformed identically. One obtains: $$\begin{array}{c}\widehat{B}_{𝐍_n}^{(n)}\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\widehat{B}_{𝐍_m}^{(m)}=\widehat{B}_{𝐍_n}^{(n)}\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\mathrm{exp}\left\{\widehat{L}t\right\}\mathrm{exp}\left\{\widehat{L}t\right\}\widehat{B}_{𝐍_m}^{(m)}\hfill \\ =_p_{𝐍_p}\widehat{B}_{𝐍_n}^{(n)}\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\mathrm{exp}\left\{\widehat{L}t\right\}\widehat{B}_{𝐍_p}^{(p)}\widehat{B}_{𝐍_p}^{(p)}\mathrm{exp}\left\{\widehat{L}t\right\}\widehat{B}_{𝐍_m}^{(m)}\hfill \\ =_p_{𝐍_p}\widehat{B}_{𝐍_n}^{(n)}\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\mathrm{exp}\left\{\widehat{L}t\right\}\widehat{B}_{𝐍_p}^{(p)}\widehat{B}_{𝐍_p}^{(p)}(t)\widehat{B}_{𝐍_m}^{(m)}\hfill \end{array}$$ (47) The averages $`\widehat{B}_{𝐍_p}^{(p)}(t)\widehat{B}_{𝐍_m}^{(m)}`$ are usual many point correlation functions. ### C Mode coupling approximation #### 1 Short time evolution of the memory An exact determination of $`\widehat{B}_{𝐍_n}^{(n)}\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\mathrm{exp}\left\{\widehat{L}t\right\}\widehat{B}_{𝐍_p}^{(p)}`$ and $`\widehat{B}_{𝐍_p}^{(p)}(t)\widehat{B}_{𝐍_m}^{(m)}`$ may be possible only for some few special cases. Therefore, we need a suitable approximation for a further treatment. In a first step we analyze the function: $$\mathrm{\Psi }_{𝐍_n𝐍_p}^{(np)}(t)=\widehat{B}_{𝐍_n}^{(n)}\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\mathrm{exp}\left\{\widehat{L}t\right\}\widehat{B}_{𝐍_p}^{(p)}$$ (48) (47) can be interpreted as a separation of fast and slow time scales. The operator $`\widehat{Q}\widehat{L}\widehat{Q}`$ is related to a dynamics completely different to the dynamics of $`\widehat{L}`$. In general, it can be expected that $`\widehat{B}_{𝐍_n}^{(n)}\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}`$ shows a significant evolution on a very short time scale in comparison to the characteristic time scale related to $`\widehat{B}_{𝐍_n}^{(n)}\mathrm{exp}\left\{\widehat{L}t\right\}`$. Therefore, we come to the rough conclusion: while the evolution operator $`\widehat{L}`$ contains all relevant time scales, the operator $`\widehat{Q}\widehat{L}\widehat{Q}`$ is mainly determined by contributions related to short time scales, i.e. the long time contributions are partially cancelled by the projection procedure. Consequently, the term $`\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\mathrm{exp}\left\{\widehat{L}t\right\}`$ and therefore $`\mathrm{\Psi }_{𝐍_n𝐍_p}^{(np)}(t)`$ should be dominated by long time scales, because the fast time scales are eliminated at least partially by the factor $`\mathrm{exp}\left\{\widehat{L}t\right\}`$. Clearly, this is a very rough interpretation, but it gives an explanation for the assumption that the time dependence of $`\mathrm{\Psi }_{𝐍_n𝐍_p}^{(np)}(t)`$ is weak in comparison to the decay of the correlation function $`\widehat{B}_{𝐍_p}^{(p)}(t)\widehat{B}_{𝐍_m}^{(m)}`$ which is connected only with the time evolution factor $`\mathrm{exp}\left\{\widehat{L}t\right\}`$. However, it seems to be reasonable to expand $`\mathrm{\Psi }_{𝐍_n𝐍_p}^{(np)}(t)`$ in powers of the time $`t`$. The determination of all Taylor coefficients is out of the question. But an approximative analysis is possible by using a finite number of coefficients. We hope that the error of this perturbation expansion is sufficiently strong suppressed by the corresponding factor $`\widehat{B}_{𝐍_p}^{(p)}(t)\widehat{B}_{𝐍_m}^{(m)}`$ (see eq.47). One obtains: $`\mathrm{\Psi }_{𝐍_n𝐍_p}^{(np)}(t)={\displaystyle \underset{M=0}{\overset{\mathrm{}}{}}}\mathrm{\Lambda }_{𝐍_n,𝐍_p}^{(np),M}{\displaystyle \frac{t^M}{M!}}`$ with: $`\mathrm{\Lambda }_{𝐍_n,𝐍_p}^{(np),M}=\widehat{B}_{𝐍_n}^{(n)}{\displaystyle \frac{^M}{t^M}}\left[\mathrm{exp}\left\{\widehat{Q}\widehat{L}\widehat{Q}t\right\}\mathrm{exp}\left\{\widehat{L}t\right\}\right]_{t=0}\widehat{B}_{𝐍_p}^{(p)}`$ The first coefficients $`\mathrm{\Lambda }_{𝐍_n,𝐍_p}^{(np),M}`$ can be determined by simple calculations, e.g. $$\begin{array}{ccc}\mathrm{\Lambda }_{𝐍_n,𝐍_p}^{(np),0}\hfill & =\hfill & \widehat{B}_{𝐍_n}^{(n)}\widehat{B}_{𝐍_p}^{(p)}=\delta ^{np}\delta _{𝐍_n,𝐍_p}\hfill \\ \mathrm{\Lambda }_{𝐍_n,𝐍_p}^{(np),1}\hfill & =\hfill & \widehat{B}_{𝐍_n}^{(n)}\left(\widehat{Q}\widehat{L}\widehat{Q}\widehat{L}\right)\widehat{B}_{𝐍_p}^{(p)}\hfill \end{array}$$ (49) In principle, the discussed expansion is an exact representation (The radius of convergence of the exponential function is infinite large). The true approximation consists in the breaking of the Taylor expansion after a finite power of $`t`$. We restrict our investigation to the simplest case, i.e. we assume $`\mathrm{\Lambda }_{𝐍_n,𝐍_p}^{(np),M}=0`$ for $`M1`$. Hence, one obtains $$\mathrm{\Psi }_{𝐍_n𝐍_p}^{(np)}(t)\delta ^{np}\delta _{𝐍_n,𝐍_p}$$ (50) But it should be remarked that an extension to higher terms is possible without any problems. We abstain from a consideration of higher terms with respect to the clarity of the calculations. Furthermore, the obtained results (see below) using (50) show already a reasonable agreement with numerical simulations. #### 2 Decomposition of the many point correlation functions The main problem is a reasonable approximation of the function $`\widehat{B}_{𝐍_p}^{(p)}(t)\widehat{B}_{𝐍_m}^{(m)}`$. This function decays in products of simple pair correlation functions if the distances between the corresponding lattice points (defined by the vectors $`𝐍_p`$ and $`𝐍_m`$) are sufficiently large: $$\widehat{B}_{(i_1i_2\mathrm{}i_p)}^{(p)}(t)\widehat{B}_{(j_1j_2\mathrm{}j_p)}^{(m)}\frac{1}{p!}(\mathrm{\Phi }_{i_1j_1}(t)\mathrm{\Phi }_{i_2j_2}(t)\mathrm{}.\mathrm{\Phi }_{i_pj_p}(t)+\mathrm{perm})$$ (51) This asymptotic limit is correct for infinitely large (or at least sufficiently large) distances between the lattice cells $`i_1`$, $`i_2`$, …. We use this borderline case as an approximation for an arbitrary set of lattice cells $`\{𝐍_p,𝐍_m\}`$. This approximation is equivalent to the decomposition of higher static correlation functions into simple pair correlation functions. For example, a similar approach was used to create self consistent equations for the static structure factor . Furthermore, this approximation is also the kernel of the well known mode coupling approach . #### 3 Reduction of the basis The third approximation consists in a reduction of the basis $`\stackrel{~}{B}=\left\{\widehat{B}_{𝐍_n}^{(n)}\right\}`$. It was demonstrated that only base operators with $`n5`$ are necessary for a representation of (46). Because of the approximations (50) and (51) all other factors of (45) contains also no higher base operators $`\widehat{B}_{𝐍_n}^{(n)}`$. The consideration of all relevant $`\widehat{B}_{𝐍_n}^{(n)}`$ is no general problem. However, we restrict the calculations only to elements with $`n2`$, also with respect to clarity. Really, the neglected terms containing base operators $`\widehat{B}_{𝐍_n}^{(n)}`$ with $`5n3`$ yield only small additional contributions to the final results. Thus, the complete representation of the memory (45) is given by the approximation: $$K_{kj}^{(\gamma \beta )}(t)\frac{1}{2}\underset{\alpha ,i}{}\underset{i_1i_2j_1j_2}{}H_{k,(i_1i_2)}^{\gamma (2)}\left[\mathrm{\Phi }_{i_1j_1}(t)\mathrm{\Phi }_{i_2j_2}(t)+\mathrm{\Phi }_{i_1j_2}(t)\mathrm{\Phi }_{i_2j_1}(t)\right]\stackrel{~}{H}_{(j_1j_2).i}^{(2)\alpha }g_{ij}^{\alpha \beta }$$ (52) ### D Macroscopic scale As above mentioned, the $`q`$–dependence of the frequency matrices is very small. This weak dependence can be expected also for the memory terms, i.e. we restrict out investigations only to the macroscopic scale $`q0`$. The assumption of a weak dependence on $`q`$ leads to: $$\mathrm{\Phi }_{nm}(t)=\phi (t)\delta _{nm}=\frac{1}{N}\underset{𝐪}{}\phi (t)\mathrm{exp}\left\{i𝐪(𝐧𝐦)\right\}$$ (53) This approximation is related to the fact that all correlation functions between different lattice cells $`\mathrm{\Phi }_{nm}(t)`$ with $`nm`$ (see eq.21) vanish for $`t=0`$ ($`\mathrm{\Phi }_{nm}(0)`$ is the equilibrium average $`\widehat{\eta }_n\widehat{\eta }_m`$ which decouples as a result of the simple Hamiltonian $`\widehat{H}_i\widehat{\eta }_i(t)`$, i.e. $`\widehat{\eta }_n\widehat{\eta }_m=\widehat{\eta }_n\widehat{\eta }_m=0`$) and for $`t\mathrm{}`$ (because of the ergodicity follows again $`\mathrm{\Phi }_{nm}(\mathrm{})=\widehat{\eta }_n(\mathrm{})\widehat{\eta }_m(0)=\widehat{\eta }_n(\mathrm{})\widehat{\eta }_m(0)=0`$). Furthermore, one obtains only a very small correlation $`\mathrm{\Phi }_{nm}(t)`$ between different cells for finite times $`t`$ differences which can be checked by numerical investigation. From (38) and (32) it follows for $`q0`$ the relation $`g_0(0)=\mathrm{\Gamma }_0^1/\tau _0`$. Thus, the evolution equation on a macroscopic scale is given by (Note, that there is the initial condition: $`\phi (0)=1`$): $$\phi (z)=\left[z+\frac{N\mathrm{\Omega }^{(10)}(0)NK^{(10)}(0,z)+z\mathrm{\Gamma }_0^1/\tau _0}{z+N\mathrm{\Omega }^{(11)}(0)NK^{(11)}(0,z)\mathrm{\Gamma }_0^1/\tau _0}\right]^1$$ (54) The determination of the frequency matrices was realized above. One obtains in the macroscopic limit by using (35) and (36): $$N\mathrm{\Omega }^{(10)}(0)=\frac{1}{\tau _0^2}\frac{\mathrm{\Gamma }_0^1\mathrm{\Gamma }_0^3\left(\mathrm{\Gamma }_0^2\right)^2}{\mathrm{\Gamma }_0^2(\mathrm{\Gamma }_0^1)^2}\mathrm{and}N\mathrm{\Omega }^{(11)}(0)=\frac{1}{\tau _0}\frac{\mathrm{\Gamma }_0^3\mathrm{\Gamma }_0^1\mathrm{\Gamma }_0^2}{\mathrm{\Gamma }_0^2(\mathrm{\Gamma }_0^1)^2}$$ (55) It remains the determination of the memory. Using the definition: $`\begin{array}{ccc}h^\alpha (𝐪,𝐪^{})\hfill & =\hfill & \frac{1}{N}_{j,k,l}H_{j,(kl)}^{\alpha (2)}\mathrm{exp}\left\{i𝐪(𝐤𝐣)+i𝐪^{}(𝐥𝐣)\right\}\hfill \\ \stackrel{~}{h}^\alpha (𝐪,𝐪^{})\hfill & =\hfill & \frac{1}{N}_{j,k,l}\stackrel{~}{H}_{(kl),j}^{(2)\alpha }\mathrm{exp}\left\{i𝐪(𝐤𝐣)+i𝐪^{}(𝐥𝐣)\right\}\hfill \end{array}`$ and (53), it follows from (52): $`\begin{array}{ccc}NK^{(\gamma \beta )}(0,t)\hfill & =\hfill & \frac{1}{N}\underset{\alpha }{}_𝐪h^\gamma (𝐪,𝐪)\phi (t)^2\stackrel{~}{h}^\alpha (𝐪,𝐪)(Ng^{\alpha \beta }(0))\hfill \\ & \hfill & \underset{\alpha }{}h^\gamma (0,0)\phi (t)^2\stackrel{~}{h}^\alpha (0,0)Ng^{\alpha \beta }(0)\hfill \end{array}`$ The substitution $`h^\gamma (𝐪,𝐪)h^\gamma (0,0)`$ and $`\stackrel{~}{h}^\alpha (𝐪,𝐪)\stackrel{~}{h}^\alpha (0,0)`$ is possible because the $`q`$–dependence of these quantities is again relatively weak. We need $`h^1(0,0)`$, $`\stackrel{~}{h}^0(0,0)`$ and $`\stackrel{~}{h}^1(0,0)`$ for the following investigations. A simple calculation leads to $`\stackrel{~}{h}^0(0,0)=0`$. Using (34) the memory terms can be written as: $`\begin{array}{ccc}NK^{(10)}(0,t)\hfill & =\hfill & \tau _0\phi (t)^2h^1(0,0)\stackrel{~}{h}^1(0,0)\frac{\mathrm{\Gamma }_0^1}{\mathrm{\Gamma }_0^2(\mathrm{\Gamma }_0^1)^2}\hfill \\ NK^{(11)}(0,t)\hfill & =\hfill & \tau _0^2\phi (t)^2h^1(0,0)\stackrel{~}{h}^1(0,0)\frac{1}{\mathrm{\Gamma }_0^2(\mathrm{\Gamma }_0^1)^2}\hfill \end{array}`$ It should be remarked that the ratio between both memory terms is given by the relation: $`{\displaystyle \frac{K^{(10)}(0,t)}{K^{(11)}(0,t)}}={\displaystyle \frac{\mathrm{\Gamma }_0^1}{\tau _0}}`$ Finally, we must determine the quantity $`\lambda =\tau _0^2h^1(0,0)\stackrel{~}{h}^1(0,0)/(4(\mathrm{\Gamma }_0^2(\mathrm{\Gamma }_0^1)^2))`$. This is again an algebraic procedure which can be realized straightforwardly. The final results are very unwieldy. Therefore, we give only the explicit expressions for the asymptotic case of low temperatures, i.e. for $`\overline{\sigma }_{\mathrm{eq}}0`$. One obtains for a square lattice ($`z_c=4`$) and a cubic lattice ($`z_c=6`$), respectively: $`\begin{array}{ccccc}\lambda & =& \frac{32\overline{\sigma }_{\mathrm{eq}}^3}{3}(1+o(\overline{\sigma }_{\mathrm{eq}}))& \mathrm{for}& z_c=4\\ \lambda & =& \frac{64\overline{\sigma }_{\mathrm{eq}}^3}{5}(1+o(\overline{\sigma }_{\mathrm{eq}}))& \mathrm{for}& z_c=6\end{array}`$ The behavior $`\lambda \overline{\sigma }_{\mathrm{eq}}^3`$ is characteristic for $`T0`$ and $`\overline{\sigma }_{\mathrm{eq}}0`$, respectively. (54) can now be written as: $$\phi (z)=\left[z+\frac{\mathrm{\Gamma }_0^1}{\tau _0}\frac{1}{\tau _0}\frac{\mathrm{\Gamma }_0^1N\mathrm{\Omega }^{(11)}(0)\tau _0N\mathrm{\Omega }^{(10)}(0)\tau _0^2\left(\mathrm{\Gamma }_0^1\right)^2}{z\tau _0+N\mathrm{\Omega }^{(11)}(0)\tau _0\lambda \mathrm{\Xi }(z)\mathrm{\Gamma }_0^1}\right]^1$$ (56) with $`\mathrm{\Xi }(z)=_0^{\mathrm{}}\left(dt/\tau _0\right)\phi (t)^2\mathrm{exp}\left\{zt\right\}`$. ## VII Discussion Now we are able to analyze the characteristic slowing down of the dynamics of the SFM$`[2,d]`$ for decreasing temperature. The central equation for the following discussion is (56). The first question is the existence of ergodicity and nonergodicity. Exists there a critical temperature $`T^{}`$, so that the correlation function $`\phi (t)`$ shows an incomplete decay $`\phi (t\mathrm{})=f_{\mathrm{}}0`$ for $`TT^{}`$? In other words, has the function $`\phi (z)`$ a pole at $`z=0`$ below the critical temperature $`T^{}`$? This question is equivalent to the determination of a kinetic phase transition from an ergodic state into a nonergodic state for supercooled liquids . To this aim we split the correlation function into a nonergodicity part $`f_{\mathrm{}}`$ and a contribution $`\phi _{\mathrm{erg}}(t)`$: $$\phi (t)=f_{\mathrm{}}+\phi _{\mathrm{erg}}(t)$$ (57) The function $`\phi _{\mathrm{erg}}(t)`$ describes the remaining ergodic part of the SFM$`[2,d]`$, i.e. $`\phi _{\mathrm{erg}}(t\mathrm{})=0`$. The Laplace transformation leads to $$\phi (z)=\frac{f_{\mathrm{}}}{z}+\phi _{\mathrm{erg}}(z)$$ (58) with $`\underset{z0}{lim}z\phi _{\mathrm{erg}}(z)=0`$. The memory term $`\mathrm{\Xi }(z)`$ can expressed by: $$\mathrm{\Xi }(z)=\frac{f_{\mathrm{}}^2}{z\tau _0}+\mathrm{\Xi }_{\mathrm{erg}}(z)$$ (59) with $`\underset{z0}{lim}z\mathrm{\Xi }_{\mathrm{erg}}(z)=0`$. (58), (59) and (56) yield: $$\underset{z0}{lim}z\phi (z)=f_{\mathrm{}}=\underset{z0}{lim}\left[1+\frac{\mathrm{\Gamma }_0^1}{z\tau _0}+\frac{\mathrm{\Gamma }_0^1N\mathrm{\Omega }^{(11)}(0)\tau _0N\mathrm{\Omega }^{(10)}(0)\tau _0^2\left(\mathrm{\Gamma }_0^1\right)^2}{\lambda f_{\mathrm{}}^2}\right]^1$$ (60) It follows immediately that the nonergodicity part $`f_{\mathrm{}}`$ has a nonvanishing value only if $`\mathrm{\Gamma }_0^1=0`$ . Otherwise, the only solution of (60) is given by $`f_{\mathrm{}}=0`$, i.e. the SFM$`[2,d]`$ is an ergodic system if $`\mathrm{\Gamma }_0^10`$. The value of $`\mathrm{\Gamma }_0^1`$ vanishes only for $`T=0`$: $`\mathrm{\Gamma }_0^1=0T=0`$ i.e. the SFM$`[2,d]`$ realizes a kinetic phase transition from an ergodic system to a nonergodic system at the critical temperature $`T^{}=0`$. Thus, the only nonergodic state can be observed at zero–point temperature. Additionally, we obtain from (55): $`\mathrm{\Omega }^{(11)}(0)=0\text{ and }\mathrm{\Omega }^{(10)}(0)=0\mathrm{for}T=0`$ Thus, the nonergodic part is given by: $`f_{\mathrm{}}=1`$ i.e. an initial equilibrium configuration at $`T=0`$ shows no structure relaxations during the total observation time. This behavior is a consequence of the fact that the kinetic phase transition occurs at the absolute zero at temperature. In other words, each arbitrary equilibrium configuration is frozen at $`T=T^{}=0`$. We obtain the important result that the SFM$`[2,d]`$ is ergodic for all finite temperatures $`T>0`$. This is a real contradiction to the statements of the original papers which predict a kinetic phase transition at a finite critical temperature. Of course, it is possible to explain $`T^{}=0`$ by a simple picture. A finite fraction of active cells (i.e. cells which are able to change their state without any previous change of states of neighbored cells) exists at each temperature $`T>0`$. The concentration of these active or nonfrozen cells is proportional to $`\overline{\sigma }_{\mathrm{eq}}^2`$ ($`\overline{\sigma }_{\mathrm{eq}}^2`$ is equivalent to the probability that a cell has two neighbored cells with $`\sigma =1`$). Active cells are mainly isolated at sufficiently low temperatures, but an annihilation of an isolated active cell is never possible. The property to be an active cell can be transferred to a neighbored cell by some few elementary flip processes, (this procedure can be interpreted as a diffusion of active states), a new active cell can be created in the nearest environment of an initially active cell (creation process) and two neighbored active cells can be unified to one active region and further this region can be reduced to one active cell (annihilation process). Diffusion, annihilation and creation are well–balanced in the thermodynamical equilibrium. On the other hand, this three processes realize a motion of active states through the whole volume, i.e. each cell is able to change their state after a sufficiently long waiting time. Note that this considerations must be modified for a finite volume because some configurations of the SFM$`[2,d]`$ show a self blockade in finite geometries . But it should be denoted also, that a special choice of the initial configuration excludes any type of self blockades . However, the main result of our discussion is the ergodicity of the SFM$`[2,d]`$ for $`T>0`$. To analyze the relaxation behavior at finite temperatures $`T>0`$ we introduce the relaxation time $`\tau _c`$ $$\tau _c=\tau _0\frac{N\mathrm{\Omega }^{(11)}(0)\tau _0\lambda \mathrm{\Xi }(0)\mathrm{\Gamma }_0^1}{N\mathrm{\Omega }^{(10)}(0)\tau _0^2\lambda \mathrm{\Gamma }_0^1\mathrm{\Xi }(0)}$$ (61) and the coefficient $$\varsigma =N\mathrm{\Omega }^{(11)}(0)\tau _0\lambda \mathrm{\Xi }(0)\mathrm{\Gamma }_0^1$$ (62) Using these notations, we get from (56): $$\phi (z)=\left[z+\tau _c^1+\left(\frac{\mathrm{\Gamma }_0^1}{\tau _0}\tau _c^1\right)\left[1+\frac{\varsigma }{z\tau _0+\lambda \left[\mathrm{\Xi }(0)\mathrm{\Xi }(z)\right]}\right]^1\right]^1$$ (63) (61), (62) and (63) are a closed, nonlinear system of equations, which can be solved by numerical standard methods. Fig.4 shows $`\phi (t)`$ for a SFM$`[2,3]`$ on a square lattice ($`z_c=6`$) for various relative temperatures $`T/\epsilon `$. We see that the correlation function $`\phi (t)`$ shows with decreasing temperature a pronounced stretched decay over some decades, while an exponential–like decay is obtained for high temperatures. This stretching can be illustrated by a simple argument. Short times ($`t0`$ or $`z\mathrm{}`$) are related to a behavior $`\phi (z)(z+\mathrm{\Gamma }_0^1/\tau _0)^1`$ or $`\phi (t)\mathrm{exp}\left\{\mathrm{\Gamma }_0^1t/\tau _0\right\}`$. On the other hand, the long time regime ($`t\mathrm{}`$ or $`z0`$) is characterized by $`\phi (z)(z+\tau _c^1)^1`$ or $`\phi (t)\mathrm{exp}\left\{t/\tau _c\right\}`$. As there is $`\mathrm{\Gamma }_0^1/\tau _0\tau _c^1`$ we expect a typical crossover between both regimes characterized by a stretched decay, see fig.4. But it should be remarked, that the decay of $`\phi (t)`$ is no Kohlrausch–Williams–Watts function ($`\mathrm{exp}((t/\tau )^\gamma )`$. There exists only over a finite interval a reasonable fit with such a stretched exponential function . The spectral density $`S(\lambda _s)`$ of the correlation function is defined as the set of amplitudes of exponential decays which contribute to $`\phi (t)`$. Thus, $`\phi (t)`$ is the Laplace transformed spectral density with respect to the Laplace variable $`t`$: $$\phi (t)=\underset{0}{\overset{\mathrm{}}{}}𝑑\lambda _sS(\lambda _s)\mathrm{exp}(\lambda _st)$$ (64) One obtains the remarkable result that the spectral density is positive definite, (fig.5). The knowledge of the spectral densities allows the determination of other interesting properties, for example the susceptibility $`\chi (\omega )`$, see fig. 6. Finally, the averaged relaxation time $`\tau (T)`$ can be obtained by using $$\tau (T)=\frac{\underset{0}{\overset{\mathrm{}}{}}t\phi (t)𝑑t}{\underset{0}{\overset{\mathrm{}}{}}\phi (t)𝑑t}=\frac{\underset{0}{\overset{\mathrm{}}{}}S(\lambda _s)\lambda _s^2𝑑\lambda _s}{\underset{0}{\overset{\mathrm{}}{}}S(\lambda _s)\lambda _s^1𝑑\lambda _s}$$ (65) This relaxation time shows with decreasing temperature an increasing deviation from the simple relaxation times (41). Especially, one obtains now a typical non–Arrhenius behavior (fig.7). ## VIII Conclusions It was shown that irreversible master equations can be easily transformed into projection equations by using the Fock space representation. Whereas the usual projection formalisms, which leads to the well known Mori–Zwanzig equations, start from a reversible Liouville equation, the master equations are already irreversible. As a consequence of this initial irreversibility, one obtains additional frequency matrices and memory terms. Thus, these additional contributions are caused mainly by the loss of the invariance against an inversion of the time. A second important property follows from a general analysis of the frequency matrices. The corresponding poles of the Laplace transformed correlation function $`\stackrel{~}{\mathrm{\Phi }}(q,z)`$ are always located on the negative real axis. Especially, the imaginary part of the poles vanishes identically. This behavior is related to the general structure of the master equation. The dynamical matrix $`L(\sigma ,\sigma ^{})`$ of a master equation is always negative definite (or better semidefinite because at least one eigen value is zero as result of the conservation of the probability). Thus, only relaxation processes should be observed, i.e. the evolution of the probability $`P(\sigma ,t)`$ can be approached by a probably infinitely large expansion in terms of exponential functions: $`P(\sigma ,t)=P_{\mathrm{eq}}(\sigma ,t)+_mA_m\left(\sigma \right)\mathrm{exp}\left\{\mathrm{\Lambda }_mt\right\}`$ with ($`\mathrm{\Lambda }_m>0`$ for all $`m`$). One obtains no oscillations in contradiction to microscopical systems which bases on a Liouville equation. The absent of oscillations is directly connected with the Markov property of the underlying master equation. This behavior corresponds also to the fact that the observation of a spin wave propagation (corresponding to density waves in a glass or a supercooled liquid) is not possible for the SFM$`[n,d]`$. From this point of view, the traditional notation ’frequency matrix’ can be misleading. It seems to be possibly favorable to use the notion relaxation matrix. However, we have used the traditional terminology to avoid misunderstandings and conflicts with other well defined quantities. In principle, the consideration of all possible time derivations $`\widehat{\eta }_i^{(\beta )}(t)=^\beta \widehat{\eta }_i(t)/t^\beta =\widehat{\eta }_i(t)\widehat{L}^\beta `$ ($`\beta =0,1,\mathrm{},\mathrm{}`$) as relevant operators leads to an infinite continuous fraction for the correlation function $`\mathrm{\Phi }_{nm}(t)`$, determined by frequency matrices of various order. This representation allows a systematic analysis of the short time behavior of the SFM$`[2,d]`$, because an infinite continuous fraction contains no memory term. Unfortunately, a general explicit determination of the frequency matrices, e.g. by a successive rule, cannot be obtained. On the other hand, each approximation using a finite number of frequency matrices (e.g. $`\mathrm{\Omega }^{(10)}`$ and $`\mathrm{\Omega }^{(11)}`$ in the present case) fails for sufficiently long times. However, some important properties of the spin facilitated kinetic Ising model can be verified for this relatively rough approximation, e.g. a weak non–Arrhenius behavior of the relaxation times. Note, that this result is in agreement with various numerical simulations . On the other hand, the typical stretched decay of the correlation function can not be explained by using this simple approximation. A satisfactory treatment is possible by an approximative consideration of the memory terms. This procedure leads to an equation, which is partially similar to the well known mode coupling equation of supercooled liquids. The memory terms yield the main contributions to the typical stretching of the correlation function. Some small remaining deviations from the numerical results are caused by the approximations (50) and (51) of the memory terms. From this point of view, we come back to the initial question. Which processes of the usual glass transition can be described by the spin facilitated kinetic Ising models? Obviously, the spin facilitated kinetic Ising model is not adequate for a description of the fast processes inside a supercooled liquid. This statement is supported by both, numerical simulations and the presented analytical investigations, which show that no fast ($`\beta `$–) processes can be observed. On the other hand, spin facilitated kinetic Ising models allow a more or less reasonable, quantitative description of the slow ($`\alpha `$–) process in supercooled liquids below the critical temperature $`T_c`$ of the usual mode coupling theory . Note that for this low temperature regime the time scales of $`\alpha `$– and $`\beta `$–process are well separated, i.e. a separation of the underlying dynamic is actually possible and it is considered in the structure of the master equations related to the SFM$`[n,d]`$, see above. The fast dynamics, corresponding to the $`\beta `$–process, determines the thermodynamical noise. This noise is not explicitly contained in the previous equations, but it is the underlying cause for the irreversibility of the master equations. Additionally, master equations and stochastic evolution equations are principally equivalent . However, the remaining slow dynamics of a supercooled liquid ($`\alpha `$–process) is represented in the kinetic scenario of the SFM$`[n,d]`$. ACKNOWLEDGEMENTS This work has been supported by the Deutsche Forschungsgemeinschaft DFG (SFB 418 and schu 934/3-1). ## A Representation of $`\mathrm{\Gamma }`$ The rigorous values $`A^\alpha `$, $`B^\alpha `$, $`C^\alpha `$ and $`D^\alpha `$ of (33) can be obtained straightforwardly by using the evolution operator $`\widehat{L}`$ (6), the commutation relations (3) and the definition of $`\mathrm{\Gamma }_{ik}^\beta `$ (29). Under consideration of the coordination number $`z_c`$ and the thermodynamical equilibrium $`\overline{\sigma }_{\mathrm{eq}}`$ of the cell state we obtain the following results $`A^\alpha `$-terms: $$\begin{array}{ccc}A^0\hfill & =\hfill & 1\hfill \\ A^1\hfill & =\hfill & 2\overline{\sigma }_{\mathrm{eq}}^2\left(\begin{array}{c}z_c\\ 2\end{array}\right)\hfill \\ A^2\hfill & =\hfill & 4\overline{\sigma }_{\mathrm{eq}}^2\left[\left(\begin{array}{c}z_c\\ 2\end{array}\right)+6\overline{\sigma }_{\mathrm{eq}}\left(\begin{array}{c}z_c\\ 3\end{array}\right)+6\overline{\sigma }_{\mathrm{eq}}^2\left(\begin{array}{c}z_c\\ 4\end{array}\right)\right]\hfill \\ A^3\hfill & =\hfill & 8\overline{\sigma }_{\mathrm{eq}}^2\left[\left(\begin{array}{c}z_c\\ 2\end{array}\right)+24\overline{\sigma }_{\mathrm{eq}}\left(\begin{array}{c}z_c\\ 3\end{array}\right)+114\overline{\sigma }_{\mathrm{eq}}^2\left(\begin{array}{c}z_c\\ 4\end{array}\right)+180\overline{\sigma }_{\mathrm{eq}}^3\left(\begin{array}{c}z_c\\ 5\end{array}\right)+90\overline{\sigma }_{\mathrm{eq}}^4\left(\begin{array}{c}z_c\\ 6\end{array}\right)\right]\hfill \\ & +\hfill & 4\overline{\sigma }_{\mathrm{eq}}^3\left(1\overline{\sigma }_{\mathrm{eq}}\right)z_c(z_c1)^2\left(1+\overline{\sigma }_{\mathrm{eq}}(z_c2)\right)\left(2+\overline{\sigma }_{\mathrm{eq}}(z_c4)\right)\hfill \end{array}$$ $`B^\alpha `$-terms: $$\begin{array}{ccc}B^0\hfill & =\hfill & 0\hfill \\ B^1\hfill & =\hfill & 0\hfill \\ B^2\hfill & =\hfill & 4\overline{\sigma }_{\mathrm{eq}}^3(1\overline{\sigma }_{\mathrm{eq}})(z_c1)^2\hfill \\ B^3\hfill & =\hfill & 16\overline{\sigma }_{\mathrm{eq}}^3(1\overline{\sigma }_{\mathrm{eq}})z_c(z_c1)^2\left(1+3\overline{\sigma }_{\mathrm{eq}}(z_c2)+\overline{\sigma }_{\mathrm{eq}}^2(z_c2)(z_c3)\right)\hfill \end{array}$$ $`C^\alpha `$-terms: $$C^0=C^1=C^2=0$$ and $$C^3=8\overline{\sigma }_{\mathrm{eq}}^3(1\overline{\sigma }_{\mathrm{eq}})^2\left(2+2\overline{\sigma }_{\mathrm{eq}}z_c(z_c2)\right)$$ $`D^\alpha `$-terms: $$D^0=D^1=D^2=0$$ and $$D^3=8\overline{\sigma }_{\mathrm{eq}}^4(1\overline{\sigma }_{\mathrm{eq}})^2(z_c1)^2$$ ## B Completeness of the basis $`\stackrel{~}{B}`$ The proof consists in two parts. At first we analyze an operator $`\widehat{X}^{}`$, which is a multilinear form of the operators $`\widehat{\eta }_i`$ , i.e. $$\widehat{X}^{}=\underset{n}{}\underset{𝐍_n}{}\beta _{𝐍_n}^{(n)}\widehat{B}_{𝐍_n}^{(n)}$$ (B1) with arbitrary coefficients $`\beta _{𝐍_n}^{(n)}`$. The orthogonality of the basis $`\stackrel{~}{B}=\left\{\widehat{B}_{𝐍_n}^{(n)}\right\}`$ leads immediately to $`\beta _{𝐍_n}^{(n)}=\widehat{X}^{}\widehat{B}_{𝐍_n}^{(n)}`$, i.e. all coefficients of $`\widehat{X}^{}`$ can be determined by a successive procedure. In other words, the basis $`\left\{\widehat{B}_{𝐍_n}^{(n)}\right\}`$ is complete with respect all operators $`\widehat{X}^{}`$, i.e. the basis forms a space containing all operators of type (B1). The second part of the proof analyses arbitrary operators $`\widehat{X}`$ acting on the Fock–space. $`\widehat{X}`$ consists in operators $`\widehat{\eta }_k`$ as well as annihilation operators $`\widehat{d}_k`$ and creation operators $`\widehat{d}_k^{}`$. A representation like (B1), extended by the operators $`\widehat{d}_k`$ and $`\widehat{d}_k^{}`$ is always possible: $$\widehat{X}=\underset{𝐌^1,𝐌^2,𝐌^3,𝐌^4}{}\theta _{𝐌^1,𝐌^2,𝐌^3,𝐌^4}\underset{k=1}{\overset{N}{}}\left[\widehat{\eta }_k^{m_k^1}\widehat{d}_i^{m_k^2}\left(\widehat{d}_i^{}\right)^{m_k^3}\widehat{1}^{m_k^4}\right]$$ (B2) ($`\widehat{1}`$ is the simple unit operator $`\widehat{1}1`$, $`N`$ is the number of lattice cells). The vectors $`𝐌^\gamma `$ ($`\gamma =1,\mathrm{}4`$) contains $`N`$ integer numbers, i.e. $`𝐌^\gamma =\{m_1^\gamma ,\mathrm{},m_N^\gamma \}`$ with $`m_k^\gamma 0`$ for all possible $`k`$ and $`\gamma `$. The commutation relations (3) can be used to write $`\widehat{X}`$ as a representation with the internal restrictions: $`m_k^1+m_k^2+m_k^3+m_k^4=1`$ for all lattice cells $`k`$, i.e. each contribution to the sum (B2) contains exactly one of the four operators $`\widehat{1}`$, $`\widehat{d}_k`$, $`\widehat{d}_k^{}`$ and $`\widehat{\eta }_k`$ with respect to any cell $`k`$. Furthermore, the commutation relation (3) allows a shift of all operators related to a given lattice cell $`i`$ to the left hand side. It is simple to show that : $$\begin{array}{cc}s\|1=s\|1\hfill & s\|\widehat{d}_i=s\|\left(1\overline{\sigma }_{\mathrm{eq}}\widehat{\eta }_i\sqrt{\overline{\sigma }_{\mathrm{eq}}(1\overline{\sigma }_{\mathrm{eq}})}\right)\hfill \\ s\|\widehat{\eta }_i=s\|\widehat{\eta }_i\hfill & s\|\widehat{d}_j^{}=s\|\left(\overline{\sigma }_{\mathrm{eq}}+\widehat{\eta }_i\sqrt{\overline{\sigma }_{\mathrm{eq}}(1\overline{\sigma }_{\mathrm{eq}})}\right)\hfill \end{array}$$ Obviously, the application of the operator $`\widehat{X}`$ onto the reference state $`s\|`$ is equivalent to the application of a corresponding reduced operator $`\widehat{X}^{}`$ containing only operators $`\widehat{\eta }_k`$ (and the trivial operators $`\widehat{1}=1`$) onto the state $`s\|`$, i.e. there is a definitely mapping $$\widehat{X}\widehat{X}^{}\mathrm{with}s\|\widehat{X}=s\|\widehat{X}^{}$$ Therefore, one obtains (see eq.4): $$\widehat{X}\widehat{B}_{𝐍_n}^{(n)}=s\left\|\widehat{X}\widehat{B}_{𝐍_n}^{(n)}\right\|F=s\left\|\widehat{X}^{}\widehat{B}_{𝐍_n}^{(n)}\right\|F=\widehat{X}^{}\widehat{B}_{𝐍_n}^{(n)}$$ and consequently by using (B1): $$s\|\widehat{X}=s\|\widehat{X}^{}=\underset{n}{}\underset{𝐍_n}{}\widehat{X}^{}\widehat{B}_{𝐍_n}^{(n)}s\|\widehat{B}_{𝐍_n}^{(n)}=\underset{n}{}\underset{𝐍_n}{}\widehat{X}\widehat{B}_{𝐍_n}^{(n)}s\|\widehat{B}_{𝐍_n}^{(n)}$$ i.e. an arbitrary operator $`\widehat{X}`$ of the Fock–space can be completely presented by using the basis $`\stackrel{~}{B}=\left\{\widehat{B}_{𝐍_n}^{(n)}\right\}`$ under the consideration that this operator acts into the left direction on the reference state $`s\|`$.
warning/0004/hep-th0004192.html	ar5iv	text	# Transverse Spin in QCD. I. Canonical Structure ## I Introduction From the early days of quantum field theory, it has been recognized that the issues associated with the spin of a composite system in an arbitrary reference frame are highly complex and non-trivial. The familiar Pauli-Lubanski operators readily qualify for spin operators only in the rest frame of the particle. For a single particle in a moving frame it is known how to construct the appropriate spin operators starting from the Pauli-Lubanski operators. How to construct the spin operators for a composite system in an arbitrary reference frame is a nontrivial problem. In equal-time quantization, the complexities arise from the facts that for a moving composite object, Pauli-Lubanski operators are necessarily interaction dependent and, further, it is quite difficult to separate the center of mass and internal variables which is mandatory in the calculation of spin. Due to these difficulties there has been rarely any attempt to study the spin of a moving composite system in the conventional equal time formulation of even simple field theoretic models, let alone Quantum Chromo Dynamics (QCD). From the early days of light-front field theory, the complications associated with transverse rotation operators $`F^i`$ have been recognized. They are interaction dependent just like the Hamiltonian. Furthermore, together with the third component of the rotation operator $`J^3`$, which is kinematical, $`F^i`$ do not obey the angular momentum algebra. Instead they obey the algebra of two dimensional Euclidean group which is appropriate only for massless particles. For massive particles, one can define transverse spin operators which together with the third component (helicity) obey the angular momentum algebra. However, they cannot be separated into orbital and spin parts unlike the helicity operator. Most of the studies of the transverse spin operators in light-front field theory, so far, are restricted to free field theory. Even in this case the operators have a complicated structure. However, one can write these operators as a sum of orbital and spin parts, which can be achieved via a unitary transformation, the famous Melosh transformation. In interacting theory, presumably this can be achieved order by order in a suitable expansion parameter which is justifiable only in a weakly coupled theory. Knowledge about transverse rotation operators and transverse spin operators is mandatory for addressing issues concerning Lorentz invariance in light-front theory. Unfortunately, very little is known regarding the field theoretic aspects of the interaction dependent spin operators, We emphasize that in a moving frame, the spin operators are interaction dependent irrespective of whether one considers equal-time field theory or light-front field theory. To the best of our knowledge, in gauge field theory, the canonical structure of spin operators of a composite system in a moving frame has never been studied. In this work we initiate a systematic investigation of the spin of a composite system in a moving frame in QCD. A brief summary of some of our results has been presented in Ref. . We show that, in spite of the complexities, light-front field theory offers a unique opportunity to address the issue of relativistic spin operators in an arbitrary reference frame since boost is kinematical in this formulation. The plan of this paper is as follows. In Sec. II, first, we briefly review the complexities associated with the description of the spin of a composite system in a moving frame in the conventional equal time quantization. Then we give the canonical structure of light-front Lorentz algebra and light-front spin operators. In this section we also provide a detailed discussion of the transverse spin operators for a massless particle of arbitrary transverse momentum. The explicit form of transverse spin operators in light-front QCD is derived in Sec. III. Summary and conclusions are presented in Sec. IV. For the sake of completeness and clarity, in Appendix A we review the intrinsic spin operators in relativistic quantum mechanics. The explicit form of the kinematical operators and the Hamiltonian in light-front QCD starting from the gauge invariant, symmetric, interaction dependent, energy momentum tensor is derived in Appendix B. A complete discussion of transverse spin operators in free fermion field theory and free massless, spin one boson field theory is presented in detail in Appendices C and D. ## II Preliminaries In this section, first we briefly review the intrinsic spin operator in equal-time quantization. We highlight the difficulties one encounters in constructing the spin operator of a composite system in an arbitrary reference frame in this case. Next, we give the Lorentz generators in light-front formulation and show that with the help of the kinematical boost in the light-front formalism, a relativistic spin operator for a composite system can be defined in an arbitrary reference frame for massive as well as massless particles. We also compare and contrast the spin operators in equal-time and light-front quantization. ### A Intrinsic Spin in Equal Time Quantization Intrinsic spin operators in an arbitrary reference frame in equal-time quantization are given in terms of the Poincare generators by (see Appendix A for details) $`𝐒=`$ $`{\displaystyle \frac{1}{M}}\left[𝐖{\displaystyle \frac{𝐏W^0}{M+H}}\right]`$ (2) $`=𝐉{\displaystyle \frac{P^0}{M}}𝐊\times {\displaystyle \frac{𝐏}{M}}{\displaystyle \frac{(𝐉𝐏)}{M+P^0}}{\displaystyle \frac{𝐏}{M}}`$ where $`𝐖`$ are the space components of the Pauli-Lubanski operator, $`W^\mu =\frac{1}{2}ϵ^{\mu \nu \rho \lambda }M_{\nu \rho }P_\lambda `$. $`H`$,$`\stackrel{}{P}`$ are equal time Hamiltonian and momentum operators respectively obtained by integrating the energy momentum tensor over a spacelike surface and $`\stackrel{}{J}`$ and $`\stackrel{}{K}`$ are the equal time rotation and boost generators respectively, which are obtained by integrating the angular momentum density over a spacelike surface. Since boost $`𝐊`$ is dynamical, all the three components of $`𝐒`$ are interaction dependent in the equal time quantization. Nevertheless, the component of S along P remains kinematical. This is to be compared with light-front quantization where the third component of the light-front spin operator $`𝒥^3`$ is kinematical (see Sec. IIB). This arises from the facts that boost operators are kinematical on the light front, the interaction dependence of light-front spin operators $`𝒥^i`$ arises solely from the rotation operators, and the third component of the rotation operator $`J^3`$ is kinematical on the light front. A further essential complication arises in equal time quantization. In order to describe the intrinsic spin of a composite system, one should be able to separate the center of mass motion from the internal motion. Even in free field theory, this turns out to be quite involved (See Ref. and references therein). On the other hand, in light-front theory, since transverse boosts are simply Galilean boosts, separation of center of mass motion and internal motion is as simple as in non-relativistic theory. (See Appendix A, of Ref. for a detailed example). ### B Light-Front Lorentz Generators and Algebra In terms of the gauge invariant, symmetric energy momentum tensor $`\mathrm{\Theta }^{\mu \nu }`$, the four-vector $`P^\mu `$ and the tensor $`M^{\mu \nu }`$ are given by $`P^\mu `$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\mathrm{\Theta }^{+\mu }},`$ (3) $`M^{\mu \nu }`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left[x^\mu \mathrm{\Theta }^{+\nu }x^\nu \mathrm{\Theta }^{+\mu }\right]}.`$ (4) The boost operators are $`M^+=2K^3`$ and $`M^{+i}=E^i`$. The rotation operators are $`M^{12}=J^3`$ and $`M^i=F^i`$. The Hamiltonian $`P^{}`$ and the transverse rotation operators $`F^i`$ are dynamical (depend on the interaction) while other seven operators are kinematical. The rotation operators obey the $`E(2)`$-like algebra of two dimensional Euclidean group, namely, $`[F^1,F^2]=0,[J^3,F^i]=iϵ^{ij}F^j`$ (5) where $`ϵ^{ij}`$ is the two-dimensional antisymmetric tensor. Thus $`F^i`$ do not qualify as angular momentum operators. Moreover, $`F^i`$ are not translationally invariant and hence they do not qualify as intrinsic spin. ### C Transverse Spin Operators: Massive particle The Pauli-Lubanski spin operator $`W^\mu ={\displaystyle \frac{1}{2}}ϵ^{\mu \nu \rho \sigma }M_{\nu \rho }P_\sigma `$ (6) with $`ϵ^{+12}=2`$. For a massive particle, the transverse spin operators $`𝒥^i`$ in light-front theory are given in terms of Poincare generators by $`M𝒥^1`$ $`=W^1P^1𝒥^3={\displaystyle \frac{1}{2}}F^2P^++K^3P^2{\displaystyle \frac{1}{2}}E^2P^{}P^1𝒥^3,`$ (7) $`M𝒥^2`$ $`=W^2P^2𝒥^3={\displaystyle \frac{1}{2}}F^1P^+K^3P^1+{\displaystyle \frac{1}{2}}E^1P^{}P^2𝒥^3.`$ (8) The first term in Eqs. (7) and (8) contains both center of mass motion and internal motion and the next three terms in these equations serve to remove the center of mass motion. The helicity operator is given by $`𝒥^3`$ $`={\displaystyle \frac{W^+}{P^+}}=J^3+{\displaystyle \frac{1}{P^+}}(E^1P^2E^2P^1).`$ (9) Here, $`J^3`$ contain both center of mass motion and internal motion and the other two terms serve to remove the center of mass motion. The operators $`𝒥^i`$ obey the angular momentum commutation relations $`[𝒥^i,𝒥^j]=iϵ^{ijk}𝒥^k.`$ (10) They commute with $`P^\mu `$. ### D Transverse Spin Operators: Massless case Again, we start from the Pauli-Lubanski spin operator, $`W^\mu ={\displaystyle \frac{1}{2}}ϵ^{\mu \nu \rho \sigma }M_{\nu \rho }p_\sigma .`$ (11) For the light-like vector $`p^\mu `$, usually the collinear choice is made, namely, $`p^+0`$, $`p^{}=0`$. Then we get, $`W^{}=0`$, $`W^+=J^3p^+`$, $`W^1=\frac{1}{2}F^2p^+`$, $`W^2=\frac{1}{2}F^1p^+`$. In free field theory, we have explicitly constructed the Poincare generators for a massless spin one particle in $`A^+=0`$ gauge in Appendix D. Consider the single particle state $`p\lambda `$ with $`p^{}=0`$. From the explicit form of the operators, we find that $`J^3p\lambda `$ $`=\lambda p\lambda ,`$ (12) $`F^ip\lambda `$ $`=0,i=1,2`$ (13) since $`p^{}=0`$. For calculations with composite states (dressed partons, for example) we have to consider light-like particles with arbitrary transverse momenta. Let us try a light like momentum $`P^\mu `$ with $`P^{}0`$, but $`P^{}=\frac{(P^{})^2}{P^+}`$ so that $`P^2=0`$. Then we get, as in the case of massive particle, $`W^+`$ $`=J^3P^++E^1P^2E^2P^1,`$ (14) $`W^1`$ $`={\displaystyle \frac{1}{2}}F^2P^++K^3P^2{\displaystyle \frac{1}{2}}E^2P^{},`$ (15) $`W^2`$ $`={\displaystyle \frac{1}{2}}F^1P^+K^3P^1+{\displaystyle \frac{1}{2}}E^1P^{},`$ (16) $`W^{}`$ $`=F^2P^1F^1P^2J^3P^{}.`$ (17) Thus even though $`W^1`$ and $`W^2`$ do not annihilate the state, we do get $`W^\mu W_\mu (=\frac{1}{2}W^+W^{}+\frac{1}{2}W^{}W^+(W^1)^2(W^2)^2)k\lambda =0`$ as it should be for a massless particle. Just as in the case of massive particle, we have the helicity operator for the massless particle, $`𝒥^3`$ $`={\displaystyle \frac{W^+}{P^+}}=J^3+{\displaystyle \frac{1}{P^+}}(E^1P^2E^2P^1).`$ (18) In analogy with the transverse spin for massive particles, we define the transverse spin operators for massless particles as $`𝒥^i=W^iP^i𝒥^3.`$ (19) They do satisfy $`𝒥^ik\lambda `$ $`=0,`$ (20) $`𝒥^3k\lambda `$ $`=\lambda k\lambda ,`$ (21) where $`k`$ is an arbitrary momentum. The operators $`𝒥^i`$ and $`𝒥^3`$ obey the $`E(2)`$-like algebra $`[𝒥^1,𝒥^2]=0,[𝒥^3,𝒥^1]=i𝒥^2,[𝒥^3,𝒥^2]=i𝒥^1.`$ (22) ### E Comments In order to calculate the transverse spin operators, first we need to construct the Poincare generators $`P^+`$, $`P^i`$, $`P^{}`$, $`K^3`$, $`E^i`$, $`J^3`$ and $`F^i`$ in light-front QCD. The explicit form of the operator $`J^3`$ is given Ref. . The construction of $`F^i`$ which is algebraically quite involved is carried out in the next section. The construction of the rest of the kinematical operators is given in Appendix B. In this appendix we also present the Hamiltonian in a manifestly Hermitian form. In order to have a physical picture of the complicated situation at hand it is instructive to calculate the spin operator in free field theory. The case of free massive fermion is carried out in Appendix C. In free field theory one can indeed show that (see Appendix C) $`𝒥^ik\lambda =\frac{1}{2}_\lambda ^{}\sigma _{\lambda ^{}\lambda }^ik\lambda ^{}`$. The case of free massless spin one particle is carried out in Appendix D. ## III The transverse rotation operator in QCD In this section we derive the expressions for interaction dependent transverse rotation operators in light-front QCD starting from the manifestly gauge invariant energy momentum tensor. It is extremely interesting to compare and contrast the situation in the equal time and light-front case. The angular momentum density $`^{\alpha \mu \nu }=x^\mu \mathrm{\Theta }^{\alpha \nu }x^\nu \mathrm{\Theta }^{\alpha \mu }.`$ (23) In equal time theory, generalized angular momentum $`M^{\mu \nu }={\displaystyle d^3x^{0\mu \nu }}.`$ (24) The rotation operators are $`J^i=ϵ^{ijk}M^{jk}`$. Thus in a non-gauge theory, all the three components of the rotation operators are manifestly interaction independent. However, the spin operators $`S^i`$ for a composite system in a moving frame involves, in addition to $`J^i`$, the boost operators $`K^i=M^{0i}`$ which are interaction dependent. Thus all the three components of $`S^i`$ become interaction dependent. A gauge invariant separation of the nucleon angular momentum is performed in Ref. . However, as far the spin operator in an arbitrary reference frame is concerned, the analysis of this reference is valid only in the rest frame where spin coincides with total angular momentum operator and in an arbitrary reference frame the need to project out the center of mass motion, which is quite complicated in equal time theory is not emphasized. Moreover, the distinction between the longitudinal and transverse components of the spin is never made. It is crucial to make this distinction since physically the longitudinal and transverse components of the spin carry quite distinct information (as is clear, for example, from the spin of a massless particle). Moreover, even for the third component of the spin of a composite system in a moving frame, there is crucial difference between equal time and light front cases. $`𝒥^3`$ (helicity) is interaction independent whereas $`S^3`$ is interaction dependent in general except when measured along the direction of P. In light-front theory, generalized angular momentum $`M^{\mu \nu }={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}^{+\mu \nu }}.`$ (25) $`J^3`$ which is related to the helicity is given by $`J^3=M^{12}={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^1d^2x^{}[x^1\mathrm{\Theta }^{+2}x^2\mathrm{\Theta }^{+1}]}`$ (26) and is interaction independent. On the other hand, the transverse rotation operators which are related to the transverse spin are given by $$F^i=M^i=\frac{1}{2}𝑑x^{}d^2x^{}[x^{}\mathrm{\Theta }^{+i}x^i\mathrm{\Theta }^+].$$ They are interaction dependent even in a non-gauge theory since $`\mathrm{\Theta }^+`$ is the Hamiltonian density. In light-front theory we set the gauge $`A^+=0`$ and eliminate the dependent variables $`\psi ^{}`$ and $`A^{}`$ using the equations of constraint. In this paper we restrict to the topologically trivial sector of the theory and set the boundary condition $`A^i(x^{},x^i)0`$ as $`x^{,i}\mathrm{}`$. This completely fixes the gauge and put all surface terms to zero. The transverse rotation operator $`F^i={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left[x^{}\mathrm{\Theta }^{+i}x^i\mathrm{\Theta }^+\right]}.`$ (27) The symmetric, gauge invariant energy momentum tensor $`\mathrm{\Theta }^{\mu \nu }`$ $`={\displaystyle \frac{1}{2}}\overline{\psi }\left[\gamma ^\mu iD^\nu +\gamma ^\nu iD^\mu \right]\psi F^{\mu \lambda a}F_\lambda ^{\nu a}g^{\mu \nu }\left[{\displaystyle \frac{1}{4}}(F_{\lambda \sigma a})^2+\overline{\psi }(\gamma ^\lambda iD_\lambda m)\psi \right],`$ (28) where $`iD^\mu `$ $`={\displaystyle \frac{1}{2}}\stackrel{}{i^\mu }+gA^\mu ,`$ (29) $`F^{\mu \lambda a}`$ $`=^\mu A^{\lambda a}^\lambda A^{\mu a}+gf^{abc}A^{\mu b}A^{\lambda c},`$ (30) $`F_\lambda ^{\nu a}`$ $`=^\nu A_\lambda ^a_\lambda A^{\nu a}+gf^{abc}A^{\nu b}A_\lambda ^c.`$ (31) First consider the fermionic part of $`\mathrm{\Theta }^{\mu \nu }`$: $`\mathrm{\Theta }_F^{\mu \nu }={\displaystyle \frac{1}{2}}\overline{\psi }\left[\gamma ^\mu iD^\nu +\gamma ^\nu iD^\mu \right]\psi g^{\mu \nu }\overline{\psi }(\gamma ^\lambda iD_\lambda m)\psi .`$ (32) The coefficient of $`g^{\mu \nu }`$ vanishes because of the equation of motion. Explicitly, the contribution to $`F^2`$ from the fermionic part of $`\mathrm{\Theta }^{\mu \nu }`$ is given by $`F_F^2`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left[x^{}\frac{1}{2}\overline{\psi }(\gamma ^+iD^2+\gamma ^2iD^+)\psi x^2\frac{1}{2}\overline{\psi }(\gamma ^+iD^{}+\gamma ^{}iD^+)\psi \right]},`$ (34) $`=F_{F(I)}^2+F_{F(II)}^2,`$ where $`F_{F(I)}^2={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}x^{}\left[\psi _{}^{+}{}_{}{}^{}\frac{1}{2}\stackrel{}{i^2}\psi ^++\psi _{}^{+}{}_{}{}^{}gA^2\psi ^++\frac{1}{4}\overline{\psi }\gamma ^i\stackrel{}{i^+}\psi \right]},`$ (35) $`F_{F(II)}^2={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}x^2\left[\psi _{}^{+}{}_{}{}^{}\left(\frac{1}{2}\stackrel{}{i^{}}+gA^{}\right)\psi ^++\frac{1}{4}\psi _{}^{}{}_{}{}^{}\gamma ^i\stackrel{}{i^+}\psi ^{}\right]}.`$ (36) We have the equation of constraint $`i^+\psi ^{}=\left[\alpha ^{}(i^{}+gA^{})+\gamma ^0m\right]\psi ^+,`$ (37) and the equation of motion $`i^{}\psi ^+=gA^{}\psi ^++\left[\alpha ^{}(i^{}+gA^{})+\gamma ^0m\right]{\displaystyle \frac{1}{i^+}}\left[\alpha ^{}(i^{}+gA^{})+\gamma ^0m\right]\psi ^+.`$ (38) Using the Eqs. (37) and (38) we arrive at free ($`g`$ independent) and interaction ($`g`$ dependent) parts of $`F_F^2`$. The free part of $`F_F^2`$ is given by $`F_{F(free)}^2`$ $`={\displaystyle \frac{1}{2}}{\displaystyle }dx^{}d^2x^{}\{x^{}[\xi ^{}\left[i^2\xi \right]\left[i^2\xi ^{}\right]\xi ]`$ (42) $`x^2\left[\xi ^{}\left[{\displaystyle \frac{(^{})^2+m^2}{i^+}}\xi \right]\left[{\displaystyle \frac{(^{})^2+m^2}{i^+}}\xi ^{}\right]\xi \right]`$ $`+\left[\xi ^{}\left[\sigma ^3^1+i^2\right]{\displaystyle \frac{1}{^+}}\xi +\left[{\displaystyle \frac{1}{^+}}(^1\xi ^{}\sigma ^3i^2\xi ^{})\right]\xi \right]`$ $`+m[\xi ^{}\left[{\displaystyle \frac{\sigma ^1}{i^+}}\xi \right]\left[{\displaystyle \frac{1}{i^+}}\xi ^{}\sigma ^1\right]\xi ]\}.`$ We have introduced the two-component field $`\xi `$, $`\psi ^+=\left[\begin{array}{c}\xi \\ 0\end{array}\right].`$ (45) The interaction dependent part of $`F_{F(I)}^2`$ is $`F_{F(I)int}^2`$ $`=g{\displaystyle 𝑑x^{}d^2x^{}x^{}\xi ^{}A^2\xi }`$ (47) $`+{\displaystyle \frac{1}{4}}g{\displaystyle 𝑑x^{}d^2x^{}\left[\xi ^{}\frac{1}{^+}[(i\sigma ^3A^1+A^2)\xi ]+\frac{1}{^+}[\xi ^{}(i\sigma ^3A^1+A^2)]\xi \right]}.`$ The interaction dependent part of $`F_{F(II)}^2`$ is $`F_{F(II)int}^2={\displaystyle \frac{1}{4}}g{\displaystyle 𝑑x^{}d^2x^{}\left[\xi ^{}\frac{1}{^+}[(i\sigma ^3A^1+A^2)\xi ]+\frac{1}{^+}[\xi ^{}(i\sigma ^3A^1+A^2)]\xi \right]}`$ (48) $`{\displaystyle \frac{1}{2}}g{\displaystyle }dx^{}d^2x^{}x^2[{\displaystyle \frac{^{}}{^+}}[\xi ^{}(\stackrel{~}{\sigma }^{}A^{})]\stackrel{~}{\sigma }^{}\xi +\xi ^{}(\stackrel{~}{\sigma }^{}A^{}){\displaystyle \frac{1}{^+}}(\stackrel{~}{\sigma }^{}^{})\xi `$ (49) $`+({\displaystyle \frac{^{}}{^+}}\xi ^{})\stackrel{~}{\sigma }^{}(\stackrel{~}{\sigma }^{}A^{})\xi +\xi ^{}{\displaystyle \frac{1}{^+}}(\stackrel{~}{\sigma }^{}^{})(\stackrel{~}{\sigma }^{}A^{})\xi `$ (50) $`m{\displaystyle \frac{1}{^+}}[\xi ^{}(\stackrel{~}{\sigma }^{}A^{})]\xi +m\xi ^{}(\stackrel{~}{\sigma }^{}A^{}){\displaystyle \frac{1}{^+}}\xi `$ (51) $`+m({\displaystyle \frac{1}{^+}}\xi ^{})(\stackrel{~}{\sigma }^{}A^{})\xi m\xi ^{}{\displaystyle \frac{1}{^+}}[(\stackrel{~}{\sigma }^{}A^{})\xi ]]`$ (52) $`{\displaystyle \frac{1}{2}}g^2{\displaystyle 𝑑x^{}d^2x^{}x^2\left[\xi ^{}\stackrel{~}{\sigma }^{}A^{}\frac{1}{i^+}\stackrel{~}{\sigma }^{}(A^{}\xi )\frac{1}{i^+}(\xi ^{}\stackrel{~}{\sigma }^{}A^{})\stackrel{~}{\sigma }^{}A^{}\xi \right]}.`$ (53) We have introduced $`\stackrel{~}{\sigma }^1=\sigma ^2`$ and $`\stackrel{~}{\sigma }^2=\sigma ^1`$. Next consider the gluonic part of the operator $`F^2`$: $`F_g^2={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left[x^{}\mathrm{\Theta }_g^{+2}x^2\mathrm{\Theta }_g^+\right]},`$ (54) where $`\mathrm{\Theta }_g^{+2}`$ $`=F^{+\lambda a}F_\lambda ^{2a},`$ (55) $`\mathrm{\Theta }_g^+`$ $`=F^{+\lambda a}F_\lambda ^a+{\displaystyle \frac{1}{4}}g^+(F_{\lambda \sigma a})^2.`$ (56) Using the constraint equation $`{\displaystyle \frac{1}{2}}^+A^a=^iA^{ia}+gf^{abc}{\displaystyle \frac{1}{^+}}(A^{ib}^+A^{ic})+2g{\displaystyle \frac{1}{^+}}\left(\xi ^{}T^a\xi \right),`$ (57) we arrive at $`F_g^2=F_{g(free)}^2+F_{g(int)}^2`$ (58) where $`F_{g(free)}^2`$ $`={\displaystyle \frac{1}{2}}{\displaystyle }dx^{}d^2x^{}\{x^{}(A^{ja}^+^jA^{2a}A^{2a}^+^jA^{ja}+A^{ja}^+^2A^{ja})`$ (61) $`x^2\left(A^{ka}(^j)^2A^{ka}\right)\}`$ $`2{\displaystyle 𝑑x^{}d^2x^{}A^{2a}^1A^{1a}}.`$ The interaction part $`F_{g(int)}^2`$ $`={\displaystyle \frac{1}{2}}{\displaystyle }dx^{}d^2x^{}x^{}\{gf^{abc}^+A^{ia}A^{2b}A^{ic}`$ (66) $`+g(f^{abc}{\displaystyle \frac{1}{^+}}(A^{ib}^+A^{ic})+2{\displaystyle \frac{1}{^+}}(\xi ^{}T^a\xi ))^+A^{2a}\}`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }dx^{}d^2x^{}x^2\{2gf^{abc}^iA^{ja}A^{ib}A^{jc}+{\displaystyle \frac{g^2}{2}}f^{abc}f^{ade}A^{ib}A^{jc}A^{id}A^{je}`$ $`+2g^iA^{ia}{\displaystyle \frac{1}{^+}}\left(f^{abc}A^{jb}^+A^{jc}+2\xi ^{}T^a\xi \right)`$ $`+g^2(f^{abc}{\displaystyle \frac{1}{^+}}(A^{ib}^+A^{ic})+2{\displaystyle \frac{1}{^+}}\xi ^{}T^a\xi )(f^{ade}{\displaystyle \frac{1}{^+}}(A^{jd}^+A^{je})+2{\displaystyle \frac{1}{^+}}\xi ^{}T^a\xi )\}.`$ So the full transverse rotation operator in QCD can be written as, $`F^2=F_I^2+F_{II}^2+F_{III}^2,`$ (67) where $`F_I^2`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}[x^{}𝒫_0^2x^2(_0+𝒱)]},`$ (68) $`F_{II}^2`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left[\xi ^{}\left[\sigma ^3^1+i^2\right]\frac{1}{^+}\xi +\left[\frac{1}{^+}(^1\xi ^{}\sigma ^3i^2\xi ^{})\right]\xi \right]}`$ (71) $`+{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}m\left[\xi ^{}\left[\frac{\sigma ^1}{i^+}\xi \right]\left[\frac{1}{i^+}\xi ^{}\sigma ^1\right]\xi \right]}`$ $`+{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}g\left[\xi ^{}\frac{1}{^+}[(i\sigma ^3A^1+A^2)\xi ]+\frac{1}{^+}[\xi ^{}(i\sigma ^3A^1+A^2)]\xi \right]},`$ $`F_{III}^2`$ $`={\displaystyle 𝑑x^{}d^2x^{}2(^1A^1)A^2}`$ (73) $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}g\frac{4}{^+}(\xi ^{}T^a\xi )A^{2a}}{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}gf^{abc}\frac{2}{^+}(A^{ib}^+A^{ic})A^{2a}}`$ where $`𝒫_0^i`$ is the free momentum density, $`_o`$ is the free Hamiltonian density and $`𝒱`$ are the interaction terms in the Hamiltonian in manifestly Hermitian form (see Appendix B). The operators $`F_{II}^2`$ and $`F_{III}^2`$ whose integrands do not explicitly depend upon coordinates arise from the fermionic and bosonic parts respectively of the gauge invariant, symmetric, energy momentum tensor in QCD. The above separation is slightly different from that in . From Eq. (7) in Sec. II it follows that the transverse spin operators $`𝒥^i`$, ($`i=1,2`$) can also be written as the sum of three parts, $`𝒥_I^i`$ whose integrand has explicit coordinate dependence, $`𝒥_{II}^i`$ which arises from the fermionic part, and $`𝒥_{III}^i`$ which arises from the bosonic part of the energy momentum tensor. ## IV Summary and Conclusions We have initiated the study of spin operators in QCD. In equal time quantization, one encounters two major difficulties in the description of the spin of a composite system in an arbitrary reference frame. They are 1) the complicated interaction dependence arising from dynamical boost operators and 2) the difficulty in the separation of center of mass motion from the internal motion. Due to these severe difficulties, there have been hardly any attempt to study spin operators of a moving composite system in the conventional equal time formulation of quantum field theory. In light-front theory, on the other hand, the longitudinal spin operator (light-front helicity) is interaction independent and the interaction dependence of transverse spin operators arises solely from that of transverse rotation operators. Moreover, in this case the separation of center of mass motion from internal motion is trivial since light-front transverse boosts are simple Galilean boosts. We have investigated the case of transverse spin operators for both massive and massless particles. A novel feature here is the introduction of transverse spin operators for massless particles with arbitrary transverse momentum. To the best of our knowledge, this is done for the first time in light-front field theory. To provide physical intuition for transverse spin operators which have a complicated structure in interaction theory, we have provided the explicit form of these operators in Fock space basis for both free fermion field theory and free massless spin one field theory. In QCD, our starting point is the formula for transverse rotation operators expressed as the integral of generalized angular momentum density given in terms of gauge invariant, symmetric, energy momentum tensor. We have emphasized the differences between spin operators in field theory in equal time and light-front quantization schemes. Appropriate to light-front quantization, we choose the light-front gauge. We use the constraint equations for $`\psi ^{}`$ and $`A^{}`$ to eliminate them in favor of dynamical degrees of freedom. In this initial study, we restrict to topologically trivial sector of QCD and set the requirement that the transverse gauge fields vanish as $`x^{,i}\mathrm{}`$. This eliminates the surface terms and completely fixes the gauge. In the gauge fixed theory we found that the transverse rotation operators can be decomposed as the sum of three distinct terms: $`F_I^i`$ which has explicit coordinate dependence in its integrand, and $`F_{II}^i`$ and $`F_{III}^i`$ which have no explicit coordinate dependence in their integrand. Further, $`F_{II}^i`$ and $`F_{III}^i`$ arise from the fermionic and bosonic parts of the energy momentum tensor. Since transverse spin is responsible for the helicity flip of the nucleon in light-front theory, we now have identified the complete set of operators responsible for the helicity flip of the nucleon. It is extremely interesting to contrast the cases of longitudinal and transverse spin operators in light-front field theory. In the case of longitudinal spin operator (light-front helicity), in the gauge fixed theory, the operator is interaction independent and can be separated into orbital and spin parts for quarks and gluons. It is known for a long time that the transverse spin operators in light-front field theory cannot be separated into orbital and spin parts except in the trivial case of free field theory. In this work, we have shown that, in spite of the complexities, a physically interesting separation is indeed possible for the transverse spin operators which is quite different from the separation into orbital and spin parts in the rest frame familiar in the equal time picture. In light-front theory, in addition to the Hamiltonian, transverse spin operators also contain interactions and have a complicated structure. Since transverse rotational symmetry is not manifest in light-front theory a study of these operators is essential for questions regarding Lorentz invariance in the theory. An important issue in the case of transverse spin operators concerns renormalization. Since they are interaction dependent, they will acquire divergences in perturbation theory just like the Hamiltonian. It is of interest to find the physical meaning of these divergences and their renormalization. We address these issues in Ref. by computing the expectation value of the transverse spin operators in a dressed quark state. In this work we have explored in detail the theoretical aspects of spin operators in quantum field theory in the context of QCD and their consequences. Our construction and decomposition of the transverse spin operators in QCD also have important phenomenological consequences. Elsewhere, we have shown that nucleon expectation values of $`F_{II}^i`$ and $`F_{III}^i`$ are directly related to the integrals of quark and gluon distribution functions that appear in transversely polarized deep inelastic scattering. Our results show that one can relate nucleon expectation values of operators appearing in the transverse spin to transversely polarized deep inelastic scattering. It is interesting to establish a transverse spin sum rule in analogy to the helicity sum rule and explore its phenomenological consequences. ###### Acknowledgements. We acknowledge helpful conversations with Rajen Kundu, Samir Mallik, Partha Mitra, Jianwei Qiu and James P. Vary. RR gratefully acknowledges the financial assistance of the Council of Scientific and Industrial Research (CSIR), India. ## A Intrinsic Spin in Relativistic Quantum Mechanics In this appendix, for the sake of clarity and completeness we review the intrinsic spin operators in relativistic quantum mechanics . The unitary representations of the Poincare group can be usefully classified on the basis of sign of $`M^2`$, where $`M^2=P^\mu P_\mu `$ ( and further by the sign of $`H`$ in case $`M^20`$). We consider two classes of representations which are of physical importance: * Positive time-like representations: $`M^2>0H>0`$ * Positive light-like representations: $`M^2=0H>0`$ In either cases we do not demand that the representations be irreducible (this allows us to deal with elementary and composite systems simultaneously). #### a Positive time-like representations Beginning from the basic generators $`H,𝐏,𝐉`$, and $`𝐊`$ one can construct an operator $`𝐒`$ such that it is translationally invariant, transforms as a three vector under pure rotations and within itself obeys $`SU(2)`$ commutation relations. $`[S^j,P^\mu ]=0,[J^j,S^k]=iϵ^{jkl}S^l,[S^j,S^k]=iϵ^{jkl}S^l.`$ (A1) A suitable solution to the above requirements is provided by $`𝐒=`$ $`{\displaystyle \frac{1}{M}}\left[𝐖{\displaystyle \frac{𝐏W^0}{M+H}}\right]`$ (A3) $`=𝐉{\displaystyle \frac{P^0}{M}}𝐊\times {\displaystyle \frac{𝐏}{M}}{\displaystyle \frac{(𝐉𝐏)}{M+P^0}}{\displaystyle \frac{𝐏}{M}}`$ where $`𝐖`$ are the space components of the Pauli-Lubanski operator, $`W^\mu =\frac{1}{2}ϵ^{\mu \nu \rho \lambda }M_{\nu \rho }P_\lambda `$. The operators $`𝐒`$ cease to be defined when $`M`$ tends to zero. The commutation relations among $`𝐏,𝐒`$ and $`M`$ are given by $`[P^j,S^k]=0,[S^j,S^k]=iϵ^{jkl}S^l,[S^j,M]=0.`$ (A4) Since $`𝐏`$ and $`M`$ stand for the momentum and invariant mass of the system, the above relations make clear that $`𝐒`$ should represent ‘intrinsic spin’ of the system. The invariant $`W^2`$ can be completely expressed in terms of $`M`$ and $`𝐒`$ as $`W^2=M^2𝐒^2.`$ (A5) #### b Positive light-like representations Begining from the basic generators $`𝐏`$, $`𝐉`$ and $`𝐊`$ (here $`H=\|𝐏\|`$) one has to construct operators $`S`$, $`𝒯^1`$ and $`𝒯^2`$ such that they commute with four momentum $`P^\mu `$ and amongst themselves satisfy $`E(2)`$ commutation relations: $$\begin{array}{ccc}[S,𝒯^1]=i𝒯^2,\hfill & [S,𝒯^2]=i𝒯^1,\hfill & [𝒯^1,𝒯^2]=0.\hfill \end{array}$$ (A6) A suitable solution consistent with the above requirements is: $`S={\displaystyle \frac{W^0}{𝐏}},`$ (A7) $`𝒯^1=W^1P^1{\displaystyle \frac{(W^3+W^0)}{(𝐏+P^3)}},`$ (A8) $`𝒯^2=W^2P^2{\displaystyle \frac{(W^3+W^0)}{(𝐏+P^3)}}.`$ (A9) Note that although $`𝒯_1`$ and $`𝒯_2`$ coincide with the front definitions, the difference lies in the remaining component. Note that here $`S`$ is the component of angular momentum in the direction of motion. To further bring out the difference, we note in passing that $`S`$ is a scalar under pure spatial rotation, while shows complicated behaviour under pure boosts. Contrast this with the fact that $`𝒥^3`$ is front boost invariant. #### c Comments The generators for a multi-particle relativistic system have been analyzed by several authors . The expressions obtained are too complicated to be used in any practical calculations and the generators cannot be easily separated into the center of mass and internal variables. Moreover, the derivations have been done neglecting the field theoretical effects such as pair creation and crossing and so are expected to be valid in the relatively low energy region where an expansion in $`\frac{v}{c}`$ is permissible. Interactions are to be incorporatated by introducing an effective potential which vanish sufficiently rapidly for large distance. ## B Poincare generators in light-front QCD In this appendix we derive the manifestly hermitian kinematical Poincare generators (except $`J^3`$) and the Hamiltonian in light-front QCD starting from the gauge invariant symmetric energy momentum tensor $`\mathrm{\Theta }^{\mu \nu }`$. To begin with, $`\mathrm{\Theta }^{\mu \nu }`$ is interaction dependent. In the gauge fixed theory we find that the seven kinematical generators are manifestly independent of the interaction. We shall work in the gauge $`A^+=0`$ and ignore all surface terms. Thus we are working in the completely gauge fixed sector of the theory. The explicit form of the operator $`J^3`$ in this case is given in Ref. which is manifestly free of interaction at the operator level. The rotation operators are given in Sec. III. At $`x^+=0`$, the operators $`K^3`$ and $`E^i`$ depend only on the density $`\mathrm{\Theta }^{++}`$. A straightforward calculation leads to $`\mathrm{\Theta }^{++}=\psi _{}^{+}{}_{}{}^{}\stackrel{}{i^+}\psi ^++^+A^i^+A^i.`$ (B1) Then, longitudinal momentum operator, $`P^+`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\mathrm{\Theta }^{++}}`$ (B3) $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left[\psi _{}^{+}{}_{}{}^{}\stackrel{}{i^+}\psi ^++^+A^j^+A^j\right]}.`$ Generator of longitudinal scaling, $`K^3`$ $`={\displaystyle \frac{1}{4}}{\displaystyle 𝑑x^{}d^2x^{}x^{}\mathrm{\Theta }^{++}},`$ (B5) $`={\displaystyle \frac{1}{4}}{\displaystyle 𝑑x^{}d^2x^{}x^{}\left[\psi _{}^{+}{}_{}{}^{}\stackrel{}{i^+}\psi ^++^+A^j^+A^j\right]}.`$ Transverse boost generators, $`E^i`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}x^i\mathrm{\Theta }^{++}},`$ (B7) $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}x^i\left[\psi _{}^{+}{}_{}{}^{}\stackrel{}{i^+}\psi ^++^+A^j^+A^j\right]}.`$ The transverse momentum operator $`P^i={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\mathrm{\Theta }^{+i}}`$ (B8) which appears to have explicit interaction dependence. Using the constraint equations for $`\psi ^{}`$ and $`A^{}`$, we have $`\mathrm{\Theta }^{+i}`$ $`=\mathrm{\Theta }_F^{+i}+\mathrm{\Theta }_G^{+i},`$ (B9) $`\mathrm{\Theta }_F^{+i}`$ $`=2\psi _{}^{+}{}_{}{}^{}i^i\psi ^++2g\psi _{}^{+}{}_{}{}^{}A^i\psi ^+,`$ (B10) $`\mathrm{\Theta }_G^{+i}`$ $`=^+A^j^iA^j^+A^j^jA^i+^+A^i^jA^j2g\psi _{}^{+}{}_{}{}^{}A^i\psi ^+.`$ (B11) Thus $`P^i={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left[\psi _{}^{+}{}_{}{}^{}\stackrel{}{i^i}\psi ^++A^j^+^jA^iA^i^+^jA^jA^j^+^iA^j\right]}.`$ (B12) Thus we indeed verify that all the kinematical operators are explicitly independent of interactions. Lastly, the Hamiltonian operator can be written in the manifestly Hermitian form as, $`P^{}={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\mathrm{\Theta }^+}={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}(_0+_{int})}`$ (B13) where $`_0`$ is the free part given by, $`_0=A_a^j(^i)^2A_a^j+\xi ^{}\left[{\displaystyle \frac{(^{})^2+m^2}{i^+}}\right]\xi \left[{\displaystyle \frac{(^{})^2+m^2}{i^+}}\xi ^{}\right]\xi .`$ (B14) The interaction terms are given by, $`_{int}=_{qqg}+_{ggg}+_{qqgg}+_{qqqq}+_{gggg},`$ (B15) where, $`_{qqg}=4g\xi ^{}{\displaystyle \frac{1}{^+}}(^{}.A^{})\xi +g{\displaystyle \frac{^{}}{^+}}[\xi ^{}(\stackrel{~}{\sigma }^{}A^{})]\stackrel{~}{\sigma }^{}\xi +g\xi ^{}(\stackrel{~}{\sigma }^{}A^{}){\displaystyle \frac{1}{^+}}(\stackrel{~}{\sigma }^{}^{})\xi `$ (B16) $`+g({\displaystyle \frac{^{}}{^+}}\xi ^{})\stackrel{~}{\sigma }^{}(\stackrel{~}{\sigma }^{}A^{})\xi +g\xi ^{}{\displaystyle \frac{1}{^+}}(\stackrel{~}{\sigma }^{}^{})(\stackrel{~}{\sigma }^{}A^{})\xi `$ (B17) $`mg{\displaystyle \frac{1}{^+}}[\xi ^{}(\stackrel{~}{\sigma }^{}A^{})]\xi +mg\xi ^{}(\stackrel{~}{\sigma }^{}A^{}){\displaystyle \frac{1}{^+}}\xi `$ (B18) $`+mg({\displaystyle \frac{1}{^+}}\xi ^{})(\stackrel{~}{\sigma }^{}A^{})\xi mg\xi ^{}{\displaystyle \frac{1}{^+}}[(\stackrel{~}{\sigma }^{}A^{})\xi ],`$ (B19) $`_{ggg}=2gf^{abc}\left[^iA_a^jA_b^iA_c^j+(^iA_a^i){\displaystyle \frac{1}{^+}}(A_b^j^+A_c^j)\right],`$ (B20) $`_{qqgg}`$ $`=g^2[\xi ^{}(\stackrel{~}{\sigma }^{}.A^{}){\displaystyle \frac{1}{i^+}}(\stackrel{~}{\sigma }^{}.A^{})\xi {\displaystyle \frac{1}{i^+}}(\xi ^{}\stackrel{~}{\sigma }^{}.A^{})\stackrel{~}{\sigma }^{}.A^{}\xi `$ (B22) $`+4{\displaystyle \frac{1}{^+}}(f^{abc}A_b^i^+A_c^i){\displaystyle \frac{1}{^+}}(\xi ^{}T^a\xi )],`$ $`_{qqqq}=4g^2{\displaystyle \frac{1}{^+}}(\xi ^{}T^a\xi ){\displaystyle \frac{1}{^+}}(\xi ^{}T^a\xi ),`$ (B23) $`_{gggg}=`$ $`{\displaystyle \frac{g^2}{2}}f^{abc}f^{ade}[A_b^iA_c^jA_d^iA_e^j`$ (B25) $`+2{\displaystyle \frac{1}{^+}}(A_b^i^+A_c^i){\displaystyle \frac{1}{^+}}(A_d^j^+A_e^j)].`$ ## C Transverse Spin in free fermion field theory ### 1 Poincare Generators: Operator Forms The symmetric energy momentum tensor $`\mathrm{\Theta }^{\mu \nu }=\left[\overline{\psi }\gamma ^\mu {\displaystyle \frac{1}{4}}\stackrel{}{i^\nu }\psi +\overline{\psi }\gamma ^\nu {\displaystyle \frac{1}{4}}\stackrel{}{i^\mu }\psi \right].`$ (C1) The momentum operators are given by $`P^+`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\overline{\psi }\gamma ^+\frac{1}{2}}\stackrel{}{i^+}\psi `$ (C3) $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}[\xi ^{}i^+(i^+\xi ^{})]\xi }.`$ $`P^i`$ $`={\displaystyle \frac{1}{2}}{\displaystyle }dx^{}d^2x^{}\left[\overline{\psi }\{\gamma ^+{\displaystyle \frac{1}{4}}\stackrel{}{i^i}+\gamma ^i{\displaystyle \frac{1}{4}}\stackrel{}{i^+}\}\psi \right]`$ (C5) $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}[\xi ^{}i^i(i^i\xi ^{})]\xi }.`$ The Hamiltonian operator is $`P^{}`$ $`={\displaystyle \frac{1}{2}}{\displaystyle }dx^{}d^2x^{}\overline{\psi }[\gamma ^{}{\displaystyle \frac{1}{4}}\stackrel{}{i^+}+\gamma ^+{\displaystyle \frac{1}{4}}\stackrel{}{i^{}}]\psi `$ (C7) $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left[\xi ^{}\frac{1}{i^+}[m_F^2(^{})^2](\frac{1}{i^+}[m_F^2(^{})^2]\xi ^{})\right]\xi }.`$ The longitudinal scaling operator (at $`x^+=0`$) is $`K^3`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}x^{}\left[\overline{\psi }\gamma ^+\frac{1}{4}\stackrel{}{i^+}\psi \right]}`$ (C9) $`={\displaystyle \frac{i}{4}}{\displaystyle 𝑑x^{}d^2x^{}x^{}\left[\xi ^{}^+\xi (^+\xi ^{})\xi \right]}.`$ The transverse boost operators are $`E^i`$ $`={\displaystyle \frac{1}{4}}{\displaystyle 𝑑x^{}d^2x^{}x^i\left[\overline{\psi }\gamma ^+\frac{1}{4}\stackrel{}{i^+}\psi \right]}`$ (C11) $`={\displaystyle \frac{1}{4}}{\displaystyle 𝑑x^{}d^2x^{}x^i\left[\xi ^{}i^+(i^+\xi ^{})\right]\xi }.`$ The generators of rotations are $`J^3`$ $`={\displaystyle \frac{1}{2}}{\displaystyle }dx^{}d^2x^{}\{x^1\left[\overline{\psi }\{\gamma ^+{\displaystyle \frac{1}{4}}\stackrel{}{i^2}+\gamma ^2{\displaystyle \frac{1}{4}}\stackrel{}{i^+}\}\psi \right]`$ (C15) $`x^2\left[\overline{\psi }\{\gamma ^+{\displaystyle \frac{1}{4}}\stackrel{}{i^1}+\gamma ^1{\displaystyle \frac{1}{4}}\stackrel{}{i^+}\}\psi \right]\}`$ $`={\displaystyle }dx^{}d^2x^{}[\xi ^{}[{\displaystyle \frac{i}{2}}(x^1\stackrel{}{}^2x^2\stackrel{}{}^1)\xi \left[{\displaystyle \frac{i}{2}}(x^1\stackrel{}{}^2x^2\stackrel{}{}^1)\xi ^{}\right]\xi `$ $`+\xi ^{}{\displaystyle \frac{\sigma _3}{2}}\xi ].`$ and $`F^i`$ $`={\displaystyle \frac{1}{2}}{\displaystyle }dx^{}d^2x^{}\{x^{}\left[\overline{\psi }\{\gamma ^+{\displaystyle \frac{1}{4}}\stackrel{}{i^i}+\gamma ^i{\displaystyle \frac{1}{4}}\stackrel{}{i^+}\}\psi \right]`$ (C21) $`x^i\left[\overline{\psi }\{\gamma ^+{\displaystyle \frac{1}{4}}\stackrel{}{i^{}}+{\displaystyle \frac{1}{4}}\gamma ^{}\stackrel{}{i^+}\}\psi \right]\}`$ $`={\displaystyle \frac{i}{2}}{\displaystyle }dx^{}d^2x^{}\xi ^{}[x^i(m^2(^{})^2){\displaystyle \frac{1}{^+}}x^{}{\displaystyle \frac{}{x^i}}`$ $`+{\displaystyle \frac{1}{^+}}\{{\displaystyle \frac{}{x^i}}iϵ^{ij}\sigma ^3{\displaystyle \frac{}{x^j}}+ϵ^{ij}m\sigma ^j\}]\xi `$ $`{\displaystyle \frac{i}{2}}{\displaystyle }dx^{}d^2x^{}[[x^i(m^2(^{})^2){\displaystyle \frac{1}{^+}}x^{}{\displaystyle \frac{}{x^i}}`$ $`+{\displaystyle \frac{1}{^+}}\{{\displaystyle \frac{}{x^i}}+iϵ^{ij}\sigma ^3{\displaystyle \frac{}{x^j}}+ϵ^{ij}m\sigma ^j\}\xi ^{}]]\xi .`$ ### 2 Fock Representation Free spin-half field operator is $`\xi (x)={\displaystyle \underset{\lambda }{}}\chi _\lambda {\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3\sqrt{k^+}}[b(k,\lambda )e^{ik.x}+d^{}(k,\lambda )e^{ik.x}]}.`$ (C22) In terms of Fock space operators $`P^+={\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}k^+\underset{\lambda }{}\left[b^{}(k,\lambda )b(k,\lambda )+d^{}(k,\lambda )d(k,\lambda )\right]}.`$ (C23) $`P^i={\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}k^i\underset{\lambda }{}\left[b^{}(k,\lambda )b(k,\lambda )+d^{}(k,\lambda )d(k,\lambda )\right]}.`$ (C24) $`P^{}={\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\frac{m_F^2+(k^{})^2}{k^+}\underset{\lambda }{}\left[b^{}(k,\lambda )b(k,\lambda )+d^{}(k,\lambda )d(k,\lambda )\right]}.`$ (C25) $`K^3`$ $`={\displaystyle \frac{i}{2}}{\displaystyle }{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}}k^+{\displaystyle \underset{\lambda }{}}([{\displaystyle \frac{b^{}(k,\lambda )}{k^+}}b(k,\lambda )+{\displaystyle \frac{d^{}(k,\lambda )}{k^+}}d(k,\lambda )]`$ (C27) $`[b^{}(k,\lambda ){\displaystyle \frac{b(k,\lambda )}{k^+}}+d^{}(k,\lambda ){\displaystyle \frac{d(k,\lambda )}{k^+}}]).`$ $`E^i`$ $`={\displaystyle \frac{i}{2}}{\displaystyle }{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}}{\displaystyle \underset{\lambda }{}}k^+([{\displaystyle \frac{b^{}(k,\lambda )}{k^i}}b(k,\lambda )+{\displaystyle \frac{d^{}(k,\lambda )}{k^i}}d(k,\lambda )]`$ (C29) $`[b^{}(k,\lambda ){\displaystyle \frac{b(k,\lambda )}{k^i}}+d^{}(k,\lambda ){\displaystyle \frac{d(k,\lambda )}{k^i}}]).`$ $`J^3`$ $`={\displaystyle \frac{i}{2}}{\displaystyle }{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}}{\displaystyle \underset{\lambda }{}}[\left([k^1{\displaystyle \frac{}{k^2}}k^2{\displaystyle \frac{}{k^1}}]b^{}(k,\lambda )\right)b(k,\lambda )b^{}(k,\lambda )[k^1{\displaystyle \frac{}{k^2}}k^2{\displaystyle \frac{}{k^1}}]b(k,\lambda )`$ (C32) $`+\left([k^1{\displaystyle \frac{}{k^2}}k^2{\displaystyle \frac{}{k^1}}]d^{}(k,\lambda )\right])d(k,\lambda )d^{}(k,\lambda )[k^1{\displaystyle \frac{}{k^2}}k^2{\displaystyle \frac{}{k^1}}]d(k,\lambda )]`$ $`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\underset{\lambda }{}\lambda \left[b^{}(k,\lambda )b(k,\lambda )+d^{}(k,\lambda )d(k,\lambda )\right]}`$ with $`\lambda =\pm 1`$. $`F^i`$ $`=i{\displaystyle }{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}}k^i{\displaystyle \underset{\lambda }{}}([{\displaystyle \frac{b^{}(k,\lambda )}{k^+}}b(k,\lambda )+{\displaystyle \frac{d^{}(k,\lambda )}{k^+}}d(k,\lambda )]`$ (C38) $`[b^{}(k,\lambda ){\displaystyle \frac{b(k,\lambda )}{k^+}}+d^{}(k,\lambda ){\displaystyle \frac{d(k,\lambda )}{k^+}}])`$ $`+{\displaystyle \frac{i}{2}}{\displaystyle }{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}}{\displaystyle \frac{m_F^2+(k^{})^2}{k^+}}{\displaystyle \underset{\lambda }{}}([{\displaystyle \frac{b^{}(k,\lambda )}{k^i}}b(k,\lambda )+{\displaystyle \frac{d^{}(k,\lambda )}{k^i}}d(k,\lambda )]`$ $`[b^{}(k,\lambda ){\displaystyle \frac{b(k,\lambda )}{k^i}}+d^{}(k,\lambda ){\displaystyle \frac{d(k,\lambda )}{k^i}}])`$ $`{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\frac{ϵ^{ij}}{k^+}k^j\underset{\lambda \lambda ^{}}{}\sigma _{\lambda \lambda ^{}}^3\left[b^{}(k,\lambda )b(k,\lambda ^{})d^{}(k,\lambda ^{})d(k,\lambda )\right]}`$ $`{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\frac{ϵ^{ij}}{k^+}m_F\underset{\lambda \lambda ^{}}{}\sigma _{\lambda \lambda ^{}}^j\left[b^{}(k,\lambda )b(k,\lambda ^{})+d^{}(k,\lambda )d(k,\lambda ^{})\right]}.`$ ### 3 Transverse Spin of a Single Fermion For a single fermion of mass $`m`$ and momenta ($`k^+,k^{}`$), we have, $`P^+k\lambda `$ $`=k^+k\lambda ,P^1k\lambda =k^1k\lambda ,P^2k\lambda =k^2k\lambda ,`$ (C39) $`P^{}k\lambda `$ $`={\displaystyle \frac{(k^{})^2+m^2}{k^+}}k\lambda ,𝒥^3k\lambda ={\displaystyle \frac{1}{2}}\lambda k\lambda ,`$ (C40) $`K^3k\lambda `$ $`=ik^+{\displaystyle \frac{}{k^+}}k\lambda ,E^1k\lambda =ik^+{\displaystyle \frac{}{k^1}}k\lambda ,E^2k\lambda =ik^+{\displaystyle \frac{}{k^2}}k\lambda ,`$ (C41) $`F^1k\lambda `$ $`=\left(2ik^1{\displaystyle \frac{}{k^+}}+i{\displaystyle \frac{(k^{})^2+m^2}{k^+}}{\displaystyle \frac{}{k^1}}{\displaystyle \frac{k^2}{k^+}}\lambda \right)k\lambda {\displaystyle \frac{m}{k^+}}{\displaystyle \underset{\lambda ^{}}{}}\sigma _{\lambda ^{}\lambda }^2k\lambda ^{},`$ (C42) $`F^2k\lambda `$ $`=\left(2ik^2{\displaystyle \frac{}{k^+}}+i{\displaystyle \frac{(k^{})^2+m^2}{k^+}}{\displaystyle \frac{}{k^2}}+{\displaystyle \frac{k^1}{k^+}}\lambda \right)k\lambda +{\displaystyle \frac{m}{k^+}}{\displaystyle \underset{\lambda ^{}}{}}\sigma _{\lambda ^{}\lambda }^1k\lambda ^{}.`$ (C43) We arrive at $`m𝒥^1k\lambda `$ $`=\left({\displaystyle \frac{1}{2}}F^2P^++K^3P^2{\displaystyle \frac{1}{2}}E^2P^{}P^1𝒥^3\right)k\lambda `$ (C45) $`=m{\displaystyle \underset{\lambda ^{}}{}}{\displaystyle \frac{\sigma _{\lambda ^{}\lambda }^1}{2}}k\lambda ^{},`$ $`m𝒥^2k\lambda `$ $`=\left({\displaystyle \frac{1}{2}}F^1P^+K^3P^1+{\displaystyle \frac{1}{2}}E^1P^{}P^2𝒥^3\right)k\lambda `$ (C47) $`=m{\displaystyle \underset{\lambda ^{}}{}}{\displaystyle \frac{\sigma _{\lambda ^{}\lambda }^2}{2}}k\lambda ^{}.`$ ## D Transverse spin in free massless spin one field theory ### 1 Poincare Generators: Operator Forms The symmetric gauge invariant energy momentum tensor $`\mathrm{\Theta }^{\mu \nu }=F^{\lambda \mu }F_\lambda ^\nu g^{\mu \nu }.`$ (D1) where the Lagrangian density $`={\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu }`$ (D2) with $`F^{\mu \nu }=^\nu A^\mu ^\mu A^\nu .`$ (D3) We choose $`A^+=0`$ gauge. Only the transverse fields $`A^i`$ are dynamical variables. The momentum operators are given by $`P^+={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}^+A^j^+A^j},`$ (D4) $`P^i={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left(A^j^+^jA^iA^i^+^jA^jA^j^+^iA^j\right)}.`$ (D5) The Hamiltonian operator is $`P^{}`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left[\frac{1}{4}(^+A^{})^2+\frac{1}{2}F^{ij}F_{ij}\right]}`$ (D7) $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}^iA^j^iA^j}={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}A^j(^i)^2A^j}`$ The longitudinal scale generator (at $`x^+=0`$) is $`K^3={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}x^{}^+A^j^+A^j}.`$ (D8) The transverse boost generators are $`E^i={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}x^i^+A^j^+A^j}.`$ (D9) The generators of rotations are $`J^3`$ $`={\displaystyle \frac{1}{2}}{\displaystyle }dx^{}d^2x^{}(x^1[^+A^2^iA^i+^+A^1(^2A^1^1A^2)]`$ (D13) $`x^2[^+A^1^iA^i+^+A^2(^2A^1+^1A^2)])`$ $`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}\left(x^1[^+A^1^2A^1+^+A^2^2A^2]x^2[^+A^1^1A^1+^+A^2^1A^2]\right)}`$ $`+{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x^{}[A^1^+A^2A^2^+A^1]}.`$ and $`F^i`$ $`={\displaystyle \frac{1}{2}}{\displaystyle }dx^{}d^2x^{}(x^{}(A^{ja}^+^jA^iA^i^+^jA^jA^j^+^iA^j)`$ (D15) $`x^i[A^k(^j)^2A^k])2{\displaystyle }dx^{}d^2x^{}A^i\eta ^{ij}^jA^j,\mathrm{no}\mathrm{summation}\mathrm{over}i,j,`$ with $`\eta ^{12}=\eta ^{21}=1`$, $`\eta ^{11}=\eta ^{22}=0`$. ### 2 Fock Representation The dynamical components of the free massless spin field operator in $`A^+=0`$ gauge are $`A^i(x)={\displaystyle \underset{\lambda =1}{\overset{2}{}}}{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\delta ^{i\lambda }[a(k,\lambda )e^{ik.x}+a^{}(k,\lambda )e^{ik.x}]}.`$ (D16) In terms of Fock space operators, we have, $`P^+={\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}k^+\underset{\lambda }{}a^{}(k,\lambda )a(k,\lambda )}.`$ (D17) $`P^i={\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}k^i\underset{\lambda }{}a^{}(k,\lambda )a(k,\lambda )}.`$ (D18) $`H={\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\frac{k_{}^{}{}_{}{}^{2}}{k^+}\underset{\lambda }{}a^{}(k,\lambda )a(k,\lambda )}.`$ (D19) $`K^3={\displaystyle \frac{i}{2}}{\displaystyle }{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}}k^+{\displaystyle \underset{\lambda }{}}[\left({\displaystyle \frac{a^{}(k,\lambda )}{k^+}}\right)a(k,\lambda )a^{}(k,\lambda ){\displaystyle \frac{a(k,\lambda )}{k^+}})].`$ (D20) $`E^i={\displaystyle \frac{i}{2}}{\displaystyle }{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}}k^+{\displaystyle \underset{\lambda }{}}[\left({\displaystyle \frac{a^{}(k,\lambda )}{k^i}}\right)a(k,\lambda )a^{}(k,\lambda ){\displaystyle \frac{a(k,\lambda )}{k^i}})].`$ (D21) $`J^3`$ $`={\displaystyle \frac{i}{2}}{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\underset{\lambda }{}\left[\left((k^1\frac{}{k^2}k^2\frac{}{k^1})a^{}(k,\lambda )\right)a(k,\lambda )a^{}(k,\lambda )(k^1\frac{}{k^2}k^2\frac{}{k^1})a(k,\lambda )\right]}`$ (D23) $`+i{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\left(a^{}(k,2)a(k,1)a^{}(k,1)a(k,2)\right)}.`$ Introduce creation and annihilation operators $`a(k,)={\displaystyle \frac{1}{\sqrt{2}}}[a(k,1)ia(k,2)],a(k,)={\displaystyle \frac{1}{\sqrt{2}}}[a(k,1)+ia(k,2)].`$ (D24) Then $`J^3`$ $`={\displaystyle \frac{i}{2}}{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\underset{\lambda }{}\left[\left((k^1\frac{}{k^2}k^2\frac{}{k^1})a^{}(k,\lambda )\right)a(k,\lambda )a^{}(k,\lambda )(k^1\frac{}{k^2}k^2\frac{}{k^1})a(k,\lambda )\right]}`$ (D26) $`+{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\underset{\lambda }{}\lambda a^{}(k,\lambda )a(k,\lambda )}.`$ where $`\lambda `$ now denotes circular polarization. $`F^i=`$ $`i{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}k^i\underset{\lambda }{}\left(\frac{a^{}(k,\lambda )}{k^+}a(k,\lambda )a^{}(k,\lambda )\frac{a(k,\lambda )}{k^+}\right)}`$ (D29) $`+{\displaystyle \frac{i}{2}}{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\frac{(k^{})^2}{k^+}\underset{\lambda }{}\left(\frac{a^{}(k,\lambda )}{k^i}a(k,\lambda )a^{}(k,\lambda )\frac{a(k,\lambda )}{k^i}\right)}`$ $`2ϵ^{ij}{\displaystyle \frac{dk^+d^2k^{}}{2(2\pi )^3k^+}\frac{k^j}{k^+}\underset{\lambda }{}\lambda a^{}(k,\lambda )a(k,\lambda )}.`$ ### 3 Transverse Spin Using the explicit form of the operators, we get for a state of momentum $`k(k^+,k^{})`$ and helicity $`\lambda `$, $`𝒥^3k\lambda ={\displaystyle \frac{W^+}{P^+}}k\lambda `$ $`=\lambda k\lambda ,`$ (D30) $`W^1k\lambda `$ $`=k^1\lambda k\lambda ,`$ (D31) $`W^2k\lambda `$ $`=k^2\lambda k\lambda ,`$ (D32) $`W^{}k\lambda `$ $`={\displaystyle \frac{(k^{})^2}{k^+}}\lambda k\lambda .`$ (D33) $`𝒥^ik\lambda =0.`$ (D34)
warning/0004/cond-mat0004164.html	ar5iv	text	# Energy Dependence of Quasi-Particle Relaxation in a Disordered Fermi Liquid ## Abstract A spectroscopic method is applied to measure the inelastic quasi-particle relaxation rate in a disordered Fermi liquid. The quasi-particle relaxation rate, $`\gamma `$ is deduced from the magnitude of fluctuations in the local density of states, which are probed using resonant tunneling through a localized impurity state. We study its dependence on the excitation energy $`E`$ measured from the Fermi level. In a disordered metal (heavily doped GaAs) we find that $`\gamma E^{3/2}`$ within the experimentally accessible energy interval, in agreement with the Altshuler-Aronov theory for electron-electron interactions in diffusive conductors. The quasi-particle description of electrons represents a common approach to understanding kinetic and thermodynamic phenomena in metals. It relies on a certain stability of quasiparticle excitations in a many-body system, which requires that the quasi-particle excitation energy, $`E`$, determined with respect to the Fermi level, $`E_F`$, exceeds the broadening of such a single-particle state, $`\mathrm{}\gamma (E)`$, due to electron-electron interactions. In ideally pure metals, the relaxation rate, $`\gamma ,`$ of ballistic quasi-particles at low energies is kept low by a diminishing of the phase space available for inelastic electron-electron collision processes . The phase space argument results in $`\gamma E^2`$ dependence, which has been confirmed by electron emission spectroscopy measurements which probe electron states at energies $`E\mathrm{}/\tau `$, where $`\tau `$ is the elastic mean free path time due to residual impurities. Quasi-particle decay at excitation energies $`E<\mathrm{}/\tau `$ is accelerated by the presence of disorder (which makes the electron states chaotically random at length scales longer than the elastic mean free path, $`l=\tau v_F`$). The theory of interaction effects in diffusive media predicts that, at zero temperature, the decay of quasi-particle excitations with $`E\mathrm{}/\tau `$ in a disordered metal or heavily doped semiconductor is slow enough to guarantee the existence of properly defined quasi-particles. In particular, in bulk three-dimensional (d=3) and two-dimensional (d=2) conductors, the quasi-particle relaxation rate is expected to obey a power-law dependence, $$\gamma (E)=aE^{d/2}E_F^{(1d/2)}(\lambda _F/l)^{d/2}.$$ (1) Despite the fundamental interest and clear theoretical predictions , there are very few direct experimental measurements of the quasi-particle decay time in disordered metals at small excitation energies $`\epsilon <\mathrm{}/\tau `$. Information about the inelastic decay of non-equilibrium quasi-particles in dirty metals is often extracted from energy relaxation rates in the electron thermalization process , $`\gamma _E`$. Otherwise, one studies the temperature dependence of a dephasing rate, $`\gamma _\phi (T)`$, of coherent carriers (determined using weak localisation or universal conductance fluctuations analysis), which can be treated as a measure of the efficiency of the interactions of an electron with energy $`ET`$ with equilibrium fluctuations of charges produced by other electrons thermally distributed near the Fermi level in a diffusive conductor . At high temperatures, both phase and energy relaxation experiments show a certain agreement with theoretical estimations . However, recently reported data on two dynamical parameters mentioned above, $`\gamma _E(E)`$ and $`\gamma _\phi (T)`$, in Au wires and films, and also in semiconductor heterostructures have indicated a certain disagreement between theoretically predicted and experimentally observed values of these two quantities, which has refuelled both theoretical and experimental interest in the problem of quasi-particle lifetimes in a disordered metal . In the present paper, we report the results of a direct measurement of the energy dependence of the inelastic decay rate $`\gamma (E)`$ of a quasi-particle state in a disordered conductor. This study employs the method of resonant tunneling spectroscopy using a discrete localised state in a double-barrier structure, which has been applied earlier by A.Geim et al to study 2D electrons in a heterostucture and by U.Sivan et al to investigate the discrete spectra of quantum pillars. It has been shown previously that, by measuring the current-voltage (IV) characteristics and by deriving the differential conductance in a system where the current passes through a single resonant impurity state in the barrier, one can study features of the single-particle spectrum of a disordered metal (playing the role of an emitter). In a bulk material with a continuous spectrum, individual chaotic quantum states formed by the interference of elastically scattered electrons produce an effect known as fluctuations of the local density of states (LDOS) . It consists of a random and coordinate-specific energy dependence of the local density of states in a diffusive metal, $`\nu (E)`$ with a correlation energy limited by inelastic broadening of quasi-particle states, $`\mathrm{}\gamma `$. We observe a random pattern in $`\nu (E)`$ by sweeping a single resonant impurity level against the electron spectrum in the emitter within a finite range of excitation energies for a quasi-hole (an empty state below $`E_F`$ in the emitter) left behind by the tunneled electron. When the energetic width of an impurity level used in this process, $`\mathrm{\Gamma }`$, is smaller than the inelastic broadening of single-particle levels in the emitting electrode, $`\mathrm{}\gamma `$, one can extract the latter characteristic from the analysis of the amplitude of the LDOS fluctuations pattern and its auto-correlation parameters. This has enabled us to measure directly the energy dependence of the inelastic relaxation rate, $`\gamma (E)`$, in bulk degenerate heavily doped GaAs at low temperatures, which we find to agree with the Altshuler-Aronov theory predicting $`\mathrm{}\gamma (E)=a\left(E^{3/2}/\sqrt{E_F}\right)\left(\lambda _F/l\right)^{3/2}`$ for a 3D system. \[Note, that under the condition of $`\mathrm{}\gamma (E)E`$ provided in Eq. (1) for $`\lambda _Fl`$, the inelastic decay rate $`\gamma `$ of a non-equilibrium quasi-particle coincides with its decoherence rate.\] As mentioned above, the LDOS can be measured via resonant tunneling through an impurity in a strongly asymmetric double-barrier heterostructure. Our microstructure consists of a 10 nm wide GaAs quantum well and 5 and 8 nm wide Al<sub>0.3</sub>Ga<sub>0.7</sub>As barriers sandwiched between doped GaAs contact layers with a donor concentration of $`3.3\times 10^{17}`$ cm<sup>-3</sup>, as sketched in Fig. 2(a). From this material we fabricated a 2 $`\mu `$m diameter mesa, as depicted in the scanning electron micrograph of Fig. 2(b). The mesa contains a small number of residual impurities in the nominally undoped quantum well. The energetically lowest impurity state $`S`$ in the well is used as a spectrometer for the LDOS imaging in the metallic emitter adjacent to the thicker barrier as illustrated in Fig. 2(c). At zero bias, $`S`$ lies above the Fermi level in the emitter and is not available for resonant transport, resulting in $`I=0`$ and $`G=0`$. This measurement has been performed at the temperature $`T=`$20mK. Upon applying a finite bias voltage $`V`$, the energetic position of a spectrometer S is shifted down to the energy $`E=\alpha e(VV_S)`$ below the Fermi level in the emitter, where the prefactor $`\alpha =0.50`$ accounts for the fact that only part of the voltage drops between emitter and spectrometer . When S crosses $`E_F`$ from above (at $`V_S=9.8`$mV), the current acquires a finite value (plateaux) limited by the left, less transparent barrier and proportional to the LDOS of occupied states in the emitter at the spectrometer position, $`I\nu `$ . At higher bias voltages ($`V_{H1}=14.6`$mV and $`V_{H2}=15.5`$ mV), other impurities (or, maybe, other excited states from the same impurity) become involved in the current formation, which produces the next prominent current steps in the IV-characteristics. Consequently, the differential conductance $`G=dI/dVd\nu (E)/dE`$ of the device plotted in Fig. 2(d) exhibits several pronounced peaks, the ’main’ one at $`V_S=9.8`$mV followed by two at $`14.6`$ and $`15.5`$ mV, each characterising the energetic width and transparency of a resonant impurity state, whereas in between, at $`V_SVV_{H1}`$, the differential conductance displays the derivative of the LDOS with respect to energy. The energy dependence of the LDOS is the result of the energy-dependent quantum interference pattern for quasi-particles in the emitter at the coordinate of a spectrometer, see Fig. 2(c). This pattern is random and tends to reflect an individual portrait of a disordered potential in a metal surrounding the spectrometer. In the sample under investigation, such a pattern can be analyzed within the energy interval of about $`0E2`$ meV below the Fermi level, since the voltage and energy scales are related via $`E=\alpha e(VV_S),\alpha 0.5`$. The correlation energy, $`E_c=\mathrm{\Gamma }+\mathrm{}\gamma `$, of the fine structure in the differential conductance pattern is determined by either the energetic spectrometer width, $`\mathrm{\Gamma }`$, or by the inelastic broadening of states in the emitter, $`\mathrm{}\gamma `$, whichever is larger. According to the theory , one can relate $`E_c`$ to the amplitude (rms value) of the fluctuation pattern of $`G(V)d\nu (E)/dE`$. Note that, although oscillations at larger energy scales are also present in each realization of $`\nu (E)`$, their contribution to $`G(V)`$ is suppressed, due to the differentiation. For a given sample, the spectrometer width, $`\mathrm{\Gamma }`$, can be extracted from the width of the ’main’ resonance peak. For the peak at $`V_S=9.8`$mV in Fig. 2(d), we find $`\mathrm{\Gamma }e\alpha \times 72\mu `$V $`36\mu `$eV. For a broad spectrometer, with $`\mathrm{\Gamma }\mathrm{}\gamma (E)`$ at any excitation energy , both the amplitude and correlation voltage (energy) of fluctuations would be the same over the entire range of $`V_SVV_{H1}`$. For a narrow spectrometer, such as studied in the present work, inelastic broadening of states in the bulk exceeds the spectrometer width upon increasing the excitation energy of a quasi-hole left in the emitter. This results in a decrease of fluctuations upon increasing voltage, as indicated by the dashed lines in Fig. 2(d). Quantitative information about the quasiparticle decay rate, $`\gamma (E)`$, is obtained from statistical analysis of the complete fluctuation pattern $`G(V,B)=dI/dV`$studied as a function of a magnetic field, $`B`$ (applied parallel to the current flow). Figure 3 shows a color-scale image of the differential conductance measured as a function of bias voltage and magnetic field (in the region of low magnetic fields, where Landau quantization is hindered by disorder). Sharp black lines in Fig. 3 correspond to the spectrometer crossing the emitter Fermi level. The decrease of the amplitude of observed LDOS fluctuations and the increase of the correlation voltage of the fluctuation pattern, as a function of the quasi-hole excitation energy is apparent from the change in the color and contrast of this image. As a quantitative measure of the fluctuation amplitude, we calculate the variance var$`{}_{B}{}^{}G=\delta G^2(B)_B`$ using $`\delta G(B)=G(B)G_B`$, where $`\mathrm{}_B`$ indicates averaging over magnetic field in the range of $`0B1.0`$ T. Figure 4(a) shows that var$`{}_{B}{}^{}G`$ drops by more than one order of magnitude within the experimentally accessible voltage range. In our limit of classical magnetic fields, the fluctuation amplitude is related to the relaxation rate according to $$\text{var}_BG\|_V=G_N^2\times [1+\mathrm{}\gamma (E)/\mathrm{\Gamma }]^{3/2}.$$ (2) Here, $`G_N`$ is a prefactor which we determine as $`G_N^2=`$ var$`{}_{B}{}^{}G_{\|V_S}^{}`$ from Fig. 4(a) by assuming that, at $`V=V_S`$ (corresponding to $`E=0`$), $`\mathrm{}\gamma \mathrm{\Gamma }`$ and $`E=\alpha e(VV_S)`$ is the excitation energy of a quasi-hole (we remind that $`\alpha 0.5`$ in this experiment). Figure 4(b) shows the obtained energy dependence of the quasi-particle decay rate. It drops strongly upon decreasing the excitation energy, and, in contrast to some experiments measuring dephasing rates , we do not observe a saturation of $`\gamma (E)`$ at low energies. In Fig. 4(b) we also compare the experimentally determined quasi-particle relaxation rate with the values calculated using the Altshuler-Aronov theory of electron-electron interaction in disordered conductors . It is worth mentioning that the above-presented determination of the quasi-particle decay rate based upon the LDOS pattern analysis enables us to study quasi-particles with pretty small excitation energies, $`E<\mathrm{}/\tau `$ and to detect the features of their inelastic decay specific to strongly disordered systems. This make it different from the analysis of the same quantity on the basis of measurements of Landau level broadening , which requires distinct Landau quantization and a strong magnetic field (or absence of impurities). Electron-electron scattering with a large momentum transfer between ballistic quasi-particles results in a rate which is determined by the phase volume of available final states, thus leading to $`\gamma E^2`$ . In disordered Fermi liquids, where transport is diffusive, small momentum transfers play an important role, such that an additional $`E^{3/2}`$ energy dependence of $`\gamma `$ appears, as described in Eq. (1). The $`E^{3/2}`$ dependence dominates at small energies, $`E<\mathrm{}/\tau `$, while the $`E^2`$ dependence is specific to large energies $`E\mathrm{}/\tau `$. After estimating the elastic scattering time of $`\tau =0.14`$ps from the magnetic-field dependence of the LDOS fluctuations (and also from the nominal doping level of the emitter contact), we find that the latter crossover would occur at $`E5`$ meV, which is beyond the energy range accessible in the reported experiment. For a quantitative comparison, we fit the relaxation rate as $`\gamma =b\times E^2+A\times E^{3/2}`$ (dashed lines in Fig. 4(b) show separately the $`b\times E^2`$ and $`A\times E^{3/2}`$ parts determined in this fit). We have also fit the data to $`\gamma (E)=A\times E^x`$ dependence, treating exponent $`x`$ as a free parameter, which yields $`x=1.543/2`$ and $`A=8\times 10^{10}`$ meV<sup>-3/2</sup>s<sup>-1</sup>. In the theory , the prefactor $`A`$ in the $`E^{3/2}`$ dependence is $`A=(105\sqrt{3\mathrm{}})/(16\pi \tau ^{3/2}E_F^2)`$. After estimating $`E_F=26`$meV from the electron density in the emitter, we evaluate $`A1\times 10^{11}`$ meV<sup>-3/2</sup>s<sup>-1</sup> which compares well with the experimental result. In conclusion, we presented a measurement of the inelastic quasi-particle relaxation rate in a disordered Fermi liquid. This quantity was obtained from the analysis of the magnitude of disorder-induced fluctuations in the local density of states probed using the method of resonant tunneling through an impurity state. Quantitative comparison with the standard theory shows that, within the energy range available for such an analysis, the experimentally determined values of the inelastic relaxation rate can be attributed to the electron-electron interaction relaxation mechanism in diffusive conductors. We thank A. Förster and H. Lüth for growing the double-barrier heterostructure. We acknowledge financial support from BMBF, DFG, EPSRC, NATO, and TMR.
warning/0004/cond-mat0004307.html	ar5iv	text	# Weakly-interacting Bose-Einstein condensates under rotation ## I Introduction One of the basic questions about Bose-Einstein condensates in trapped alkali atom vapors is how they behave under rotation. A lot of theoretical work has been done on this subject, both analytical, and numerical , and the problem has been studied theoretically, both in the Thomas-Fermi limit of strong interactions and in the limit of weak interactions , which we consider in this paper. In Ref. Wilkin et al. considered a weakly interacting Bose gas with attractive interactions and showed that in the lowest energy state of a given angular momentum, the angular momentum is carried by the center of mass motion. Butts and Rokhsar calculated numerically the moment of inertia of a weakly-interacting trapped Bose gas with effective repulsive interactions . One of us identified the elementary modes of excitation for small angular momentum and demonstrated in Ref. that a system of rotating weakly-interacting bosons exhibits two additional kinds of condensation associated with the nature of low-lying excitations. Finally Bertsch and Papenbrock performed in Ref. exact numerical diagonalization within the subspace of states with a given angular momentum, which are degenerate in the absence of interactions. Experimentally the detection of vortex states in a two-component system has been reported by Matthews et al. , while Madison et al. have provided evidence for the formation of vortex states in a stirred one-component Bose-Einstein condensate. Our basic goal in this study is to identify the lowest energy states of a harmonically trapped, weakly interacting Bose gas for a given angular momentum $`L`$. As we show below these states are selected by the interactions. In Sec. II we describe the model and discuss the degeneracy of the many-body states for a given angular momentum in the absence of interactions. In Sec. III we use the mean-field approximation to calculate the interaction energy, and derive numerical and analytical results under various conditions. In Sec. IV we describe how one can go beyond the mean-field approximation and study as an example the specific case of small negative $`L/N1`$. Finally in Sec. V we give our conclusions. ## II The model Our starting point is the Hamiltonian $`H`$, given by $`H=H_0+V.`$ (1) Here $`H_0={\displaystyle \underset{i}{}}{\displaystyle \frac{\mathrm{}^2}{2M}}\mathbf{}_i^2+{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{2}}M\omega ^2(x_i^2+y_i^2)+f(z_i)`$ (2) includes the kinetic energy of the particles and their potential energy due to the trapping potential. The axis of rotation is taken to be the $`z`$ axis, and the trapping potential is assumed to be that of an isotropic harmonic oscillator of frequency $`\omega `$ in the $`x`$-$`y`$ plane. Also $`M`$ is the mass of the atoms. Our results do not depend on the trapping potential $`f(z)`$ in the $`z`$ direction. The interaction $`V`$ between atoms is assumed to be of zero range, $`V={\displaystyle \frac{1}{2}}U_0{\displaystyle \underset{ij}{}}\delta (𝐫_i𝐫_j),`$ (3) where $`U_0=4\pi \mathrm{}^2a/M`$ is the strength of the effective two-body interaction, with $`a`$ being the scattering length for atom-atom collisions. We assume that the interaction is repulsive, $`a>0`$. Attractive interactions have been studied in Refs. . Much theoretical work on rotating condensates has been done in the Thomas-Fermi limit of strong interactions, where the superfluid coherence length $`\xi =(8\pi na)^{1/2},`$ (4) $`n`$ being the particle density, is much less than the size of the cloud. Under these conditions the system is expected to exhibit superfluid properties much like those of liquid helium II . In this study we examine the opposite limit of weak interactions, $`nU_0\mathrm{}\omega `$ and $`nU_0\mathrm{\Delta }E_z`$, where $`\mathrm{\Delta }E_z`$ is the energy separation between the first excited state and the ground state for motion in the $`z`$ direction. Under the above conditions $`{\displaystyle \frac{\xi }{a_{\mathrm{osc}}}}\left({\displaystyle \frac{a_z}{Na}}\right)^{1/2},`$ (5) where $`N`$ is the number of atoms in the trap, $`a_{\mathrm{osc}}=(\mathrm{}/M\omega )^{1/2}`$ is the oscillator length, and $`a_z`$ is the characteristic length associated with the motion of the atoms along the $`z`$ axis. Therefore the coherence length is larger than the size of the cloud, and the situation is analogous to that for BCS pairing of nucleons in nuclei. Since we consider rotation around the $`z`$ axis, the condition $`nU_0\mathrm{\Delta }E_z`$ implies that the motion along this axis is frozen out and the problem is essentially two-dimensional. It is well known that for the harmonic oscillator potential in two dimensions the single-particle energies $`ϵ`$ are given in the absence of interactions by $`ϵ=(2n_r+\|m\|+1)\mathrm{}\omega ,`$ (6) where $`n_r`$ is the radial quantum number, and $`m`$ is the quantum number corresponding to the angular momentum. In the lowest energy state of the many boson system all particles are in states with $`n_r=0`$, and with $`m`$ being zero or having the same sign as the total angular momentum. The energy of the lowest state of a system of non-interacting bosons with angular momentum $`L`$ measured relative to that of the ground state is therefore $`E=L\mathrm{}\omega .`$ (7) There is a huge degeneracy corresponding to the many different ways of distributing $`L`$ quanta of angular momentum among $`N`$ atoms. Interactions between the atoms lift the degeneracy. We incorporate the effect of the interactions in both the mean-field approximation, as well as by diagonalization within some appropriately chosen truncated space of degenerate states. We describe the two methods separately below. ## III Mean-field approximation We start with the mean-field Gross-Pitaevskii approach. Butts and Rokhsar have used this method to derive numerical results for the moment of inertia of a Bose gas . In this scheme the many-body condensate wavefunction with $`N`$ particles and $`L`$ units of angular momentum $`\mathrm{\Psi }_{L,N}(𝐫_1,𝐫_2,\mathrm{},𝐫_N)`$ is taken to be the product of the single-particle states $`\mathrm{\Psi }(𝐫_i)`$, $`\mathrm{\Psi }_{L,N}(𝐫_1,𝐫_2,\mathrm{},𝐫_N)=\mathrm{\Psi }(𝐫_1)\times \mathrm{\Psi }(𝐫_2)\mathrm{}\mathrm{\Psi }(𝐫_N).`$ (8) The single-particle states $`\mathrm{\Psi }(𝐫_i)`$ can be expanded in terms of the harmonic-oscillator eigenstates $`\mathrm{\Phi }_m(𝐫_i)`$: $`\mathrm{\Psi }(𝐫_i)={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}c_m\mathrm{\Phi }_m(𝐫_i),`$ (9) where the $`c_m`$ are variational parameters, which are complex in general and are functions of $`L`$. The summation in Eq. (9) is restricted to positive $`m`$, since states with negative $`m`$ do not belong to the space of degenerate states. The quantity $`\|c_m\|^2`$ gives the occupation probability for state $`\mathrm{\Phi }_m`$. Also $$\mathrm{\Phi }_m(𝐫)=\frac{1}{(m!\pi a_{\mathrm{osc}}^2)^{1/2}}g(z)\left(\frac{\rho }{a_{\mathrm{osc}}}\right)^{\|m\|}e^{im\varphi }e^{\rho ^2/2a_{\mathrm{osc}}^2}.$$ (10) Here $`\rho ,z`$, and $`\varphi `$ are cylindrical polar coordinates. In the above expression we have assumed that the bosons are in their ground state $`g(z)`$ along the axis of rotation. The expectation value of the interaction energy $`V`$ in the state given by Eq. (8) is $`V={\displaystyle \frac{1}{2}}N(N1)U_0{\displaystyle \|\mathrm{\Psi }\|^4𝑑𝐫}.`$ (11) To find the lowest energy state we calculate $`V`$ as a function of the variational parameters $`c_m`$, and minimize it with respect to them under the following two constraints: the normalization condition, $`{\displaystyle \underset{m}{}}\|c_m^2\|=1,`$ (12) and the condition that the expectation value of the angular momentum per particle be fixed, $`{\displaystyle \underset{m}{}}m\|c_m\|^2=L/N.`$ (13) The parameters $`c_m`$ are complex in general, and therefore both their magnitudes, and their phases need to be determined. However Eqs. (12) and (13) impose two constraints on the magnitudes of the $`c_m`$. Furthermore, the overall phase of the wavefunction is arbitrary. Finally the rotational symmetry of the confining potential implies that the origin of the angular coordinate is arbitrary, which corresponds to the condition for conservation of angular momentum, which holds even in the presence of interactions. Therefore if the expansion (9) is truncated at a value $`m_{\mathrm{max}}`$, the number of independent variables is $`2\times (m_{\mathrm{max}}+1)4=2(m_{\mathrm{max}}1)`$. ### A Numerical results We have examined the problem numerically with $`m_{\mathrm{max}}`$ up to 9. The total number of terms in the expression for $`V`$ is 125 in this case. The result of such a calculation with $`m_{\mathrm{max}}=6`$ is shown in Fig. 1 for $`0L/N2`$. We show the results with $`m_{\mathrm{max}}=6`$, since the occupancy of states with higher $`m`$ is very low, and therefore including such states would not alter the results on this scale. Figure 2 shows the corresponding interaction energy. Also Fig. 3 shows the lines of constant density, $`\|\mathrm{\Psi }\|^2=`$ constant for $`L/N=0.1,0.6,0.8,`$ and 1.0. Figure 3 shows the gradual transition from mostly quadrupole and to a less extent octupole excitations, which are present at low angular momentum, to vortex-like structures as $`L`$ approaches $`N`$. We should also mention that the structures in Fig. 3, as well as those in Figs. 4, 5 and 6, rotate with an angular frequency $`\mathrm{\Omega }`$ given by $`\mathrm{\Omega }={\displaystyle \frac{1}{\mathrm{}}}{\displaystyle \frac{E_{\mathrm{tot}}}{L}}=\omega {\displaystyle \frac{1}{N\mathrm{}}}{\displaystyle \frac{V}{l}},`$ (14) which is lower than the trap frequency $`\omega `$. Here $`E_{\mathrm{tot}}`$ is the total energy of the system. When $`L`$ increases beyond $`N`$, the rotational invariance for $`L/N=1`$ is lost. Density contours for various values of $`L`$ between $`N`$ and $`2N`$ are shown in Fig. 4. These were calculated including the states up to $`m=6`$. For $`L1.75N`$ the optimal wavefunction has a two-fold axis of symmetry, and the odd-$`m`$ coefficients in the wavefunction vanish smoothly as the transition is approached, as shown in Fig. 1. In Figs. 5 and 6 we show contours for $`L=2N`$ and $`L=2.1N`$. There is a first-order transition from a state with two-fold symmetry to one with three-fold symmetry for $`L2.03N`$. The solution for $`L/N=2.0`$ with the states $`m=0`$, 2, 4, 6, and 8 considered has an energy of $`0.1757N^2v_0`$, whereas the one with the three-fold symmetry, with $`m=0`$, 3, 6, and 9, has $`0.1761N^2v_0`$. In contrast for $`L/N=2.1`$, the state with the three-fold symmetry has an energy $`0.1691N^2v_0`$, which is lower than the solution with the two-fold symmetry, with an energy $`0.1700N^2v_0`$. More generally, we have found that as $`L/N`$ increases, the lowest-energy states are the ones where a vortex array is formed, in agreement with the results of Ref.. ### B Analytical results for $`L/N`$ 0 We now turn to an analytic approach to the problem. One can systematically develop a power-series expansion for the occupancies $`\|c_m\|^2`$ of the states, as well as for the energy in certain limits. We start with the case of very low angular momentum, $`l=L/N0`$. Working in terms of quantum-mechanical states, one of us has calculated in Ref. the difference in the energies of two states where in the one state all $`N`$ particles have $`m=0`$ and in the other state a particle is promoted to the state with $`m=\lambda `$, and $`N1`$ particles have $`m=0`$. If we denote the states as $`\|0^{N_0},1^{N_1},2^{N_2},\mathrm{},`$ (15) where $`N_m`$ is the number of particles with angular momentum $`m\mathrm{}`$, the two states are $`\|0^N,`$ and $`\|0^{N1},\lambda ^1`$. The difference $`ϵ_\lambda `$ in the energy between these two states corresponding to this $`2^\lambda `$-pole excitation is given by $`ϵ_\lambda =\lambda \mathrm{}\omega \left(1{\displaystyle \frac{1}{2^{\lambda 1}}}\right)Nv_0+𝒪(v_0),`$ (16) where $`v_0=U_0\|\mathrm{\Phi }_0\|^4𝑑𝐫`$. One can easily see from Eq. (16) that at this level of approximation the excitations with the highest gain in interaction energy per unit of angular momentum are the ones with $`\lambda =2`$ or $`\lambda =3`$, i.e., quadrupole or octupole excitations. We now calculate the interaction energy for low values of $`l`$. The calculation of the energies of elementary excitations indicates that one would expect quadrupolar and octupolar modes to be the most important ones for small $`l`$. To determine the most energetically favorable way of giving the system angular momentum, one has to identify the behavior of $`\|c_2\|^2`$ and $`\|c_3\|^2`$ as $`l=L/N0`$, and then it is possible to build up a whole power-series expansion. Motivated by the fact that the $`\lambda =2`$ and 3 excitations are degenerate, and are the ones which give the highest gain in energy per unit of angular momentum, we assume that both $`c_2`$ and $`c_3`$ are of order $`l^{1/2}`$. As we show below, it is the mode-mode interaction that lifts this degeneracy, making the $`\lambda =2`$ mode dominant for low values of angular momentum. It is instructive to give an explicit example, so let us assume that we wish to examine the interaction energy up to order $`l^2`$. In order to minimize the interaction energy to this order, the states with $`m=1`$, 4, 5, and 6 need to be considered, since the phases of off-diagonal terms, like for example $`\|c_0\|\|c_1\|\|c_2\|\|c_3\|`$, can be chosen to have a negative sign, and thus lower the energy as compared to the case where only $`c_0`$, $`c_2`$, and $`c_3`$ are non-zero. A useful formula for the matrix elements of the potential is $`{\displaystyle \mathrm{\Phi }_k^{}(𝐫)\mathrm{\Phi }_l^{}(𝐫)\mathrm{\Phi }_m(𝐫)\mathrm{\Phi }_n(𝐫)𝑑𝐫}=`$ (17) $`\delta _{k+l,m+n}{\displaystyle \frac{(k+l)!}{2^{(k+l)}\sqrt{k!l!m!n!}}}{\displaystyle \|\mathrm{\Phi }_0(𝐫)\|^4𝑑𝐫}.`$ (18) As will become clear below, to calculate the energy up to order $`l^2`$ we must include the following terms: $`V=({\displaystyle \frac{1}{2}}\|c_0\|^4+{\displaystyle \frac{1}{2}}\|c_0\|^2\|c_2\|^2+{\displaystyle \frac{1}{4}}\|c_0\|^2\|c_3\|^2`$ (19) $`+{\displaystyle \frac{3}{16}}\|c_2\|^4+{\displaystyle \frac{5}{32}}\|c_3\|^4+{\displaystyle \frac{5}{8}}\|c_2\|^2\|c_3\|^2`$ (20) $`+\|c_0\|^2\|c_1\|^2{\displaystyle \frac{\sqrt{3}}{2}}\|c_0\|\|c_1\|\|c_2\|\|c_3\|`$ (21) $`+{\displaystyle \frac{1}{8}}\|c_0\|^2\|c_4\|^2{\displaystyle \frac{\sqrt{6}}{8}}\|c_0\|\|c_2\|^2\|c_4\|`$ (22) $`+{\displaystyle \frac{1}{16}}\|c_0\|^2\|c_5\|^2{\displaystyle \frac{\sqrt{10}}{8}}\|c_0\|\|c_2\|\|c_3\|\|c_5\|`$ (23) $`+{\displaystyle \frac{1}{32}}\|c_0\|^2\|c_6\|^2{\displaystyle \frac{\sqrt{5}}{16}}\|c_0\|\|c_3\|^2\|c_6\|)N^2v_0+𝒪(Nv_0).`$ (24) In the above expression we have chosen the phases $`\varphi _m`$ of the variational coefficients $`c_m`$ in such a way as to minimize $`V`$, and in the specific example we can arrange them so that all the off-diagonal matrix elements are negative. One of the phases can have any value, and we make the choice $`\varphi _0=0`$. The rest of them can be expressed in terms of, say, $`\varphi _1`$. We have found that up to $`m=6`$ \[$`m0]`$ the expression $`\varphi _m=m\varphi _1+(m+1)\pi `$ (25) gives the lowest energy. It is convenient to introduce the variable $`X=\|c_2\|^2+\|c_3\|^2`$, which is linear in $`l`$ to leading order, and make use of the constraints given by Eqs. (12) and (13) to get $`\|c_0\|^2`$ $`=`$ $`1X\|c_1\|^2\|c_4\|^2\|c_5\|^2\|c_6\|^2;`$ (26) $`\|c_2\|^2`$ $`=`$ $`3Xl+\|c_1\|^2+4\|c_4\|^2+5\|c_5\|^2+6\|c_6\|^2;`$ (27) $`\|c_3\|^2`$ $`=`$ $`l2X\|c_1\|^24\|c_4\|^25\|c_5\|^26\|c_6\|^2.`$ (28) Then Eq. (24) takes the form $`V=({\displaystyle \frac{1}{2}}{\displaystyle \frac{l}{4}}{\displaystyle \frac{31}{16}}X^2+{\displaystyle \frac{13}{8}}lX{\displaystyle \frac{9}{32}}l^2`$ (29) $`+{\displaystyle \frac{1}{4}}\|c_1\|^2{\displaystyle \frac{\sqrt{3}}{2}}\|c_1\|\|c_2\|\|c_3\|`$ (30) $`+{\displaystyle \frac{1}{8}}\|c_4\|^2{\displaystyle \frac{\sqrt{6}}{8}}\|c_2\|^2\|c_4\|`$ (31) $`+{\displaystyle \frac{5}{16}}\|c_5\|^2{\displaystyle \frac{\sqrt{10}}{8}}\|c_2\|\|c_3\|\|c_5\|`$ (32) $`+{\displaystyle \frac{17}{32}}\|c_6\|^2{\displaystyle \frac{\sqrt{5}}{16}}\|c_3\|^2\|c_6\|)N^2v_0+𝒪(Nv_0).`$ (33) The last four terms in the above equation can lower the energy to order $`l^2`$. For $`c_4`$, for example, the energy is minimized if \[see the third line of Eq. (33)\] $`{\displaystyle \frac{}{\|c_4\|}}{\displaystyle \frac{1}{8}}\|c_4\|^2={\displaystyle \frac{}{\|c_4\|}}{\displaystyle \frac{\sqrt{6}}{8}}\|c_2\|^2\|c_4\|,`$ (34) or $`\|c_4\|={\displaystyle \frac{\sqrt{6}}{2}}\|c_2\|^2l.`$ (35) Due to the non zero value of $`c_4`$ the energy is lowered by an amount $`\mathrm{\Delta }`$ $`=`$ $`\left({\displaystyle \frac{1}{8}}\|c_4\|^2{\displaystyle \frac{\sqrt{6}}{8}}\|c_2\|^2\|c_4\|\right)N^2v_0`$ (36) $`=`$ $`{\displaystyle \frac{3}{16}}\|c_2\|^4N^2v_0l^2N^2v_0.`$ (37) It is remarkable that the term $`\mathrm{\Delta }`$ exactly cancels the term $`3\|c_4\|^4/16`$ in the second line of Eq. (24). In a similar way $`c_1`$, $`c_5`$, and $`c_6`$ can be expressed in terms of $`c_2`$ and $`c_3`$ (and thus $`X`$), and Eq. (24) takes the form of the effective Hamiltonian $`V=({\displaystyle \frac{1}{2}}\|c_0\|^4+{\displaystyle \frac{1}{2}}\|c_0\|^2\|c_2\|^2+{\displaystyle \frac{1}{4}}\|c_0\|^2\|c_3\|^2`$ (38) $`+{\displaystyle \frac{5}{34}}\|c_3\|^4{\displaystyle \frac{1}{4}}\|c_2\|^2\|c_3\|^2)N^2v_0+𝒪(Nv_0),`$ (39) or equivalently $$V=\left[\frac{1}{2}\frac{l}{4}+\frac{27}{17}\left(X\frac{l}{2}\right)^2\right]N^2v_0+𝒪(Nv_0).$$ (40) Minimizing the above expression with respect to $`X`$ we find that $`X=l/2`$ and thus the angular momentum has to be carried by the $`m=2`$ state alone, since $`\|c_2\|^2=l/2`$ and $`\|c_3\|^2=0`$ up to terms linear in $`l`$. Also the quadratic correction to $`V`$ vanishes. Therefore for $`L/N0`$, the quadrupole ($`\lambda =2`$) excitations are dominant. This is one of the important conclusions of our study. We show in the Appendix that a diagrammatic perturbative expansion which assumes that only the states with $`m=0,2`$, and 3 are occupied by a macroscopic number of particles, while all the other states are not, \[but still contribute to the energy\] gives the same result. If one goes to higher order in $`l`$ the interaction energy has within the perturbative scheme a term of the form $`\|c_2\|^3\|c_3\|^2`$, which includes all the processes that convert three $`\lambda =2`$ excitations to two $`\lambda =3`$ excitations. This term can combine with the term $`\|c_3\|^4`$, which implies that it is possible that $`\|c_3\|^2\|c_2\|^3l^{3/2}`$, which actually turns out to be the case. Then $`c_1`$, for example, is given according to the second line of Eq. (33), by $`\|c_1\|=\sqrt{3}\|c_2\|\|c_3\|l^{5/4}.`$ (41) Using similar arguments we find that $`\|c_m\|^2l^{m/2}\mathrm{for}m1,\mathrm{and}\|c_1\|^2l^{5/2}.`$ (42) The leading terms in $`\|c_m\|^2`$ are given by $`\|c_0\|^2`$ $`=`$ $`1{\displaystyle \frac{1}{2}}l+{\displaystyle \frac{1}{3}}l^{3/2},`$ (43) $`\|c_1\|^2`$ $`=`$ $`l^{5/2}+2l^3,`$ (44) $`\|c_2\|^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}ll^{3/2},`$ (45) $`\|c_3\|^2`$ $`=`$ $`{\displaystyle \frac{2}{3}}l^{3/2},`$ (46) $`\|c_4\|^2`$ $`=`$ $`{\displaystyle \frac{3}{8}}l^2{\displaystyle \frac{3}{2}}l^{5/2}{\displaystyle \frac{1173}{816}}l^3,`$ (47) $`\|c_5\|^2`$ $`=`$ $`{\displaystyle \frac{2}{15}}l^{5/2}{\displaystyle \frac{4}{15}}l^3,`$ (48) $`\mathrm{and}\|c_6\|^2`$ $`=`$ $`{\displaystyle \frac{1}{144}}l^3,`$ (49) and the corresponding interaction energy is $`V=\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{l}{4}}+𝒪(l^4)\right]N^2v_0+𝒪(Nv_0).`$ (50) The above equation is another basic result of our study, namely that the interaction energy drops linearly with the angular momentum up to the order we have examined, for $`L/N0`$, in agreement with our numerical simulations and with those of Refs. . ### C Analytical results for $`l1`$ We now turn to the region $`L/N1`$. When the angular momentum per particle is exactly equal to 1, the lowest-energy state is the one where $`\|c_1\|^2=1`$, and corresponds to a vortex state. We consider the two cases $`l<1`$ and $`l>1`$ separately, starting with $`l<1`$. a. Analytical results for $`l<1`$ The simplest way to create a state with $`l<1`$ from that with $`l=1`$ is to transfer particles from the $`m=1`$ state to the ground state. However, the energy can be even lower if also the $`m=2`$ state is populated, because of the off-diagonal term $`\|c_0\|\|c_1\|^2\|c_2\|`$. The interaction energy up to order $`\overline{l}=1L/N`$ is found by minimizing the potential energy, retaining only the coefficients $`c_0`$, $`c_1`$, and $`c_2`$. This is $`V=({\displaystyle \frac{1}{4}}\|c_1\|^4+\|c_0\|^2\|c_1\|^2+{\displaystyle \frac{3}{4}}\|c_1\|^2\|c_2\|^2`$ (51) $`{\displaystyle \frac{\sqrt{2}}{2}}\|c_0\|\|c_1\|^2\|c_2\|)N^2v_0+𝒪(Nv_0),`$ (52) where we have used the fact that for this case too the phases may be shown to be given by Eq. (25). Equation (25) is valid for small negative $`L/N1`$ at least up to $`m=4`$. Thus in this limit $`\|c_0\|^2\|c_2\|^2\overline{l}.`$ (53) To obtain the coefficients of proportionality it is convenient to use the following parametrization: $`\|c_0\|^2`$ $`=`$ $`(1+\alpha )\overline{l},`$ (54) $`\|c_1\|^2`$ $`=`$ $`1(1+2\alpha )\overline{l},`$ (55) $`\|c_2\|^2`$ $`=`$ $`\alpha \overline{l},`$ (56) where $`\alpha `$ is a variational parameter. Minimizing the interaction energy in Eq. (52) with respect to $`\alpha `$ we find that $`\alpha =1`$. More generally using similar arguments we find that to leading order $`\|c_m\|^2\overline{l}^{\|m1\|},`$ (57) and the explicit expressions for the coefficients are $`\|c_0\|^2`$ $`=`$ $`2\overline{l}{\displaystyle \frac{3}{2}}\overline{l}^2,`$ (58) $`\|c_1\|^2`$ $`=`$ $`13\overline{l}+{\displaystyle \frac{27}{8}}\overline{l}^2,`$ (59) $`\|c_2\|^2`$ $`=`$ $`\overline{l}{\displaystyle \frac{9}{4}}\overline{l}^2,`$ (60) $`\|c_3\|^2`$ $`=`$ $`{\displaystyle \frac{3}{8}}\overline{l}^2,`$ (61) $`\mathrm{and}\|c_4\|^2`$ $`=`$ $`{\displaystyle \frac{\overline{l}^3}{12}},`$ (62) and the interaction energy to order $`\overline{l}^3`$ is, $`V=\left[{\displaystyle \frac{1}{4}}+{\displaystyle \frac{\overline{l}}{4}}+𝒪(\overline{l}^4)\right]N^2v_0+𝒪(Nv_0).`$ (63) The above equation implies that also in the region $`l<1`$ the interaction energy varies linearly with the angular momentum to the order we have examined, which is also in agreement with the numerical simulations. The coefficient of the linear term is the same as the one we found for small values of the angular momentum. Equations (50) and (63) as well as the numerical results \[see Fig. 2\] strongly suggest that the interaction energy $`V`$ drops linearly as a function of $`L/N`$ in the whole region $`0L/N1`$. The same result was derived by Butts and Rokhsar numerically within the mean-field approximation. In Ref. Bertsch and Papenbrock performed exact diagonalizations of degenerate states of a given $`L`$ and found that up to machine accuracy the energy of the lowest state for a given $`L`$ varies linearly with $`L`$ in the range $`2LN`$, in agreement with our analytic expansions. This result is specific to the contact form of the effective interaction, and is probably connected to a hidden symmetry in our Hamiltonian $`H`$, as discussed by Pitaevskii and Rosch , where the same Hamiltonian was considered in the context of breathing modes. b. Analytical results for $`l>1`$ We turn now to the case $`l>1`$. We calculate the difference in energy between the states $`\|1^N`$, and $`\|1^{N1},(\lambda +1)^1`$, by a method similar to that which for small $`l`$ led to Eqs. (15) and (16). The energy $`ϵ_\lambda `$ of this $`2^\lambda `$-pole excitation with $`L=N+\lambda `$, with $`\lambda N`$ is given by $`ϵ_\lambda =\lambda \mathrm{}\omega {\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{\lambda +2}{2^\lambda }}\right)Nv_0+𝒪(v_0).`$ (64) This formula implies that for $`l>1`$, the single-particle excitations with the lowest energy per unit of angular momentum are those with $`\lambda =4`$ or 5, which means that the actual angular momentum carried by the particles is $`m=5`$ or 6. In contrast to the low-angular momentum case, here both $`\|c_5\|^2`$ and $`\|c_6\|^2`$ vary linearly with $`\overline{l}`$, $`\|c_5\|^2\|c_6\|^2\overline{l}`$, where $`\overline{l}=L/N1`$. In addition, in this regime we find that the energy has corrections of higher order than linear. Using similar arguments to those given before for small $`l`$, we find, to order $`\overline{l}^2`$, $`\|c_0\|^2`$ $`=`$ $`0.1213\overline{l}^2,`$ (65) $`\|c_1\|^2`$ $`=`$ $`10.2241\overline{l},`$ (66) $`\|c_2\|^2`$ $`=`$ $`0.1934\overline{l}^2,`$ (67) $`\|c_5\|^2`$ $`=`$ $`0.1205\overline{l},`$ (68) $`\|c_6\|^2`$ $`=`$ $`0.1036\overline{l},`$ (69) $`\|c_9\|^2`$ $`=`$ $`1.7\times 10^3\overline{l}^2,`$ (70) $`\|c_{10}\|^2`$ $`=`$ $`1.3\times 10^3\overline{l}^2,`$ (71) $`\|c_{11}\|^2`$ $`=`$ $`8.6\times 10^3\overline{l}^2,`$ (72) and the interaction energy is $$V=\left[\frac{1}{4}\frac{5\overline{l}}{64}+6.7\times 10^3\overline{l}^2+𝒪(\overline{l}^3)\right]N^2v_0+𝒪(Nv_0).$$ (73) If we compare the expressions (63) and (73) for the interaction energy we see that there is a change in its slope, $`V/L`$, as $`L`$ passes $`N`$, from $`Nv_0/4`$ for $`L<N`$ to $`5Nv_0/64`$ for $`L>N`$. ### D Results for higher values of $`L/N`$ We mentioned earlier that as the angular momentum per particle increases even further, there are certain ranges of values of $`L/N`$ over which the state with the lowest energy has a specific symmetry. The lowest value of $`L/N`$ for which this occurs is $`1.75`$ and the symmetry of the state is two-fold, i.e., only $`c_{2m}0`$. We have examined analytically as an example the case $`L/N=2`$. Keeping only the first three non-zero coefficients, which are the dominant ones, we find to order $`\overline{l}=L/N2`$ that $`\|c_0\|^2`$ $`=`$ $`\beta (\overline{l})\overline{l}/4,`$ (74) $`\|c_2\|^2`$ $`=`$ $`12\beta (\overline{l}),`$ (75) $`\mathrm{and}\|c_4\|^2`$ $`=`$ $`\beta (\overline{l})+\overline{l}/4,`$ (76) where $`\beta (\overline{l})`$ $`=`$ $`{\displaystyle \frac{309248\sqrt{6}}{12695}}+{\displaystyle \frac{17(64\sqrt{6}+109)}{10156}}\overline{l}`$ (77) $``$ $`0.2343+0.4449\overline{l},`$ (78) or $`\|c_0\|^2`$ $``$ $`0.2343+0.1949\overline{l},`$ (79) $`\|c_2\|^2`$ $``$ $`0.53140.8897\overline{l},`$ (80) $`\|c_4\|^2`$ $``$ $`0.2343+0.6949\overline{l},`$ (81) and the corresponding interaction energy is: $`V=\left[AB\overline{l}+C\overline{l}^2+𝒪(\overline{l}^3)\right]N^2v_0+𝒪(Nv_0),`$ (82) where $`A`$ $`=`$ $`{\displaystyle \frac{3}{16}}+{\displaystyle \frac{\beta _0}{32}}(74\sqrt{6})+{\displaystyle \frac{\beta _0^2}{256}}(64\sqrt{6}109),`$ (83) $`B`$ $`=`$ $`{\displaystyle \frac{1}{512}}(4+85\beta _0),`$ (84) and $`\beta _0=\beta (0)`$. The actual numbers which appear in Eq. (82) are $`V\left[0.17730.0467\overline{l}+0.0170\overline{l}^2+𝒪(\overline{l}^3)\right]N^2v_0`$ (85) $`+𝒪(Nv_0).`$ (86) There is no change in the slope of the interaction energy as $`L`$ passes $`2N`$. Figure 5 shows lines of constant density, $`\|\mathrm{\Psi }\|^2=\mathrm{constant}`$, for $`L/N=2`$. The occurence of two nodes in the density reflects the presence of two displaced vortices, and thus we see that the lowest-energy state of the system has two separated vortices and not a doubly quantized vortex. This clearly demonstrates the instability of the double-quantized vortex state to formation of two vortices plus surface waves. ## IV Beyond the mean-field approximation Another way of approaching the problem of rotation, is to diagonalize the Hamiltonian within the space of degenerate states. This approach goes beyond the mean-field approximation, since in mean field theory the many-body wavefunction is the product of the single-particle states, whereas the diagonalization allows for the many-body state to have all kinds of correlations between the particles. This technique can be used by taking into account the whole set of states , but it is convenient and pedagogical to work in a restricted space, appropriately chosen. As an example we consider small negative $`L/N1`$. From the analysis of Sec. III C we know that in this limit the states with $`m=1,0`$, and 2 are dominant. Therefore the eigenstates $`\|\mu ,\stackrel{~}{l}=\|0^{\stackrel{~}{l}+\mu },1^{N\stackrel{~}{l}2\mu },2^\mu `$ (87) with $`N`$ particles and $`L=N\stackrel{~}{l}`$ units of angular momentum are expected to provide a good basis for describing the low-lying states for $`\stackrel{~}{l}N`$. We restrict ourselves to this limit and demonstrate how one can derive an effective Hamiltonian which can be diagonalized exactly. In the limit we consider, $`\stackrel{~}{l}`$ is $`N`$, and thus $`\mu \stackrel{~}{l}N`$. The diagonal matrix elements in the Hamiltonian are, up to terms of order $`N`$, $`\mu \|V\|\mu =\left({\displaystyle \frac{1}{4}}N(N1)+{\displaystyle \frac{1}{2}}\stackrel{~}{l}N+{\displaystyle \frac{3}{4}}\mu N\right)v_0,`$ (88) and the off-diagonal matrix elements are $`\mu +1\|V\|\mu {\displaystyle \frac{\sqrt{2}}{4}}Nv_0\sqrt{(\mu +\stackrel{~}{l}+1)(\mu +1)}.`$ (89) Ignoring for the moment the (diagonal) first term of Eq. (88), which corresponds to the interaction energy of the state $`\|1^N`$, we see from Eqs. (88) and (89) that we have to diagonalize the Hamiltonian $`\stackrel{~}{H}=\left[{\displaystyle \frac{1}{2}}a_0^{}a_0+{\displaystyle \frac{1}{4}}a_2^{}a_2+{\displaystyle \frac{\sqrt{2}}{4}}(a_2^{}a_0^{}+a_2a_0)\right]Nv_0,`$ (90) which can be done exactly by use of a Bogoliubov transformation. Here $`a_m`$ is an annihilation operator that destroys a particle with angular momentum $`m\mathrm{}`$. Introducing the operators $`c`$ and $`d`$ given by $`c=a_0^{}+\sqrt{2}a_2;d=\sqrt{2}a_0+a_2^{},`$ (91) we may write the Hamiltonian as $`\stackrel{~}{H}={\displaystyle \frac{N}{4}}(d^{}d1)v_0.`$ (92) Acting on states of the type (87) the operator $`d^{}dc^{}c`$ is diagonal, and has an eigenvalue $`NL`$. Therefore $`d^{}d`$ can be eliminated and from Eq. (92) we obtain $`\stackrel{~}{H}={\displaystyle \frac{N}{4}}(NL1)v_0+{\displaystyle \frac{Nv_0}{4}}c^{}c.`$ (93) The total interaction energy is thus the eigenenergy of $`\stackrel{~}{H}`$ plus the diagonal part $`N(N1)v_0/4`$, or $`V`$ $`=`$ $`{\displaystyle \frac{N}{4}}(N1)v_0+{\displaystyle \frac{N}{4}}(NL1+c^{}c)v_0+𝒪(v_0)`$ (94) $`=`$ $`{\displaystyle \frac{N(2NL2)}{4}}v_0+{\displaystyle \frac{1}{4}}c^{}cNv_0+𝒪(v_0).`$ (95) In the ground state $`c^{}c=0`$, and the energy given by Eq. (95) is the same as that derived numerically in Ref. . The presence of the term $`c^{}c`$ in Eq. (95) implies that the excited states in this limit of small negative $`L/N1`$ are separated from the ground state by an amount $`\mathrm{\Delta }E={\displaystyle \frac{N}{4}}v_0+𝒪(v_0).`$ (96) We have also performed numerical diagonalization, and we have confirmed the above result (95), as well as Eq. (96). The average occupancy of the states of the non-interacting problem is, for the lowest-energy state and $`L=N`$: $$\|c_1\|^2=1\frac{2}{N}+𝒪\left(\frac{1}{N^2}\right);\|c_0\|^2=\|c_2\|^2=\frac{1}{N}+𝒪\left(\frac{1}{N^2}\right),$$ (97) and thus in the limit $`N\mathrm{}`$ there is agreement between the mean-field approximation and the present one. ## V Summary and conclusions To summarize, we have studied the lowest-energy states of a system of rotating, weakly interacting harmonically trapped bosons. Within the mean-field approximation, for $`L/N0`$ we find that the angular momentum is carried mainly by quadrupole $`(\|m\|=2)`$ excitations. We have demonstrated that diagrammatic perturbation theory also leads to the same results as the method we have used here. For $`L/N=1`$ the angular momentum is carried by particles in the $`m=1`$ state, while for small negative $`L/N1`$ the $`m=0`$ and $`m=2`$ states are also populated. In the limits $`L/N0`$ and $`L/N1`$ the energy is a linear function of the angular momentum up to the order we have explored, while numerically this linearity persists in the whole region $`0L/N1`$. This result is specific to the contact form of the effective interaction, and does not hold for more general interactions. For small positive $`L/N1`$, the states which carry the additional angular momentum are those with $`m=5`$ and $`m=6`$. In addition, as $`L`$ passes $`N`$ the derivative of the interaction energy with respect to the angular momentum changes abruptly. We have also found that for $`L/N1.75`$ there is a second order phase transition and for $`1.75L/N2.03`$ the lowest energy state has two-fold symmetry. At $`L/N2.03`$, there is a first order phase transition to a state with three-fold symmetry. More generally, for higher values of $`L/N`$ a vortex array develops. The Gross-Pitaevskii wavefunction is a power series in $`\stackrel{~}{z}=x+iy`$. Thus if one truncates the series at $`m=m_{\mathrm{max}}`$, the wavefunction will have $`m_{\mathrm{max}}`$ nodes. In the vicinity of a node at $`\stackrel{~}{z}=\stackrel{~}{z}_0`$ the wavefunction varies as $`\stackrel{~}{z}\stackrel{~}{z}_0`$ and therefore each node corresponds to a singly quantized vortex having the same sense as the total angular momentum. It is instructive to study how the vortex lines move as the angular momentum is increased. For low angular momentum, the condensate wavefunction has only $`m=0`$ and $`m=2`$ components, and it is therefore proportional to $`[1(l^{1/2}/2)(\rho /a_{\mathrm{osc}})^2e^{2i\varphi }]\mathrm{exp}(\rho ^2/2a_{\mathrm{osc}}^2)`$ for the choice of $`\varphi _2=\pi `$, according to Eq. (25) with $`\varphi _1=0`$. This has vortices on the $`x`$ axis at $`x=\pm (2/l^{1/2})^{1/2}a_{\mathrm{osc}}`$. With increasing angular momentum, components of the wavefunction with odd $`m`$ grow, and the two-fold symmetry of the cloud is broken, as may be seen in Fig. 1 for $`L/N=0.1`$, one of the vortices moving to larger distances, and the other to smaller ones. The $`c_3`$ term leads to a third vortex at large distances from the origin. For $`L/N=1`$ there is only one vortex, which is at the origin. With further increase in $`L/N`$, the velocity field is at first still dominated by a vortex close to the origin, but subsequently a second vortex moves into the cloud until at $`L/N1.75`$ the two-fold symmetry is restored. As $`L/N`$ increases towards the value 2.03, at which the first-order transition to the state with three-fold symmetry mentioned above occurs, the separation of the two vortices changes little. In this paper we have also investigated effects not included in mean-field theory by diagonalizing a model Hamiltonian for $`L`$ close to, but less from, $`N`$. We find that for $`L=N`$ the occupancy of the $`m=1`$ state is 1, with corrections of order $`1/N`$. We have calculated the energy up to terms of order $`N`$. Finally we also found that the low-lying excited states are separated from the lowest state by energies of order $`Nv_0`$ of the same angular momentum. In this study we have examined the limit of weak interactions. When the interaction energy per particle $`nU_0`$ becomes comparable to or greater than $`\mathrm{}\omega `$, components of the wavefunction that are not members of the lowest multiplet in the absence of interaction must be included. Calculations for this regime based on the Gross-Pitaevskii equation have been carried out by Isoshima and Machida . Comparison of our results with theirs is difficult because these authors calculated the lowest energy state in a rotating frame, rather than the lowest energy state for a given angular momentum. One question of importance both conceptually and because of its relevance to experiment is whether or not the states are stable to small perturbations, and if they are not, what is the lifetime of the state. The answer to these questions depends on the nature of the perturbation, whether it is due to a deformation of the trap, or to interactions with particles outside the condensate, and we shall discuss it elsewhere. In our calculations above we have shown for a particular example that the Gross-Pitaevskii approach gives correctly the contribution to the energy of order $`N^2`$. This result, which is alluded to in Ref. , is more general, and in a future publication it will be shown how the Gross-Pitaevskii approach is recovered as the first term in an expansion in powers of $`1/N`$. The method may be extended to calculate contributions to the energy of order $`N`$ which are in excellent agreement with results obtained by numerical diagonalization of the Hamiltonian. ## Perturbation theory approach We show in this appendix that one can use perturbation theory to derive an effective Hamiltonian in the region $`L/N0`$, corresponding to Eq. (39). We assume that only the states with $`m=0,2`$, and 3 are macroscopically occupied. However, other states (the $`m=1,4,5`$ and 6 in this case) give corrections to the energy that can be treated perturbatively. Let us demonstrate how this works by considering the interaction energy up to $`l^2`$. As long as both $`c_2`$ and $`c_3`$ vary as $`l^{1/2}`$, the only processes that contribute to the interaction energy up to $`l^2`$ are shown in Fig. 7. Let us consider the first process on the left as an example. The matrix element $`M`$ corresponding to the vertex, where two particles with $`m=2`$ scatter to states with $`m=0`$ and $`m=4`$, is equal to $`M={\displaystyle \frac{\sqrt{6}}{16}}N_2\sqrt{N}v_0.`$ (98) From Eq. (16) we find that the difference in the energy between the intermediate state and the initial state is $`\delta ϵ={\displaystyle \frac{Nv_0}{8}}+𝒪(v_0),`$ (99) and thus perturbation theory implies that the correction to the energy is $`{\displaystyle \frac{\|M\|^2}{\delta ϵ}}={\displaystyle \frac{3}{16}}N_2^2v_0={\displaystyle \frac{3}{16}}\|c_2\|^4N^2v_0,`$ (100) which is precisely the correction $`\mathrm{\Delta }`$ given by Eq. (37) (plus terms of order $`Nv_0`$). Similarly the other diagrams shown below give $`5\|c_3\|^4/544`$, $`\|c_2\|^2\|c_3\|^3/8`$, and $`3\|c_2\|^2\|c_3\|^2/4`$ in units of $`N^2v_0`$, respectively, and are identical to the corrections given by the terms in the last four lines of Eq. (24). G.M.K. was supported by the European Commission, TMR program, contract No. ERBFMBICT 983142. Helpful discussions with A. Jackson and S. Reimann are gratefully acknowledged. G.M.K. would like to thank the Foundation of Research and Technology, Hellas (FORTH) for its hospitality.
warning/0004/hep-ph0004264.html	ar5iv	text	# IRB-TH-2/00April 2000 Lifetime-difference pattern of heavy hadrons ( Theoretical Physics Division, Rudjer Bošković Institute, P.O.Box 180, HR-10002 Zagreb, Croatia) PACS: 14.20.Lq; 13.20.Mr; 14.40.Lb; 14.40.Nd Keywords: Heavy hadrons; Lifetimes; Inclusive decays; Four-quark operators; Isospin symmetry; Heavy-quark symmetry It is astonishing that a lot of physical observables, decay rates of semileptonic, nonleptonic, and radiative decays are described in the framework of the inverse heavy-quark mass expansion in terms of relatively few basic quantities, e.g., quark masses and hadronic expectation values ($`HEV`$) of several leading local operators . The underlying theory is based on a few field-theory basics, such as the operator product expansion, quark-hadron duality, and certain well-known symmetries. Even more surprisingly, the theory seems to work rather well in the case of charmed hadrons, where the expansion parameter $`\sqrt{\mu _G^2(D)/m_c^2}0.5`$ is by no means very small. A systematic analysis leads to very clear lifetime-pattern predictions, the agreement with experiment being reasonable even for absolute lifetime values . In beauty hadron decays, one expects that the whole theory should work much better, since the expansion parameter, $`\sqrt{\mu _G^2(B)/m_b^2}0.131`$, is significantly smaller than in the case of charmed hadrons. Still, in spite of the overall agreement between theory and experiment , there are some questions to be answered. The main question is: does the quark-hadron duality work? Unfortunately, for a reliable study of duality, one needs complete control over nonperturbative phenomena, which is out of reach of the present theory. Nevertheless, the validity of duality was studied in exactly solvable ’t Hooft model in $`QCD_2`$ , leading to a perfect matching, which is certainly encouraging. Still, one would be happy to find stronger support for duality in the $`(3+1)`$ theory, too. The question is: could one try to reduce the uncertanties of the calculation, like the strong sensibility to the value of $`m_Q`$, and/or find the way to disentangle various preasymptotic effects which cause lifetime differences? The decay rate of a decaying hadron is generally of the type $$\mathrm{\Gamma }(H_Qf)=\frac{G_\mathrm{F}^2m_Q^5}{192\pi ^3}\|V\|^2\frac{1}{2M_{H_Q}}\left[\underset{D=3}{\overset{D_c}{}}c_D^f\frac{H_Q\|O_D\|H_Q}{m_Q^{D3}}+𝒪(1/m_Q^{D_c2})\right],$$ (1) where $`c_D^f`$ are the Wilson coefficients and $`H_Q\|O_D\|H_Q`$ are the matrix elements of the $`D`$-dimensional operators which are suppressed by the inverse power of mass $`1/m_Q^{D3}`$. The following comments are in order: 1. It would be desirable to extract a combination of rates $`\mathrm{\Gamma }^{}`$ that depends only on the operators with $`D=6`$, since they are the operators responsible for lifetime differences. 2. Since the $`D=6`$ operators are suppressed by $`m_Q^3`$ with respect to the leading $`D=3`$ operator, this would reduce the mass dependence of $`\mathrm{\Gamma }^{}`$ to $`m_Q^2`$, thus significantly reducing the errors coming from the uncertainty of $`m_Q`$. 3. Even truncating the sum in (1) at $`D=6`$, one would presumably be able to test the role of preasymptotic effects with satisfactory accuracy, relying on heavy-quark symmetry. That is the aim of this paper. As both the decays in the $`b`$ and $`c`$ sectors of hadrons are described by the same formalism, it seems worthwhile to investigate the possibility of establishing connections between these two sectors. In a recent paper , Voloshin demonstrated one of such possible connections. The strength of that approach is in avoiding the model dependence assuming $`SU(3)`$ flavor and heavy-quark symmetry (HQS). One can extend this approach to obtain more predictivity. Using the matrix elements extracted from charmed baryons, a parameter $`F_B^{eff}`$ (parametrizing four-quark contributions to beauty baryons) can be determined and concrete numerical predictions can be obtained in this way . This procedure brings the ratio $`\tau (\mathrm{\Lambda }_b)/\tau (B_d^0)`$ to better agreement with experiment and gives predictions for the lifetimes of beauty baryons. In the present paper we adopt a similar, yet different strategy from that followed in . In , the author, when establishing the connection between $`c`$ and $`b`$ baryons, always relates baryons with the same light-quark content, i.e., mutually related baryons differ in the flavor of the heavy quark. Also, the author explicitly extracts values of four-quark operator matrix elements and then applies them elsewhere. Our approach differs in both these respects. The heavy hadrons that we relate neither contain the same heavy quark(s) nor the same light (anti)quarks – they are affected by the same type of the four-quark operator contribution. Furthermore, we form such combinations that the four-quark operator matrix elements get reduced. In forming such combinations, we consider the leading modes of the decay of $`c`$ and $`b`$ quarks with respect to the CKM matrix elements. For the $`c`$ quark, we then have only one nonleptonic mode, $`cs\overline{d}u`$, and one semileptonic mode per lepton family, $`cs\overline{l}\nu _l`$. In the case of $`b`$ quark, there are two nonleptonic modes, $`bc\overline{u}d`$ and $`bc\overline{c}s`$, as well as one semileptonic mode per lepton family, $`bcl\overline{\nu _l}`$. As we consider $`b`$ baryons containing only light quarks along with $`b`$ quarks, the semileptonic Cabibbo leading modes for the decay of the $`b`$ quark do not appear. We also disregard mass corrections of four-quark operators coming from the massive particles in final decay states. This is a better approximation for $`c`$ decays (as $`m_s^2/m_c^20.01`$) than for $`b`$ decays (where $`m_c^2/m_b^20.1`$). Traditionally, the effects of four-quark operators are called the positive Pauli interference, the negative Pauli interference and, $`W`$ exchange or annihilation . We consider pairs of one charmed and one beauty particle which contain the same type of four-quark contribution. It is important to notice that in $`b`$ and $`c`$ decays different light-quark flavors participate in the same type of four-quark contributions. Using isospin ($`SU(2)`$ flavor) symmetry and HQS we are able to form appropriate decay-rate differences which lead us to the final result. We start our considerations with heavy mesons. The first pair of considered particles are $`D^+(c\overline{d})`$ and $`B^{}(b\overline{u})`$. In both these mesons the effect of negative Pauli interference occurs. The contributions of four-quark operators to the decay rates of these two particles are described by the same type of operators. Moreover, the operators for the $`b`$ case can be obtained from those in the $`c`$ case by making the substitution $`cb`$, $`\overline{d}\overline{u}`$. The second pair of particles form $`D^0(c\overline{u})`$ and $`B^0(b\overline{d})`$. In both these mesons the effect of $`W`$ exchange occurs. Next, we would like to isolate the effects of four-quark operators, i.e., we have to eliminate the contributions of the operators of dimensions $`5`$ and lower. We achieve this goal by taking differences of the decay rates of $`D`$ and $`B`$ mesons assuming the isospin symmetry. This procedure leads us to the following relations: $$\mathrm{\Gamma }(D^+)\mathrm{\Gamma }(D^0)=\frac{G_F^2m_c^2}{4\pi }\|V_{cs}\|^2\|V_{ud}\|^2[D^+\|P^{cd}\|D^+D^0\|P^{cu}\|D^0],$$ (2) $$\mathrm{\Gamma }(B^{})\mathrm{\Gamma }(B^0)=\frac{G_F^2m_b^2}{4\pi }\|V_{cb}\|^2\|V_{ud}\|^2[B^{}\|P^{bu}\|B^{}B^0\|P^{bd}\|B^0],$$ (3) where $`P`$’s denote the appropriate four-quark terms from the heavy-quark effective Lagrangian, see, e.g., Ref. . Using HQS and isospin symmetry we can express the terms describing the negative Pauli interference as $$P^{cd}=\underset{i=1}{\overset{2}{}}c_i^{cd}(\mu )O_i^{cd}(\mu )=\underset{i=1}{\overset{2}{}}c_{i,negint}^{Qq}(\mu )O_{i,negint}^{Qq}(\mu )+𝒪(1/m_c),$$ (4) $$P^{bu}=\underset{i=1}{\overset{2}{}}c_i^{bu}(\mu )O_i^{bu}(\mu )=\underset{i=1}{\overset{2}{}}c_{i,negint}^{Qq}(\mu )O_{i,negint}^{Qq}(\mu )+𝒪(1/m_b),$$ (5) where $`Q`$ denotes a heavy quark in the heavy-quark limit $`m_Q\mathrm{}`$. For the matrix elements of the above mentioned operators we obtain (by applying HQS and isospin symmetry again) $`D^+\|P^{cd}\|D^+=A_{negint}+𝒪^{}(1/m_c)`$ and $`B^{}\|P^{bu}\|B^{}=A_{negint}+𝒪^{}(1/m_b)`$, where the matrix element is $`A_{negint}=M_{Q\overline{q}}\|_{i=1}^2c_{i,negint}^{Qq}(\mu )O_{i,negint}^{Qq}(\mu )\|M_{Q\overline{q}}`$ and $`M_{Q\overline{q}}`$ denotes the mesonic state with one heavy quark $`Q`$ and one light antiquark $`\overline{q}`$ in the heavy-quark limit. In the preceding relations we have also stated that heavy-quark mass suppressed corrections to the operators and their matrix elements do not have to be the same. In an analogous manner we can obtain the expressions for the operators describing $`W`$ exchange using the notation explained above: $$P^{cu}=\underset{i=1}{\overset{4}{}}c_i^{cu}(\mu )O_i^{cu}(\mu )=\underset{i=1}{\overset{4}{}}c_{i,exch}^{Qq}(\mu )O_{i,exch}^{Qq}(\mu )+𝒪(1/m_c),$$ (6) $$P^{bd}=\underset{i=1}{\overset{4}{}}c_i^{bd}(\mu )O_i^{bd}(\mu )=\underset{i=1}{\overset{4}{}}c_{i,exch}^{Qq}(\mu )O_{i,exch}^{Qq}(\mu )+𝒪(1/m_b).$$ (7) The matrix elements of the operators given by (6) and (7) are $`D^0\|P^{cu}\|D^0=A_{exch}+𝒪^{}(1/m_c)`$ and $`B^0\|P^{bd}\|B^0=A_{exch}+𝒪^{}(1/m_b)`$, with $`A_{exch}=M_{Q\overline{q}}\|_{i=1}^4c_{i,exch}^{Qq}(\mu )O_{i,exch}^{Qq}(\mu )\|M_{Q\overline{q}}`$. In all relations given above, the dependence of the coefficients $`c_i`$ and the operators $`O_i`$ on the scale $`\mu `$ is explicitly displayed. Since the dependences of the operators and the coefficients on $`\mu `$ cancel each other, there is no need for an explicit value for the $`\mu `$ scale, and also no hybrid renomalization is required. That scale, however, should be small compared with the heavy-quark masses, to allow for the proper heavy-quark expansion and the reduction of the four-quark operator contribution in the heavy-quark limit, as shown below. Combining all preceding expressions, we obtain for the decay rate differences $$\mathrm{\Gamma }(D^+)\mathrm{\Gamma }(D^0)=\frac{G_F^2m_c^2}{4\pi }\|V_{cs}\|^2\|V_{ud}\|^2[A_{negint}A_{exch}+𝒪(1/m_c)],$$ (8) $$\mathrm{\Gamma }(B^{})\mathrm{\Gamma }(B^0)=\frac{G_F^2m_b^2}{4\pi }\|V_{cb}\|^2\|V_{ud}\|^2[A_{negint}A_{exch}+𝒪(1/m_b)].$$ (9) As $`m_bm_c`$, it is clear that the approximation made above should work much better for beauty particles than for charmed particles. Still, we assume that the approximations made are justified in both $`c`$ and $`b`$ decays and that the corrections of order $`𝒪(1/m_{c,b})`$ are not large. In this case, we obtain the final relation for mesons: $$r^{BD}\frac{\mathrm{\Gamma }(B^{})\mathrm{\Gamma }(B^0)}{\mathrm{\Gamma }(D^+)\mathrm{\Gamma }(D^0)}=\frac{m_b^2}{m_c^2}\frac{\|V_{cb}\|^2}{\|V_{cs}\|^2}+𝒪(1/m_c,1/m_b).$$ (10) Let us consider relation (10) in more detail. By forming the lifetime differences (2) and (3) we have not only eliminated the effects of the leading operators in the mesonic decay rates, but applied the HQS first at the subleading level of $`𝒪(1/m_{c,b}^3)`$. This has in turn enabled us to eliminate the dependence on four-quark matrix elements in relation (10). Also, the sensitivity to the choice of heavy quark masses is now significantly reduced since in expression (10) we have only the second power of masses. The procedure, however, has its limits. It enables us to reduce the matrix elements, but cannot be extended to CKM suppressed modes. Also, mass corrections due to massive particles in the final decay states spoil the procedure and, therefore, have been left out. Nevertheless, both of these corrections are under good theoretical control. Moreover, the contributions from mass corrections and suppressed modes are smaller by more than one order of magnitude than the leading contribution and cannot significantly change relation (10). Starting from expression (10), we may check the standard formalism of inclusive decays against experimental data, especially if the four-quark operator contributions are sufficient to explain the lifetime differences of heavy mesons. To perform such a check, we determine the quantity $`r^{BD}`$ using the experimental values for lifetimes and theoretically (using relation (10)) and then compare the values. Taking the experimental values from , we obtain $$r_{exp}^{BD}=0.030\pm 0.011.$$ (11) The large error of $`r_{exp}^{BD}`$ comes from the fact that the experimental errors of lifetime measurements of $`B^{}`$ and $`B^0`$ are large compared with the difference between central values of the results for $`B^{}`$ and $`B^0`$. Inspection of shows that the experimental situation in the sector of $`B`$ mesons is far from being settled and more precise measurements of $`B`$ meson lifetimes are needed to reduce the uncertainty in the value of $`r_{exp}^{BD}`$. The theoretical value $`r_{th}^{BD}`$ is calculated from expression (10). The numerical values for heavy-quark masses are taken to be $`m_c(m_c)=1.25\pm 0.1GeV`$ and $`m_b(1GeV)=4.59\pm 0.08GeV`$ for the reasons discussed in . The values for the matrix elements of the CKM matrix are taken from to be $`\|V_{cs}\|=1.04\pm 0.16`$ and $`\|V_{cb}\|=0.0402\pm 0.0019`$. Using these numerical values we obtain $$r_{th}^{BD}=0.020\pm 0.007.$$ (12) The relatively large error in (12) comes predominantly from the large experimental error of the $`\|V_{cs}\|`$ CKM matrix element. Direct comparison of the numerical results (11) and (12) indicates that these are consistent within errors. Since the relative errors of both results are rather large and experimental data are still fluid, it is possible that these results will experience some change. Nevertheless, we expect that the two results will remain comparable and even closer to each other. These results confirm, in a model-independent way, that four-quark operators can account for the greatest part of decay rate differences of heavy mesons, i.e., that the contributions of other operators of dimension $`6`$ and higher-dimensional operators are not so important. Of course, these conclusions should be interpreted in the light of the approximations made. The procedure presented for heavy mesons can be extended to heavy baryons. First, we consider the system of singly heavy baryons. These baryons contain two light quarks, which introduces two types of four-quark operator contributions for each baryon. Again, we choose two pairs of singly heavy baryons, each pair containing one charmed and one beauty baryon related by the same dominant type of four-quark contributions. The first pair contains $`\mathrm{\Xi }_c^+(cus)`$ and $`\mathrm{\Xi }_b^{}(bds)`$, while the second pair contains $`\mathrm{\Xi }_c^0(cds)`$ and $`\mathrm{\Xi }_b^0(bus)`$. The first pair exhibits the negative Pauli interference and different nonleptonic and semileptonic contributions of four-quark operators containing $`s`$ quark fields. The second pair comprises the $`W`$-exchange effect and the same contributions of four-quark operators containing $`s`$ quark fields as in the first pair of baryons. It is important to notice that the four-quark operators that contain $`s`$ quark fields do not describe the same type of effects in $`c`$ and $`b`$ decays in our pairs (and therefore are not given by the contributions of the same type in the effective Lagrangian). Nevertheless, we form decay-rate differences in such a manner that the contributions of four-quark operators containing the $`s`$-quark field cancel (assuming isospin symmetry). This cancellation is important in the case of $`b`$ decays because in the contributions of $`s`$-quark type four-quark operators significant mass corrections appear because of two massive $`c`$ quarks in the final state. In the singly heavy sector we, therefore, form the following differences: $$\mathrm{\Gamma }(\mathrm{\Xi }_c^+)\mathrm{\Gamma }(\mathrm{\Xi }_c^0)=\frac{G_F^2m_c^2}{4\pi }\|V_{cs}\|^2\|V_{ud}\|^2[\mathrm{\Xi }_c^+\|P^{cu}\|\mathrm{\Xi }_c^+\mathrm{\Xi }_c^0\|P^{cd}\|\mathrm{\Xi }_c^0],$$ (13) $$\mathrm{\Gamma }(\mathrm{\Xi }_b^{})\mathrm{\Gamma }(\mathrm{\Xi }_b^0)=\frac{G_F^2m_b^2}{4\pi }\|V_{cb}\|^2\|V_{ud}\|^2[\mathrm{\Xi }_b^{}\|P^{bd}\|\mathrm{\Xi }_b^{}\mathrm{\Xi }_b^0\|P^{bu}\|\mathrm{\Xi }_b^0].$$ (14) Following the procedure displayed in relations (4) to (10) and using the isospin symmetry and HQS again, we obtain the expressions (note that the four-quark operators which contribute to the negative Pauli interference in mesons contribute to $`W`$ exchange in baryons and vice versa ): $$\mathrm{\Gamma }(\mathrm{\Xi }_c^+)\mathrm{\Gamma }(\mathrm{\Xi }_c^0)=\frac{G_F^2m_c^2}{4\pi }\|V_{cs}\|^2\|V_{ud}\|^2[B_{negint}B_{exch}+𝒪(1/m_c)],$$ (15) and $$\mathrm{\Gamma }(\mathrm{\Xi }_b^{})\mathrm{\Gamma }(\mathrm{\Xi }_b^0)=\frac{G_F^2m_b^2}{4\pi }\|V_{cb}\|^2\|V_{ud}\|^2[B_{negint}B_{exch}+𝒪(1/m_b)],$$ (16) where the matrix elements of the four-quark operators between singly heavy baryons $`B_{Qqq^{}}`$ have the form $`B_{negint}=B_{Qqq^{}}\|_{i=1}^4c_{i,negint}^{Qq}(\mu )O_{i,negint}^{Qq}(\mu )\|B_{Qqq^{}}`$ and $`B_{exch}=B_{Qqq^{}}\|_{i=1}^2c_{i,exch}^{Qq}(\mu )O_{i,exch}^{Qq}(\mu )\|B_{Qqq^{}}`$. The final relation for singly heavy baryons looks like $$r^{bc}\frac{\mathrm{\Gamma }(\mathrm{\Xi }_b^{})\mathrm{\Gamma }(\mathrm{\Xi }_b^0)}{\mathrm{\Gamma }(\mathrm{\Xi }_c^+)\mathrm{\Gamma }(\mathrm{\Xi }_c^0)}=\frac{m_b^2}{m_c^2}\frac{\|V_{cb}\|^2}{\|V_{cs}\|^2}[1+𝒪(1/m_c,1/m_b)].$$ (17) This relation can be used to test the present model-dependent calculation of heavy-baryon lifetimes. Lifetime splittings both in the charmed and in the beauty sector have been investigated in several papers . We have recalculated our predictions from for charmed baryon lifetimes and from for beauty baryon lifetimes with the numerical parameters preferred in this paper and have used the recalculated predictions to test the relation (17), which we express in units $`\frac{m_b^2}{m_c^2}\frac{\|V_{cb}\|^2}{\|V_{cs}\|^2}`$ thus reducing the dependence on heavy-quark masses. In the case when the same approximations (neglecting mass corrections and Cabibbo-suppressed modes) are made, the above relation differs from unity by $`12\%`$. This number shows the deviation of the model-dependent calculation (where the contributions of four-quark operators are explicitly evaluated) from the model-independent prediction given only by the ratios of heavy-quark masses and Cabibbo matrix elements. The complete calculation with the mass corrections and Cabibbo-suppressed modes gives $`0.79`$ for the above ratio, which indicates the order of neglected corrections to be less than $`10\%`$. There are several important implications of relation (17). We can obtain a model-independent prediction for the still unmeasured lifetime difference between singly beauty baryons in the isospin doublet. Using the same numerical values for the parameters entering the right-hand side as before, and taking the experimentally measured lifetimes of singly-charmed baryons $`\mathrm{\Xi }_c^+`$ and $`\mathrm{\Xi }_c^0`$ from , we obtain $$\mathrm{\Gamma }(\mathrm{\Xi }_b^{})\mathrm{\Gamma }(\mathrm{\Xi }_b^0)=(0.14\pm 0.06)\mathrm{ps}^1.$$ (18) This lifetime difference can be compared with the model-independent prediction obtained by using the HQS and $`SU(3)`$ symmetry in , $`\mathrm{\Gamma }(\mathrm{\Xi }_b^{})\mathrm{\Gamma }(\mathrm{\Xi }_b^0)=(0.11\pm 0.03)\mathrm{ps}^1`$, and some moderate model-dependent prediction from , $`\mathrm{\Gamma }(\mathrm{\Xi }_b^{})\mathrm{\Gamma }(\mathrm{\Xi }_b^0)=0.094\mathrm{ps}^1`$. It is also worth mentioning that owing to the neglected mass corrections in final decay states of beauty baryons in the derivation of the ratios (17) and (18), there is no splitting between the nonleptonic rates of $`\mathrm{\Xi }_b^0`$ and $`\mathrm{\Lambda }_b`$. Since the Cabibbo-suppressed modes have also been neglected , the total decay rates of $`\mathrm{\Xi }_b^0`$ and $`\mathrm{\Lambda }_b`$ appear to be equal at this level. Therefore, the predictions from (17) are also valid for the $`\mathrm{\Gamma }(\mathrm{\Xi }_b^{})\mathrm{\Gamma }(\mathrm{\Lambda }_b)`$ difference. We can see that the results for the predicted splitting in the isospin doublet of beauty baryons obtained using different methods are all consistent. It is interesting to examine the explicit expression taken from , eq.(23), which can be rewritten as $$\mathrm{\Gamma }(\mathrm{\Xi }_b^{})\mathrm{\Gamma }(\mathrm{\Xi }_b^0)=\frac{m_b^2}{m_c^2}\frac{\|V_{cb}\|^2}{\|V_{cs}\|^2}[0.91\mathrm{\Gamma }(\mathrm{\Xi }_c^+)0.85\mathrm{\Gamma }(\mathrm{\Xi }_c^0)0.06\mathrm{\Gamma }(\mathrm{\Lambda }_c)],$$ (19) and which is obtained by applying $`SU(3)`$ symmetry, and with Cabibbo subleading effects included in charmed baryon decays. The appearance of the total decay rate for $`\mathrm{\Lambda }_c`$ is the direct consequence of the $`SU(3)`$ relations used. Otherwise, we can note a very similar structure to our relation (17) obtained within the $`SU(2)`$ approximation. The first two coefficients in front of the rates in (19) differ from one by less than $`15\%`$, where about $`510\%`$ difference comes from the Cabibbo-suppressed modes included, and the rest comes presumably from the difference between the $`SU(3)`$ and $`SU(2)`$ approximations applied. In the last case, we apply our procedure to the system of doubly heavy baryons. Again, we form pairs of one doubly charmed and one doubly beauty baryon. The first pair includes $`\mathrm{\Xi }_{cc}^{++}(ccu)`$ and $`\mathrm{\Xi }_{bb}^{}(bbd)`$. The decays of both baryons in this pair include the effect of negative Pauli interference. The second pair is formed from $`\mathrm{\Xi }_{cc}^+(ccd)`$ and $`\mathrm{\Xi }_{bb}^0(bbu)`$. Both of these baryons exhibit the effect of $`W`$ exchange. There are no Cabibbo-leading semileptonic four-quark contributions in doubly heavy baryons. Appropriate decay-rate differences are given by the following expressions: $$\mathrm{\Gamma }(\mathrm{\Xi }_{cc}^{++})\mathrm{\Gamma }(\mathrm{\Xi }_{cc}^+)=\frac{G_F^2m_c^2}{4\pi }\|V_{cs}\|^2\|V_{ud}\|^2[\mathrm{\Xi }_{cc}^{++}\|P^{cu}\|\mathrm{\Xi }_{cc}^{++}\mathrm{\Xi }_{cc}^+\|P^{cd}\|\mathrm{\Xi }_{cc}^+],$$ (20) $$\mathrm{\Gamma }(\mathrm{\Xi }_{bb}^{})\mathrm{\Gamma }(\mathrm{\Xi }_{bb}^0)=\frac{G_F^2m_b^2}{4\pi }\|V_{cb}\|^2\|V_{ud}\|^2[\mathrm{\Xi }_{bb}^{}\|P^{bd}\|\mathrm{\Xi }_{bb}^{}\mathrm{\Xi }_{bb}^0\|P^{bu}\|\mathrm{\Xi }_{bb}^0].$$ (21) Using isospin symmetry and HQS, we obtain, per analogiam with above derivations, relations for the decay-rate differences of doubly heavy baryons which lead to the final relation for doubly heavy baryons $$r^{bbcc}\frac{\mathrm{\Gamma }(\mathrm{\Xi }_{bb}^{})\mathrm{\Gamma }(\mathrm{\Xi }_{bb}^0)}{\mathrm{\Gamma }(\mathrm{\Xi }_{cc}^{++})\mathrm{\Gamma }(\mathrm{\Xi }_{cc}^+)}=\frac{m_b^2}{m_c^2}\frac{\|V_{cb}\|^2}{\|V_{cs}\|^2}[1+𝒪(1/m_c,1/m_b)].$$ (22) Again, we have achieved the reduction of the matrix elements of four-quark operators between doubly heavy baryons in the heavy quark limit. Relation (22) enables us to estimate the difference in the decay rates of doubly-beauty baryons using recent results for doubly-charmed baryon lifetimes. Using the calculated values $`\mathrm{\Gamma }(\mathrm{\Xi }_{cc}^{++})=0.952ps^1`$ and $`\mathrm{\Gamma }(\mathrm{\Xi }_{cc}^+)=5ps^1`$ with the values of the parameters $`m_c=1.35GeV`$, $`m_b=4.7GeV`$, $`\|V_{cb}\|=0.04`$ and $`\|V_{cs}\|=1.04`$ gives $$\mathrm{\Gamma }(\mathrm{\Xi }_{bb}^{})\mathrm{\Gamma }(\mathrm{\Xi }_{bb}^0)=0.073\mathrm{ps}^1.$$ (23) So far we have applied our reduction procedure separately to heavy mesons, singly-heavy baryons, and doubly-heavy baryons. Still, by inspection of relations (10), (17), and (22) we see that $`r^{BD}`$, $`r^{bc}`$, and $`r^{bbcc}`$ are the same up to $`𝒪(1/m_c,1/m_b)`$, i.e., $$\frac{\mathrm{\Gamma }(B^{})\mathrm{\Gamma }(B^0)}{\mathrm{\Gamma }(D^+)\mathrm{\Gamma }(D^0)}=\frac{\mathrm{\Gamma }(\mathrm{\Xi }_b^{})\mathrm{\Gamma }(\mathrm{\Xi }_b^0)}{\mathrm{\Gamma }(\mathrm{\Xi }_c^+)\mathrm{\Gamma }(\mathrm{\Xi }_c^0)}=\frac{\mathrm{\Gamma }(\mathrm{\Xi }_{bb}^{})\mathrm{\Gamma }(\mathrm{\Xi }_{bb}^0)}{\mathrm{\Gamma }(\mathrm{\Xi }_{cc}^{++})\mathrm{\Gamma }(\mathrm{\Xi }_{cc}^+)}=\frac{m_b^2}{m_c^2}\frac{\|V_{cb}\|^2}{\|V_{cs}\|^2}.$$ (24) Thus, we obtain a relation in (24) which clearly indicates certain universal behavior in the decays of all heavy hadrons. The existence of some universality could have been anticipated from the fact that the same expression (1) describes the decays of all heavy baryons. Still, that universality attains its concrete, model-independent form in the relation (24). This relation connects all sectors of heavy hadrons that are usually treated separately: mesonic and baryonic, charmed and beauty. Also, this relation brings some order in the otherwise rather intricate pattern of heavy-hadron lifetimes. An advantage of this relation is that by knowing some decay rates, one can calculate or give constraints on some other decay rates. Also, knowing decay rates from experiment, one can test the findings of the theory in a model-independent fashion. These results will help to establish the limitations of the present standard method of calculating inclusive processes and to test its underlying assumptions. Acknowledgements This work was supported by the Ministry of Science and Technology of the Republic of Croatia under the contract No. 00980102.
warning/0004/hep-th0004006.html	ar5iv	text	# The Lax pair for 𝐶₂-type Ruijsenaars-Schneider modelProject supported by the National Natural Science Foundation of China(Grant No. 19805006) ## I Introduction Ruijsenaars-Schneider(RS) and Calogero-Moser(CM) models as integrable many-body models recently have attracted remarkable attention and have been extensively studied. They describe one-dimensional $`n`$-particle system with pairwise interaction. Their importance lies in various fields ranging from lattice models in statistical physics$`^{\text{[1, 2]}}`$, to the field theory and gauge theory$`^{\text{[3, 4]}}`$, e.g., to the Seiberg-Witten theory $`^{\text{[5]}}`$, etc. Recently, the Lax pairs for the CM models in various root systems have been given by Olshanetsky et al$`^{\text{[6]}}`$, Inozemtsev$`^{\text{[7]}}`$, D’Hoker et al$`^{\text{[8]}}`$ and Bordner et al$`^{\text{[9]}}`$ with or without spectral parameter respectively. Further a more general algebra-geometric construction was proposed by Hurtubise et al in Ref. , while the commutative operators for the RS model based on various type Lie algebra were given by Komori$`^{\text{[11, 12]}}`$, Diejen$`^{\text{[13, 14]}}`$ and Hasegawa et al$`^{\text{[1, 15]}}`$. An interesting result is that in Ref. , the authors show that for the $`sl_2`$ trigonometric RS and CM models there exists the same non-dynamical $`r`$-matrix structure compared with the usual dynamical ones. On the other hand, similar to Hasegawa’s result that $`A_{N1}`$ RS model can be obtained as transfer matrices associated to the Sklyanin algebra, they also reveal that corresponding CM model’s integrability can be depicted by $`sl_N`$ Gaudin algebra$`^{\text{[17]}}`$. As for the $`C_n`$ type RS model, commuting difference operators acting on the space of functions on the $`C_2`$ type weight space have been constructed by Hasegawa et al in Ref. . Extending that work, the diagonalization of elliptic difference system of that type has been studied by Kikuchi in Ref. . Despite of the fact that the Lax pairs for CM models have been proposed for general Lie algebra even for all of the finite reflection groups$`^{\text{[19]}}`$, the Lax integrability of RS model are not clear except only for $`A_{N1}`$ -type$`^{\text{[20, 2, 21, 22, 23, 24]}}`$, i.e., the Lax pairs for the RS models other than $`A_{N1}`$ -type have not yet been obtained. In this paper, we concentrate on the $`C_2`$ type trigonometric Ruijsenaars-Schneider model. The basic materials about $`C_2`$ $`RS`$ model are reviewed in Section II. In Section III, we present the Lax pair without spectral parameter and its integrability in Liouville sense is also given. In Section IV, taking its non-relativistic limit, we recover the system of corresponding CM type. In Section V, we give the Lax pair for the system with spectral parameter, and show that at certain limit it will degenerate to the one without spectral parameter. The last section is a brief summary and some discussions. ## II Model and equations of motion As a relativistic-invariant generalization of the $`C_n`$-type Calogero-Moser model, the $`C_n`$-type Ruijsenaars-Schneider system is completely integrable whose integrability is shown by Ruijsenaars$`^{\text{[25]}}`$ and Diejen $`^{\text{[13, 14]}}`$. In terms of the canonical variables $`p_i`$, $`x_i(i,j=1,\mathrm{},2)`$ enjoying in the canonical Poisson bracket $`\{p_i,p_j\}=\{x_i,x_j\}=0,`$ $`\{x_i,p_j\}=\delta _{ij}`$, the Hamiltonian of $`C_2`$ RS system can be of the form $`H`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}\{e^{p_i}f(2x_i){\displaystyle \underset{ki}{\overset{2}{}}}(f(x_{ik})f(x_i+x_k))`$ (1) $`+e^{p_i}g(2x_i){\displaystyle \underset{ki}{\overset{2}{}}}(g(x_{ik})g(x_i+x_k))\},`$ where $`f(x)`$ $`:`$ $`={\displaystyle \frac{\mathrm{sin}(x\gamma )}{\mathrm{sin}(x)}},`$ $`g(x)`$ $`:`$ $`=f(x)\|_{\gamma \gamma },x_{ik}:=x_ix_k,`$ and $`\gamma `$ denotes the coupling constant. Notice that in Ref. Ruijsenaars used another “gauge” of the momenta such that the two systems are connected by the following canonical transformation: $`x_ix_i,p_ip_i+{\displaystyle \frac{1}{2}}ln{\displaystyle \underset{ji}{\overset{2}{}}}{\displaystyle \frac{f(x_{ij})f(x_i+x_j)}{g(x_{ij})g(x_i+x_j)}}{\displaystyle \frac{f(2x_i)}{g(2x_i)}}.`$ (2) The canonical equations of motion for the Hamiltonian (1) are $`\dot{x_i}`$ $`=`$ $`\{x_i,H\}=e^{p_i}b_ie^{p_i}b_i^{^{}},`$ (3) $`\dot{p_i}`$ $`=`$ $`\{p_i,H\}={\displaystyle \underset{ji}{\overset{2}{}}}(e^{p_j}b_j(h(x_{ji})h(x_j+x_i))`$ (4) $`+e^{p_j}b_j^{^{}}(\widehat{h}(x_{ji})\widehat{h}(x_j+x_i)))`$ $`e^{p_i}b_i\left(2h(2x_i)+{\displaystyle \underset{ji}{\overset{2}{}}}\left(h(x_{ij})+h(x_i+x_j)\right)\right)`$ $`e^{p_i}b_i^{^{}}\left(2\widehat{h}(2x_i)+{\displaystyle \underset{ji}{\overset{2}{}}}\left(\widehat{h}(x_{ij})+\widehat{h}(x_i+x_j)\right)\right),`$ where $`h(x):`$ $`=`$ $`{\displaystyle \frac{d\mathrm{ln}f(x)}{dx}},\text{ }\widehat{h}(x):={\displaystyle \frac{d\mathrm{ln}g(x)}{dx}},`$ $`b_i`$ $`=`$ $`f(2x_i){\displaystyle \underset{ki}{\overset{2}{}}}\left(f(x_ix_k)f(x_i+x_k)\right),`$ $`b_i^{^{}}`$ $`=`$ $`g(2x_i){\displaystyle \underset{ki}{\overset{2}{}}}\left(g(x_ix_k)g(x_i+x_k)\right).`$ (5) Here, of course $`x_i=x_i(t)`$, $`p_i=p_i(t)`$ and the dot on top denotes t-differentiation. ## III The Lax pair without spectral parameter Let us first mention some results about the integrability of Hamiltonian (1). In Ref. Ruijsenaars demonstrated that the symplectic structure of $`C_n`$ type RS system can be proved integrable by embedding its phase space to a submanifold of $`A_{2n1}`$ type RS one, while in Refs. and Ref. , Diejen and Komori, respectively, gave a series of commuting difference operators which led to its quantum integrability. However, there is not any result about its Lax representation so far. That is, the explicit form of the Lax matrix $`L`$, associated with a $`M`$ which ensure its Lax integrability, have not been proposed up to now. In this section, we restrict our treatment to the exhibition of the explicit form for $`C_2`$ RS system. Therefore, some previous results, as well as new results, could now be obtained in a more straightforward manner by using the Lax pair. Define one $`4\times 4`$ Lax matrix for $`C_2`$ RS model as follows: $`L=\left(\begin{array}{cc}A\hfill & B\hfill \\ C\hfill & D\hfill \end{array}\right),`$ (8) where $`A`$, $`B`$, $`C`$, $`D`$ are $`2\times 2`$ matrices(hereafter, we use the indices $`i,j=1,2`$) $`A_{ij}`$ $`=`$ $`e^{p_j}b_j{\displaystyle \frac{\mathrm{sin}\gamma }{\mathrm{sin}(x_{ij}+\gamma )}},B_{ij}=e^{p_j}b_j^{^{}}{\displaystyle \frac{\mathrm{sin}\gamma }{\mathrm{sin}(x_i+x_j+\gamma )}},`$ $`C_{ij}`$ $`=`$ $`e^{p_j}b_j{\displaystyle \frac{\mathrm{sin}\gamma }{\mathrm{sin}(x_ix_j+\gamma )}},D_{ij}=e^{p_j}b_j^{^{}}{\displaystyle \frac{\mathrm{sin}\gamma }{\mathrm{sin}(x_{ji}+\gamma )}}.`$ (9) For the concise expression for $`M`$, we define four auxiliary $`2\times 2`$ matrices $`\stackrel{~}{A}`$, $`\stackrel{~}{B}`$, $`\stackrel{~}{C}`$, $`\stackrel{~}{D}`$ as follows $`\stackrel{~}{A}_{ij}`$ $`=`$ $`e^{p_i}b_j^{^{}}{\displaystyle \frac{\mathrm{sin}\gamma }{\mathrm{sin}(x_{ij}\gamma )}},\stackrel{~}{B}_{ij}=e^{p_i}b_j{\displaystyle \frac{\mathrm{sin}\gamma }{\mathrm{sin}(x_i+x_j\gamma )}},`$ $`\stackrel{~}{C}_{ij}`$ $`=`$ $`e^{p_i}b_j^{^{}}{\displaystyle \frac{\mathrm{sin}\gamma }{\mathrm{sin}(x_ix_j\gamma )}},\stackrel{~}{D}_{ij}=e^{p_i}b_j{\displaystyle \frac{\mathrm{sin}\gamma }{\mathrm{sin}(x_{ji}\gamma )}},`$ (10) such that $`M`$ can be of the form $`M=\left(\begin{array}{cc}𝒜\hfill & \hfill \\ 𝒞\hfill & 𝒟\hfill \end{array}\right),`$ (13) where entries of $`M`$ are $`𝒜_{ij}`$ $`=`$ $`\mathrm{cot}(x_{ij})(A_{ij}\stackrel{~}{A}_{ij}),𝒟_{ij}=\mathrm{cot}(x_{ji})(D_{ij}\stackrel{~}{D}_{ij}),(ij),`$ $`_{ij}`$ $`=`$ $`\mathrm{cot}(x_i+x_j)(B_{ij}\stackrel{~}{B}_{ij}),𝒞_{ij}=\mathrm{cot}(x_ix_j)(C_{ij}\stackrel{~}{C}_{ij}),`$ $`𝒜_{ii}`$ $`=`$ $`{\displaystyle \underset{ki}{\overset{2}{}}}{\displaystyle \frac{A_{ik}\stackrel{~}{A}_{ik}}{\mathrm{sin}(x_{ik})}}{\displaystyle \underset{k=1}{\overset{2}{}}}{\displaystyle \frac{B_{ik}\stackrel{~}{B}_{ik}}{\mathrm{sin}(x_i+x_k)}},`$ $`𝒟_{ii}`$ $`=`$ $`{\displaystyle \underset{ki}{\overset{2}{}}}{\displaystyle \frac{D_{ik}\stackrel{~}{D}_{ik}}{\mathrm{sin}(x_{ik})}}+{\displaystyle \underset{k=1}{\overset{2}{}}}{\displaystyle \frac{C_{ik}\stackrel{~}{C}_{ik}}{\mathrm{sin}(x_i+x_k)}}.`$ (14) We have checked that $`L,M`$ satisfies the Lax equation $`\dot{L}=\{L,H\}=[M,L],`$ (15) which is equivalent to the equations of motion (3) and (4) with the help of computer. The Hamiltonian $`H`$ can be rewritten in the following form $`H={\displaystyle \underset{j=1}{\overset{2}{}}}(e^{p_j}b_j+e^{p_j}b_j^{^{}})=trL.`$ (16) The characteristic polynomial of the Lax matrix $`L`$ is $`det(LvId)`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{4}{}}}(v)^{4j}H_j`$ (17) $`=`$ $`v^4Hv^3+H_2v^2Hv+1,`$ where $`H_0=H_4=1`$, $`H_1=H_3=H`$. The function-independent Hamiltonian flows $`H`$ and $`H_2`$ are $`H`$ $`=`$ $`e^{p_1}f(2x_1)f(x_{12})f(x_1+x_2)`$ (18) $`+e^{p_1}g(2x_1)g(x_{12})g(x_1+x_2)`$ $`+e^{p_2}f(2x_2)f(x_{21})f(x_2+x_1)`$ $`+e^{p_2}g(2x_2)g(x_{21})g(x_2+x_1),`$ $`H_2`$ $`=`$ $`e^{p_1+p_2}f(2x_1)(f(x_1+x_2))^2f(2x_2)`$ (19) $`+e^{p_1p_2}g(2x_1)(g(x_1+x_2))^2g(2x_2)`$ $`+e^{p_1p_2}f(2x_1)(f(x_{12}))^2f(2x_2)`$ $`+e^{p_2p_1}g(2x_1)(g(x_{12}))^2g(2x_2)`$ $`+2f(x_{12})g(x_{12})f(x_1+x_2)g(x_1+x_2).`$ We verify that $`H`$ and $`H_2`$ Poisson commute each other $`\{H,H_2\}=0,`$ (20) which ensures the complete integrability of $`C_2`$ RS model (in Liouville sense). ## IV Nonrelativistic limit to the Calogero-Moser system The Nonrelativistic limit can be achieved by rescaling $`p_i\beta p_i`$, $`\gamma \beta \gamma `$ while letting $`\beta 0,`$ and making a canonical transformation $`p_ip_i+\gamma \left(\mathrm{cot}(2x_i)+{\displaystyle \underset{ki}{\overset{2}{}}}\left(\mathrm{cot}(x_{ik})+\mathrm{cot}(x_i+x_k)\right)\right),`$ (21) such that $`L`$ $``$ $`Id+\beta L_{CM}+O(\beta ^2),`$ (22) $`M`$ $``$ $`2\beta M_{CM}+O(\beta ^2),`$ (23) and $`H`$ $``$ $`4+2\beta ^2H_{CM}+O(\beta ^2).`$ (24) $`L_{CM}`$ can be expressed as $`L_{CM}=\left(\begin{array}{cc}A_{CM}\hfill & B_{CM}\hfill \\ B_{CM}\hfill & A_{CM}\hfill \end{array}\right),`$ (27) where $`(A_{CM})_{ii}`$ $`=`$ $`p_i,(B_{CM})_{ij}={\displaystyle \frac{\gamma }{\mathrm{sin}(x_i+x_j)}},`$ $`(A_{CM})_{ij}`$ $`=`$ $`{\displaystyle \frac{\gamma }{\mathrm{sin}(x_{ij})}},(ij).`$ (28) $`M_{CM}`$ is $`M_{CM}=\left(\begin{array}{cc}𝒜_{CM}\hfill & _{CM}\hfill \\ _{CM}\hfill & 𝒜_{CM}\hfill \end{array}\right),`$ (31) where $`(𝒜_{CM})_{ii}`$ $`=`$ $`{\displaystyle \underset{ki}{\overset{2}{}}}({\displaystyle \frac{\gamma }{\mathrm{sin}^2x_{ik}}}+{\displaystyle \frac{\gamma }{\mathrm{sin}^2(x_i+x_k)}}){\displaystyle \frac{\gamma }{\mathrm{sin}^2(2x_i)}},(_{CM})_{ij}={\displaystyle \frac{\gamma \mathrm{cos}(x_i+x_j)}{\mathrm{sin}^2(x_i+x_j)}},`$ $`(𝒜_{CM})_{ij}`$ $`=`$ $`{\displaystyle \frac{\gamma \mathrm{cos}(x_{ij})}{\mathrm{sin}^2x_{ij}}},(ij),`$ (32) which coincide with the form given in Ref. with the difference of a constant diagonalized matrix. The Hamiltonian of $`C_2`$-type CM model can be given by $`H_{CM}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=1}{\overset{2}{}}}p_k^2{\displaystyle \frac{\gamma ^2}{2}}{\displaystyle \underset{ki}{\overset{2}{}}}({\displaystyle \frac{1}{\mathrm{sin}^2x_{ik}}}+{\displaystyle \frac{1}{\mathrm{sin}^2(x_i+x_k)}}+{\displaystyle \frac{1}{\mathrm{sin}^2(2x_i)}})`$ (33) $`=`$ $`{\displaystyle \frac{1}{4}}trL^2.`$ The $`L_{CM}`$, $`M_{CM}`$ satisfies the Lax equation $`\dot{L}_{CM}=\{L_{CM},H_{CM}\}=[M_{CM},L_{CM}].`$ (34) ## V The Lax pair with spectral parameter Also, we can give the Lax pair which include spectral parameters. Define the Lax matrix for trigonometric RS model as follows: $`L=\left(\begin{array}{cc}A\hfill & B\hfill \\ C\hfill & D\hfill \end{array}\right),`$ (37) where $`A`$, $`B`$, $`C`$, $`D`$ are $`2\times 2`$ matrices($`i,j=1,2`$) $`A_{ij}`$ $`=`$ $`e^{p_j}b_j{\displaystyle \frac{\mathrm{sin}(x_{ij}+\gamma +\lambda )\mathrm{sin}\gamma }{\mathrm{sin}(x_{ij}+\gamma )\mathrm{sin}(\gamma +\lambda )}},B_{ij}=e^{p_j}b_j^{^{}}{\displaystyle \frac{\mathrm{sin}(x_i+x_j+\gamma +\lambda )\mathrm{sin}\gamma }{\mathrm{sin}(x_i+x_j+\gamma )\mathrm{sin}(\gamma +\lambda )}},`$ $`C_{ij}`$ $`=`$ $`e^{p_j}b_j{\displaystyle \frac{\mathrm{sin}(x_ix_j+\gamma +\lambda )\mathrm{sin}\gamma }{\mathrm{sin}(x_ix_j+\gamma )\mathrm{sin}(\gamma +\lambda )}},D_{ij}=e^{p_j}b_j^{^{}}{\displaystyle \frac{\mathrm{sin}(x_{ji}+\gamma +\lambda )\mathrm{sin}\gamma }{\mathrm{sin}(x_{ji}+\gamma )\mathrm{sin}(\gamma +\lambda )}}.`$ (38) $`M`$ is $`M=\left(\begin{array}{cc}𝒜\hfill & \hfill \\ 𝒞\hfill & 𝒟\hfill \end{array}\right),`$ (41) where entries of $`M`$ are $`𝒜_{ij}`$ $`=`$ $`e^{p_j}b_j{\displaystyle \frac{\mathrm{sin}(x_{ij}+\lambda )}{\mathrm{sin}\lambda \mathrm{sin}x_{ij}}}e^{p_i}b_j^{^{}}{\displaystyle \frac{\mathrm{sin}(x_{ij}+\lambda +4\gamma )}{\mathrm{sin}(\lambda +4\gamma )\mathrm{sin}x_{ij}}},`$ $`𝒟_{ij}`$ $`=`$ $`e^{p_j}b_j^{^{}}{\displaystyle \frac{\mathrm{sin}(x_{ji}+\lambda )}{\mathrm{sin}\lambda \mathrm{sin}x_{ji}}}e^{p_i}b_j{\displaystyle \frac{\mathrm{sin}(x_{ji}+\lambda +4\gamma )}{\mathrm{sin}(\lambda +4\gamma )\mathrm{sin}x_{ji}}},(ij),`$ $`_{ij}`$ $`=`$ $`e^{p_j}b_j^{^{}}{\displaystyle \frac{\mathrm{sin}(x_i+x_j+\lambda )}{\mathrm{sin}\lambda \mathrm{sin}(x_i+x_j)}}e^{p_i}b_j{\displaystyle \frac{\mathrm{sin}(x_i+x_j+\lambda +4\gamma )}{\mathrm{sin}(\lambda +4\gamma )\mathrm{sin}(x_i+x_j)}},`$ $`𝒞_{ij}`$ $`=`$ $`e^{p_j}b_j{\displaystyle \frac{\mathrm{sin}(x_i+x_j\lambda )}{\mathrm{sin}\lambda \mathrm{sin}(x_i+x_j)}}e^{p_i}b_j^{^{}}{\displaystyle \frac{\mathrm{sin}(x_i+x_j\lambda 4\gamma )}{\mathrm{sin}(\lambda +4\gamma )\mathrm{sin}(x_i+x_j)}},`$ $`𝒜_{ii}`$ $`=`$ $`(\mathrm{cot}(\gamma )+\mathrm{cot}(\lambda ))e^{p_i}b_i(\mathrm{cot}(\lambda +\gamma )\mathrm{cot}(\gamma ))e^{p_i}b_i^{^{}}`$ $`+{\displaystyle \underset{ki}{\overset{2}{}}}((\mathrm{cot}(x_{ik}+\gamma )\mathrm{cot}(x_{ik})e^{p_k}b_k)`$ $`+{\displaystyle \frac{\mathrm{sin}(x_{ik}+\lambda +4\gamma )\mathrm{sin}(x_{ki}+\lambda +\gamma )\mathrm{sin}\gamma }{\mathrm{sin}(\lambda +4\gamma )\mathrm{sin}x_{ik}\mathrm{sin}(x_{ki}+\gamma )\mathrm{sin}(\lambda +\gamma )}}e^{p_i}b_k^{^{}})`$ $`+{\displaystyle \underset{k=1}{\overset{2}{}}}((\mathrm{cot}(x_i+x_k+\gamma )\mathrm{cot}(x_i+x_k)e^{p_k}b_k^{^{}}`$ $`+{\displaystyle \frac{\mathrm{sin}(x_i+x_k+\lambda +4\gamma )\mathrm{sin}(x_i+x_k\lambda \gamma )\mathrm{sin}\gamma }{\mathrm{sin}(\lambda +4\gamma )\mathrm{sin}(x_i+x_k)\mathrm{sin}(x_i+x_k\gamma )\mathrm{sin}(\lambda +\gamma )}}e^{p_i}b_k)`$ $`𝒟_{ii}`$ $`=`$ $`(\mathrm{cot}(\gamma )+\mathrm{cot}(\lambda ))e^{p_i}b_i^{^{}}(\mathrm{cot}(\lambda +\gamma )\mathrm{cot}(\gamma ))e^{p_i}b_i`$ (42) $`+{\displaystyle \underset{ki}{\overset{2}{}}}((\mathrm{cot}(x_{ki}+\gamma )\mathrm{cot}(x_{ki})e^{p_k}b_k^{^{}})`$ $`+{\displaystyle \frac{\mathrm{sin}(x_{ki}+\lambda +4\gamma )\mathrm{sin}(x_{ik}+\lambda +\gamma )\mathrm{sin}\gamma }{\mathrm{sin}(\lambda +4\gamma )\mathrm{sin}x_{ki}\mathrm{sin}(x_{ik}+\gamma )\mathrm{sin}(\lambda +\gamma )}}e^{p_i}b_k)`$ $`+{\displaystyle \underset{k=1}{\overset{2}{}}}((\mathrm{cot}(x_i+x_k)\mathrm{cot}(x_i+x_k\gamma ))e^{p_k}b_k`$ $`+{\displaystyle \frac{\mathrm{sin}(x_i+x_k\lambda 4\gamma )\mathrm{sin}(x_i+x_k+\lambda +\gamma )\mathrm{sin}\gamma }{\mathrm{sin}(\lambda +4\gamma )\mathrm{sin}(x_i+x_k)\mathrm{sin}(x_i+x_k+\gamma )\mathrm{sin}(\lambda +\gamma )}}e^{p_i}b_k^{^{}}).`$ The $`L,M`$ satisfies the Lax equation Eq.(15) and the Hamiltonian $`H`$ can also be rewritten in the form of Eq.(16). The function-independent Hamiltonian flows can be generated by calculating the characteristic polynomial of that Lax matrix $`L`$ $$det(LvId)=\underset{j=0}{\overset{4}{}}\frac{(\mathrm{sin}\lambda )^{^{(j1)}}\mathrm{sin}(\lambda +j\gamma )}{(\mathrm{sin}(\gamma +\lambda ))^j}(v)^{4j}H_j,$$ (43) where $`H_0=H_4=1`$, $`H_1=H_3=H`$. $`H`$ and $`H_2`$ have the same forms as Eq.(18) and (19). Remarks: 1. As far as the forms of the Lax pair for the rational-type systems are concerned, we can get them by making the following substitutions $`\mathrm{sin}x`$ $``$ $`x,`$ $`\mathrm{cos}x`$ $``$ $`1,`$ for all the above statements. 2. It should be pointed out that the Lax pair given in Eqs.(8)-(14) which are without spectral parameter can be derived from the one with spectral parameters (see Eqs.(37)-(42)) by taking the following limit $$\lambda i\mathrm{},$$ up to an appropriate gauge transformation of the Lax matrix with a diagonal matrix. ## VI Summary and discussions In this paper, we propose the Lax pairs for trigonometric $`C_2`$ RS model together with its rational limit and show their integrability. Involutive Hamiltonians are shown to be generated by the characteristic polynomial of the Lax matrix. In the nonrelativistic limit, the system leads to CM system associated with the root system of $`C_2`$ which is known previously. It is expected that, in the general case of $`C_n`$ for $`n2`$, the explicit expressions of $`L`$ and $`M`$ must have similar forms as those presented here. So it would be interesting to make some progress in this respect in the very near future. ## Acknowledgement One of the authors K. Chen is grateful to professors K. J. Shi and L. Zhao for their encouragement.
warning/0004/hep-ph0004072.html	ar5iv	text	# ELECTROWEAK SYMMETRY BREAKING BY EXTRA DIMENSIONS EFI-2000-11 Fermilab-Conf-00-077-T Electroweak symmetry breaking may be naturally induced by the observed quark and gauge fields in extra dimensions without a fundamental Higgs field. We show that a composite Higgs doublet can arise as a bound state of $`(t,b)_L`$ and a linear combination of the Kaluza-Klein states of $`t_R`$, due to QCD in extra dimensions. The top quark mass depends on the number of active $`t_R`$ Kaluza-Klein modes, and is consistent with the experimental value. Understanding the origin of electroweak symmetry breaking is currently one of the most important questions in particle physics. In the Standard Model (SM), it is assumed that there exists a fundamental scalar Higgs field with a negative squared mass, whose vacuum expectation value (vev) breaks the electroweak symmetry. However, the mass squared of a fundamental scalar field receives a quadratic divergent contribution from radiative corrections, therefore suffers from the “hierarchy problem” if the cutoff scale is much higher than the weak scale. Recently, a new solution to the hierarchy problem has been proposed by postulating the existence of large extra dimensions.$`^\mathrm{?}`$ In that case, the fundamental scale (which is assumed to be the cutoff scale) in the higher dimensional theory can be reduced to the TeV range and hence it removes the large hierarchy between the weak scale and the Planck scale. However, just removing the large hierarchy, while avoiding the fine-tuning problem, does not explain why the electroweak symmetry is broken, i.e., why there is a Higgs field and why its squared mass is negative. Here we point out that the extra dimensions can also provide a natural mechanism for electroweak symmetry breaking without introducing a fundamental Higgs field. In fact, a composite Higgs field can arise in the presence of certain strongly coupled four-quark operators.$`^\mathrm{?}`$ These four-quark operators are naturally induced by the Kaluza-Klein (KK) excitations of the gluons if QCD lives in compact extra dimensions.$`^\mathrm{?}`$ The strength of these contact interactions depends on the ratio of the compactification scale, $`M_c`$, and the fundumental scale $`M_s`$ of the quantum gravitational effects which cuts off the higher dimensional gauge interactions. It has been argued that the SM gauge couplings can unify in the presence of extra dimensions due to the accelerated power-law running.$`^\mathrm{?}`$ Assuming that they unify at $`M_s`$, one typically finds that the gauge couplings already become strong at that scale for more than two extra dimensions. Thus, it is possible that the non-perturbative effects form a composite Higgs out of the quarks. The simplest 4-dimensional top-condensate model of a composite Higgs predicts a too large top-quark mass, of approximately $`500600`$ GeV, if the compositeness scale is not much higher than the weak scale.$`^\mathrm{?}`$ For a viable composite Higgs model, some vector-like quarks are required to participate in the binding mechanism together with the top quark.$`^\mathrm{?}`$ They should have the same SM quantum numbers as the top quark, so they can naturally be identified as the KK excitations of the top quark in a theory with extra dimensions.$`^\mathrm{?}`$ For instance, if we assume that only the right-handed top lives in extra dimensions, the Higgs doublet will be a bound state of the left-handed top-bottom doublet ($`\psi _L`$) and a linear combination of the KK modes of $`t_R`$, $$H\psi _L\underset{i=0}{\overset{n_{\mathrm{KK}}1}{}}\chi _R^i,$$ (1) where $`\chi _R^0=t_R`$. The top quark mass will be suppressed by $`1/\sqrt{n_{\mathrm{KK}}}`$ compared with the prediction of the simplest top condensate model. For a typical $`n_{\mathrm{KK}}1030`$, the top Yukawa coupling is around 1, in agreement with the experimental result. We now present a model for concreteness. This is not a unique choice, but just an illustrative example of the idea. We assume that SM gauge fields propagate in $`\delta `$ compact extra dimensions, whose coordinates are labeled by $`y,z_1,\mathrm{},z_{\delta 1}`$, with sizes $`L`$ and $`L_z(L)`$ respectively. The $`t_R`$ is the zero mode of a 5-dimensional fermion $`\chi `$, which is fixed at $`z=0`$, but propagates on the $`[0,L]`$ interval in the $`y`$ direction, while the $`\psi _L=(t,b)_L`$ is fixed at $`z=0`$ and $`y=y_0`$. (See Fig. 1.) The absence of a left-handed zero mode of $`\chi `$ can result from some boundary condition such as the $`S_1/Z_2`$ orbifold projection. For simplicity, we assume that all other quarks are 4-dimensional fields with left- and right-handed quarks localized at different positions in extra dimensions so that they do not form bound states. Because $`L_zL`$, we first integrate out the $`z`$ directions below the scale $`L_z^1`$ and obtain a 5-dimensional effective theory. The 4-quark interactions are induced by the KK gluons in the $`z`$ directions. After a Fierz transformation, the 4-quark interactions at the compositeness scale $`\mathrm{\Lambda }`$ are given by $$\frac{cg_5^2}{\mathrm{\Lambda }^2}\left\{\delta (yy_0)\left(\overline{\psi }_L\chi \right)\left(\overline{\chi }\psi _L\right)+\frac{5}{16}\left[\left(\overline{\chi }\chi \right)^2\frac{1}{3}\left(\overline{\chi }\gamma _5\chi \right)^2\right]\right\}+\mathrm{},$$ (2) where $`c1`$ represents the effect of suming over gluon KK modes, and the ellipsis stand for vectorial and tensorial four-quark operators, which are not relevant at low energies. From eq. 2 one can see that the attractive interactions can give rise to the following bound states: a 4-dimensional weak doublet complex scalar, $`H(x^\mu )\overline{\chi }\psi _L`$, and a 5-dimensional gauge singlet real scalar, $`\phi (x^\mu ,y)\overline{\chi }\chi `$. They obtain non-zero kinetic terms in running down to low energies and become dynamical fields. They also receive large negative contributions to their squared masses. Electroweak symmetry is broken when the squared mass of the composite Higgs becomes negative. Decomposing $`\chi `$ into 4-dimensional KK states, one can see indeed that the Higgs field is a bound state of the doublet $`\psi _L`$ and a linear combination of the KK modes of $`t_R`$. Since the observed (right-handed) top quark is only one of the $`n_{\mathrm{KK}}`$ states which participate in the electroweak symmetry breaking, the top Yukawa coupling is adequately suppressed by $`1/\sqrt{n_{\mathrm{KK}}}`$ as mentioned above. Whether the singlet $`\phi `$ affects the low energy physics at the weak scale depends on the setup. To illustrate this we consider two special cases. First, if $`\psi _L`$ is located at the boundary ($`y_0=0`$), the Higgs is more strongly bound and hence $`\phi `$ will remain heavy when $`m_H^2`$ becomes negative. There is no mixing between $`H`$ and $`\phi `$ because $`\phi `$ vanishes at the boundary. The low energy theory is simply the Standard Model with a composite Higgs boson which is expected to be heavy. It can still be consistent with the electroweak precision measurements because of the mixings between $`t,W,Z`$ and their KK states.$`^\mathrm{?}`$ Another interesting case is that $`\psi _L`$ is localized in the middle of the $`0<y<L`$ interval. In this case, the attractive interactions in the $`H`$ and $`\phi `$ channels are comparable, therefore both can obtain non-zero vevs. The mixing between $`H`$ and $`\phi `$ can make either the Higgs boson or the singlet light. If the mass of the lightest singlet is less than half of the Higgs mass, the Higgs boson may decay predominantly into two $`\phi `$’s, which subsequently decay into two gluons or two photons through the top loop, because the Higgs boson interacts strongly with $`\phi `$. This will modify the Higgs search at future colliders. The model we have summarized here represents a minimal model in extra dimensions. The only fields present at the fundamental scale are the standard fermions and gauge bosons in the higher dimensional spacetime. It is remarkable that at energies below the compactification scale this simple theory reproduces the Standard Model, with the possible addition of a light singlet scalar. We emphasize though that there is need for flavor violating operators at the fundamental scale in order to generate masses for the quarks (other than top) and leptons. These operators induce Yukawa couplings at low energy. Therefore the fermion masses are accomodated as in the SM without theoretical predictions other than the top mass.
warning/0004/physics0004040.html	ar5iv	text	# The Physical Origin of Intrinsic Bends in Double Helical DNA ## Introduction The intrinsic sequence dependent curvature of the DNA molecule is likely to be involved in fundamental mechanisms of genome regulation. The possibility of strong static bends in the B-DNA double helix has been proven for sequences containing regular repeats of $`\mathrm{A}_\mathrm{n}\mathrm{T}_\mathrm{m},\mathrm{with}\mathrm{n}+\mathrm{m}>3`$, called A-tracts . This effect plays an important role in DNA studies because it is unique in the amount and the variety of the available experimental information and, possibly, it can serve as a gate to the unknown general mechanisms of recognition and regulation of genome sequences. Every A-tract deviates the helical axis in a locally fixed direction by approximately 18, and, if the A-tracts are repeated in phase with the helical screw, a macroscopic curvature emerges. The effect was first noticed and identified in restriction fragments from the kinetoplast body of Leishmania tarentolae , and confirmed by electric birefringence decay and electron microscopy . A large variety of interesting information has been obtained by biochemical methods. It appeared that the double helix bends towards the minor grooves of A-tracts . The curvature is reduced with the temperature above 40 and in high salt, but for some sequences it is increased in presence of divalent metal ions . It depends upon the length and composition of A-tracts as well as on sequences between them . Detailed analysis of these results can be found in comprehensive reviews published in different years . According to many independent experimental observations, the structure of A-tract sequences should differ significantly from the “random” B-DNA. It is well established that, in solution, the poly-dA double helix is overwound to a twist of around 36 from around 34 of a random sequence . The models constructed from fiber diffraction data suggest consistently that the poly-dA double helix is characterized by a very narrow minor groove and a high propeller twist . Yet another distinction is an apparently large negative inclination of base pairs . Several A-tracts avaliable in single crystal structures of B-DNA oligomers have irregular conformations, but exhibit similar trends toward their centers . Even though the curvature is apparently caused by A-tracts, in Xray structures, A-tracts look generally less prone to bending than other sequences. Some indirect observations also support this view, notably, poly-dA fragments move faster than random DNA in gel migration assays and avoid wrapping around nucleosome particles . In spite of a large body of the experimental information accumulated during the last twenty years, the possible physical mechanism of this effect remains unclear. Since every base pair in a stack interacts only with the two neighbors, any sequence specificity in the DNA structure should mainly depend upon the stacking interactions in one base pair step. Non-local effects are also possible, however, due to base-backbone interactions and propagation of correlations along the backbone. The initial experimental data on A-tract bending were interpreted in terms of two alternative mechanisms, namely, the wedge model and the junction model . Both had to be modified significantly as and when new experimental data appeared and some other theories were discussed as well. The possible mechanisms of bending considered in the literature will be discussed below. Here we note only that none of them explains all experimental data and can be definitely preferred . The overall pattern has been additionally complicated when it was found that certain non A-tract sequences also exhibit distinguishable curvature . Here we report new results obtained in free unbiased molecular dynamics simulations of a DNA oligomer with phased A-tracts. We managed to find computational conditions in which stable sequence dependent static curvature emerges spontaneously in good agreement with experimental observations in terms of both the bending direction and magnitude. It is found that three independent long time trajectories converge to conformations with similar bends, but rather different local structural parameters, in evident contradiction with the common views of the origin of curvature. Analysis of these discrepancies leads to a new, significantly different hypothesis of the possible mechanism of intrinsic bends in the double helical DNA. ### Theoretical Background As a theoretical problem, the phenomenon of the sequence dependent DNA bending presents a challenge in many aspects similar to the protein folding problem. In order to understand the underlying physical mechanism, one has to analyze terribly noisy experimental data, with the noise being due to the biological diversity and, therefore, unremovable. If and when the physical mechanism is understood one will have to tackle this diversity again, because it will be necessary to analyze specific DNA sequences. The atom level molecular modeling is virtually the only theoretical approach that is potentially able to treat these difficulties. Ideally, we would like to have a model where the base pair sequence represents the only initial bias towards a specific conformation. If it could reproducibly yield curved DNA conformations in agreement with experimental data we should be able disclose the mechanism of bending in the model, and hope that a similar mechanism takes place in the nature. The foregoing scheme, however, is too difficult and, until now, most of the modeling studies used other strategies. Much has been learned about the DNA bending mechanics by using energy minimization and Monte Carlo . Unfortunately, the possibilities of such studies are limited by the multiple minima problem especially when it is necessary to take into consideration specific interaction with solvent molecules. The proposed alternative strategies commonly involved some bias towards specifically bent conformations introduced either explicitly, by imposing restraints, or implicitly, by choosing particular initial conditions, which made impossible unequivocal conclusions concerning the possible mechanism of bending. In the recent years, owing to the progress achieved in improving the full atom force fields, multi-nanosecond free MD simulations of DNA became feasible . A few such studies of phased A-tract sequences have been already reported . It has been demonstrated that, without any a priori bias, the DNA double helix bends anisotropically, and certain sequence specific features of A-tracts were at least qualitatively reproduced. At present, the free MD simulations represent the most promising line of research in this field, and we continue it here by using the recently proposed minimal model of B-DNA . The minimal B-DNA consists of a double helix with the minor groove filled with explicit water. Unlike the more widely used models, it does not involve explicit counterions and damps long range electrostatic interactions in a semi-empirical way by using distance scaling of the electrostatic constant and reduction of phosphate charges. We have earlier found that the minimal model gives B-DNA conformations which better compare with experimental structures than DNA structures obtained with other computational methods. Notably, it is free from a systematic negative bias of the average twist observed in simulations with full hydration and explicit counterions . This factor is likely to be involved in DNA bending because, as noted above, A-tracts are overwound with respect to the average B-DNA. For the standard test case of the EcoRI dodecamer structure, the dynamics of the minimal model reproducibly converged to structures slightly bent towards the minor groove which was narrowed in excellent agreement with the single crystal conformation . All these preliminary observations suggested that the minimal model was a good choice for studying the DNA bending induced by A-tracts. We report here simulations of the bending dynamics of a 25-mer B-DNA fragment. Its sequence, $`\mathrm{AAAATAGGCTATTTTAGGCTATTTT}`$, has been constructed after many preliminary tests with shorter sequence motives, and it includes three A-tracts separated by one helical turn. Our general strategy came out from the following considerations. Although the A-tract sequences that induce the strongest bends are known from experiments, probably not all of them would work in simulations. There are natural limitations, such as the precision of the model, and, in addition, the limited duration of trajectories may be insufficient for some A-tracts to adopt their specific conformation. Also, there is little experimental evidence of static curvature in short DNA fragments, notably, one may well expect the specific A-tract structure to be unstable near the ends. That is why we did not simply take the strongest experimental “benders”, but looked for sequence motives that in calculations readily adopt the characteristic local structure, with a narrow minor groove profile and high propeller twist, both in the middle and near the ends of the duplex. The complementary duplex $`\mathrm{AAAATAGGCTATTTTAGGCTATTTT}`$ has been constructed by repeating and inverting one such motive. Our goal was to find sequences that would appear statically bent in these conditions. It means that, in dynamics, the structure should fluctuate around a state with a distinguishable bend and a definite bending direction. Since any MD simulation is limited in time, there is no way to prove rigorously that some specific conformation is representative. Some degree of confidence can be achieved, however, if independent trajectories are able converge to the same state, and this is exactly what we tried to obtain. We found important to place A-tracts at both ends because of the two reasons. First, we could study only short DNA fragments, therefore, it was preferable to place A-tracts at both ends in order to maximize the possible bend. Second, in calculations with short sequences, it is important to be sure that the boundary conditions are correct. Since the A-tracts have a characteristic local structure, the boundary conditions could be at least qualitatively verified, which would not be the case for GC-rich sequences. ## Results Three long MD trajectories were computed for a complementary DNA duplex with the sequence $`\mathrm{AAAATAGGCTATTTTAGGCTATTTT}`$. The model system employed was same in all three simulations, with only the starting states varied. The first trajectory referred to below as TJBa started from the fiber canonical B-DNA structure and continued to 10 ns. When it was found that TJBa converged to a statically bent conformation, in good qualitative agreement with expectations based upon experimental data, another trajectory (TJBb) was computed in order to verify the reproducibility of the results. It started with random velocities from a re-minimized straight conformation taken from the initial phase of TJBa and was also continued to 10 ns. Simultaneously, in order to remove any initial bias implicitly involved in the choice of the starting state, the third trajectory (TJA) was obtained which started form the fiber canonical A-DNA conformation . Initially, we computed 10 ns of TJA and found that it sampled conformations rather dissimilar from those observed in TJBa and TJBb. A careful analysis revealed, however, certain slow structural trends and prompted us to continue TJA to 20 ns. The first two trajectories (TJBa and TJBb) have been the subject of our initial report . Therefore, below we describe in detail only TJA and use the corresponding data from TJBa and TJBb in comparisons. The structures referred to as the final MD states are the conformations averaged over the last nanosecond of the corresponding trajectory. The detailed computational protocols are described in Methods. ### All Three Trajectories Converge to Similar Structures within B-DNA Family Table I presents atom rmsd comparison between the final MD states of the tree trajectories and canonical A and B-forms of this 25-mer duplex. All computed structures are clearly different from the standard A-DNA, even though for TJA it was the starting point. Moreover, the TJA final state appears to be the less similar of the three, demonstrating that the trajectories were not trapped kinetically in the vicinities of their starting points. At the same time, the final MD states are evidently close to the canonical B-DNA. When only the central undecamer is considered, the rmsd values are in the same range as those observed in our earlier MD simulations of dodecamer duplexes . They are low, and the three computed conformations seem to form a single cluster around the canonical B-DNA. The rmsd naturally increases with the helix length, but it should be noted that the somewhat larger values obtained for the whole structures are much lower than ever observed in free MD simulations of long DNA helices . It appears, however, that, if taken as a whole, the TJA state is as close to the canonical B-DNA as to the TJBa and TJBb states, while the latter two are yet closer to each other. The kinetics of the structural convergence in terms of atom rmsd is illustrated in Fig. 1 for TJA. It is seen that the trajectory rapidly went from the initial A-DNA conformation towards the B-DNA form and, after the equilibration, the rmsd from B-DNA have already lost a half of the initial 8.6 Å. Starting from the second nanosecond it fluctuated between 2 and 4 Å. The corresponding kinetics for TJBa and TJBb were very similar except for the initial fall of the rmsd value . The rmsd from the TJBa state also shown in Fig. 1 falls down to similar final values, but exhibits a somewhat different kinetics. Namely, an overall negative drift occurred during the first ten nanoseconds followed by random fluctuations during the second half of the trajectory. One may say, therefore, that TJA first quickly traveled from A-DNA towards the B-DNA family and next slowly refined its position within this family coming closer to other computed structures. This refinement was not complete, however, since, according to Table I, the final TJA-TJB difference is still larger than that between TJBa and TJBb. Figure 1 suggests that a more accurate convergence, if possible, would require much longer time. Table II compares a few representative structural parameters of MD conformations with the corresponding standards. Already after the first nanosecond even TJA gave the helicoidals corresponding to the B-DNA family, and they exhibited no systematic change afterwards. All three final MD states have an overall bend of around 30. The bending direction is somewhat different between TJA and TJB, which is the main cause of the corresponding residual difference in terms of atom rmsd. Table II indicates that in all three trajectories both the magnitude and the direction of the bends changed significantly, and that very large variations in the bending direction apparently occurred in TJA. Thus, the slow rmsd kinetics considered above appear to be largely due to the bending of the double helices whereas the contribution from the variations of the helical parameters looks minor. Figure 2 shows the three last nanosecond average structures superimposed. They all are evidently curved, with the bends being nearly planar in each structure. In agreement with Table II, the TJA bending plane slightly deviates from the other two. The bending planes intersect the minor groove in five points which alternate between the inside and the outside edges of the bend, and in each case the the three A-tracts appear at the inside edge. The tracts are approximately phased with the helical turn, but, since the lower one is inverted with respect to the other two, its 5’ end is phased with the 3’ ends of the other. The three inside intersection points are shifted within the A-tracts from their middle towards the 3’ end of the upper two and the 5’ end of the lower one. On the other hand, the minor grooves of the two AGGC tetraplets appear at the outside edge of the curved axis, and it is readily seen that the minor groove is widened here, especially at the upper tetraplet. ### Quasi-Regular Rotation of The Bending Plane in TJA The two surface plots in Fig. 3 exhibit the time evolution of the shape of the helical axis for TJA. It is seen that the molecule was strongly bent after the initial equilibration, which was not observed in case of TJBa and TJBb . One should note that considerable initial deformation of the double helix is common for trajectories starting from the A-DNA conformation. Apparently, the molecule is stressed because the transition to the B-form occurs in these conditions during unphysically short time with much energy released. During the next few nanoseconds the bend reduced and the axis acquired a more complex shape with wound profiles in both projections. After the fifth nanosecond the bending became more planar, with much smaller curvature in Y projection. A planar bent may just mean that the helical axis is kinked in a single point or, alternatively, a lager number of local bends are properly directed. Figure 3 indicates that there is probably a mixture of these two effects. During the last few nanoseconds the axis had one stable bending point shifted upwards from the middle while another bend in the lower half emerged from time to time. The two bends were slightly misaligned, therefore, the overall bending plane rotated a little when the second bend emerged, and, in the Y projection, one sees alternation of straight and S-shaped profiles. Figure 4 displays kinetics of several quantitative measures of the magnitude of bending. The three parameters used, namely, the total angle, the shortening, and the average shift of the curved axis, all exhibit a coherent pattern of fluctuations, which locally correlates also with the rmsd from the canonical B-DNA (see Fig. 1). This indicates that they all are produced by the same motion, namely, the axis bending. Comparison of the data in Figs. 3 and 4 with similar plots earlier reported for TJBa and TJBb reveals little difference except the already mentioned initial deformation and the absence of a stable bend between the two lower A-tracts. Accordingly, TJBa and TJBb showed a somewhat stronger bending, with one-nanosecond average values usually beyond 35. In TJA, after the initial strong temporary bending, a comparable magnitude has been reached only during the last four nanoseconds. There is, however, a striking difference between TJA and the other two trajectories in the dynamics of the bend direction which is exhibited in Fig. 5. Both in TJBa and TJBb the final bending direction occurred early in the trajectories and remained quite stable although the molecule sometimes straightened producing broad scattering of points in Fig. 5. In contrast, during the first ten nanoseconds of TJA, the bending plane made almost a half turn with respect to the coordinate system bound to the molecule. It means that a transition occurred between the oppositely bent conformations, but, as seen in Fig. 3, the straight one was avoided. This rotation was very steady, almost regular. It gradually slowed down becoming indistinguishable in the last five nanoseconds. After this transition the directions of the bends in the three trajectories became much closer, and this quasi-regular motion is apparently responsible for the slow drift of the rmsd from the TJBa state in Fig. 1. The overall amplitude of this motion was around 150, that is the initial strong bend noticed in Fig. 3 was nearly opposite to that finally established. ### The Rotation of the Bending Plane is Not Energy Driven The overall character of motion revealed in Fig. 5, namely, the steady rotation of the bending plane, looks strange and counter-intuitive. A priori, we would rather expect to obtain random sampling of different bending directions, with the correct one statistically preferred due to its lower energy. The apparent quasi-regular dynamics exhibited in Fig. 4 might mean that our trajectory represents a downhill motion along a valley on a potential energy surface. Its steep borders would separate bent conformations from the straight one, with the bottom of this valley slightly inclined towards the preferred bending direction. In this case all bent conformations, including incorrect bends, should have been lower in energy than the straight one. Figure 6 displays the time evolution of the potential energy in all three trajectories. It is seen that the energy dropped during the first nanosecond and later remained stable. No clear correlation is seen between the instantaneous magnitude of bending displayed in Fig. 4 and the potential energy, therefore, one cannot say that straight states have significantly different energies than the bent ones. Neither can we claim that the preferred bending direction is characterized by lower energy values than other bends. In Fig. 6, a slight decrease in energy is observed during the second half of TJA, but it occurred when the regular rotation of the bending plane has essentially finished. On the other hand, the lowest energy during the first half of the trajectory was observed at around 3.2 ns when the bending direction was completely different. Note also that, during the first ten nanoseconds, the traces of TJBa and TJBb go above the last one, although in these cases the correct bending direction has already established. We have to conclude, therefore, that the simple energetic al explanation of the observed effect does not work. ### The Minor Groove Profiles of Converged Structures Are Similar But Not Identical The surface plot in Fig. 7 exhibits the evolution of the profile of the minor groove during TJA. The initial A-DNA conformation is characterized by a uniformly wide minor groove of 13.6 Å. It is seen that after the equilibration period the groove was much narrower, but still wider than in the canonical B-DNA model. Moreover, a complex profile have emerged with three local widenings at A-tracts, which is exactly opposite to the expectations. The two terminal widenings reduced during the first ten nanoseconds whereas the maximum of the middle one gradually shifted from its 3’ end to 5’ end. This shift evidently accompanies the rotation of the bending plane described above. One can note that the maximal widening of the minor groove moved for only 2-3 base pair steps, which is less than a 150rotation seen in Fig. 5. It appears that, in fact, the maximal widenings and narrowings in the minor groove profile do not always correspond to the direction of local bends. The initial bend was directed towards the minor groove of the upper TAGG tetraplet where the minor groove was narrowed. The two neighboring widenings are shifted by three base pairs only and they appear at the opposite sides of the bending plane which is approximately collinear to the pseudodiad axis at the center of the middle ATT triplet. In contrast, in the last structure shown in Figs. 2 and 3 the bending plane passes exactly through the maximum widening of the minor groove. The overall rotation, therefore, corresponds to approximately four base pair steps which gives the observed turn by 150. It can be noted, finally, that although Fig. 5 indicates that the bending stabilized after ten nanoseconds, the profile of the minor groove in Fig. 7 continues to evolve slowly till the very end of the trajectory. Figure 8 displays the minor groove profiles of the last average structures from the three trajectories. For TJBa and TJBb their kinetics was detailed in our first report , and we only note here that the corresponding profiles shown in Fig. 8 established during the first two nanoseconds and showed little variations afterwards . The three traces evidently exhibit a certain similarity, but do not coincide. The TJBa and TJBb grooves have the same number of local narrowings and widenings which differ slightly between the two both in amplitude and in position. The TJA profile is similar in the right-hand half of the figure. One can notice that the change from TJBa to TJBb and next to TJA involves the growing widening at the TTAG tetraplet accompanied by a shift of the secondary maximum, and looks rather regular and concerted. At the opposite half of the structure, the TJA conformation shows a narrow minor groove without significant modulations of the width. This difference may be related with the smaller magnitude of the bending in the case of TJA where the second bending point appeared from time to time only and was less significant than in the other two trajectories. ### Key Helicoidal Parameters Exhibit Consistent Regular Patterns Only after Window Averaging Figure 9 shows variation of some helicoidal parameters along the duplex in the three structures. The two inter base pair parameters, namely, roll and tilt, are most often quoted in the literature in relation to the static DNA curvature. If one first takes an ideal straight column of stacked parallel base pairs and next introduce a non-zero roll value at a certain step, the structure will bend at this step towards the major groove, if the roll is positive, and to the minor groove if it is negative. A similar experiment with the tilt value would result in bending in the perpendicular direction. It seems obvious that, whatever the physical origin of the curvature, in a bent double helical DNA, the roll and tilt values must exhibit systematic variations phased with the helical turn. Moreover, it is often assumed that for some short DNA sequences certain non-zero roll and tilt values are strongly preferred energetically, which produces static bending when they are repeated appropriately. However surprising, although all three average structures are smoothly curved, only a few supporting signs for the foregoing paradigm are readily seen in Fig. 9. For the tilt, the three traces are very dissimilar and the only feature that repeats is the alteration of its values between consecutive steps. Namely, if the tilt is low at a given step it normally goes up at the next one, and vice versa. In the three average structures, however, these alterations are sometimes oppositely phased even in TJBa and TJBb where the overall structures look particularly similar. The same is true for the roll and twist although, in these cases, some clear sequence preferences do exist. Note, for instance, that, in all four TpA steps, the roll is almost always positive and larger than in the neighboring steps. Paradoxically, two of these TpA steps occur almost exactly at the inside edge of the curved axis, that is a high positive roll accompanies the bending in an exactly opposite direction. This paradox is readily resolved when one looks at the roll values at the neighboring steps. A TpA step with a high positive roll is normally preceded or followed by a step with a low negative roll. The higher is the maximum, the lower is the neighboring minimum, so that the two nearly cancel each other. The other two TpA steps are found at the outside edge of the helical axis and their high roll probably contributes to bending. However, while a more or less repetitive pattern is observed around the third TpA step, the first one exhibits rather dissimilar pictures even for TJBa and TJBb structures which both have a widened minor groove here. Also, the roll values at the third TpA step differ considerably between the structures, but do not correlate with the bending magnitudes. The twist, tilt, and roll values used for the plots in Fig. 9 are the so called “global” parameters from the outputs of the Curves program . One may argue that they are not appropriate in the present context since they are computed by using local directions of already curved optimal helical axis. However, when “local” values are used instead, the amplitudes of the alternations in these profiles are only increased. The last plot in Fig. 9 exhibits the variations of the propeller twist. Again one sees that its value alternates between consecutive base pairs, with little phase similarity between the three structures. At the same time, in this case, a consistent repetitive pattern is evident, with strong negative propeller values in all A-tracts. These regular patterns look even more similar than the structures themselves. For instance, there is no evident difference between the three traces that would correspond to that in the minor grove profiles in Fig. 8. The apparent jumping alterations of the helicoidal parameters along the double helix naturally suggests that one should try to smooth them out by averaging the traces in Fig. 9 with a sliding window. Figure 10 shows the results of such treatment and also includes the corresponding data for the inclination which was, however, used without the smoothing. The difference between Figs. 9 and 10 is rather significant. Now all four helicoidal parameters considered in Fig. 9 exhibit regular, sometimes almost sinusoidal, oscillations. The phasing of these oscillations with the helical turn, however, is not always evident. The propeller and the inclination both exhibit approximately 2.5 periods, that is the dominating Fourier component has a wave length of approximately ten base pairs corresponding to one helical turn. For the roll, the dominating wave length apparently corresponds to 5-6 base pairs, that is a half of a helical turn. When different structures are compared, however, it is seen that only for propeller the maxima and minima coincide well. A more complex pattern is observed for the twist, and one can notice a correlation between its traces and the minor groove profiles shown in Fig. 8. Namely, the twist is lower in the widenings of the groove and higher in its narrowings. These results suggest that the relationship between the helicoidal parameters and the bending is complex and cannot be reduced to simple models of roll-like or tilt-like bends outlined above. Accumulation of the regular variations revealed in Fig. 10 probably gives the correct overall bend angles and directions, but neither can be easily predicted just by looking at these traces. ### The Distributions of $`𝐁_𝐈\mathrm{𝐚𝐧𝐝}𝐁_{\mathrm{𝐈𝐈}}`$ Backbone Conformers in Bent Structures are Surprisingly Dissimilar In all three trajectories, dynamics of $`\mathrm{B}_\mathrm{I}\mathrm{B}_{\mathrm{II}}`$ backbone transitions was qualitatively similar in a few aspects. Consider Fig. 11a, for instance, where the results are shown for TJA. The overall pattern reveals rather slow dynamics, suggesting that MD trajectories in the 10 ns time scale are not long enough to sample all relevant conformations. A somewhat higher $`\mathrm{B}_\mathrm{I}\mathrm{B}_{\mathrm{II}}`$ activity was observed during the first half of the trajectory, when the rotation of the bending plane occurred. It is seen that, in A-tracts, the B<sub>II</sub> conformers are preferably found in ApA steps and that they tend to alternate with B<sub>I</sub> within the same strand. There are many examples of concerted $`\mathrm{B}_\mathrm{I}\mathrm{B}_{\mathrm{II}}`$ transitions, when a given step switches from $`\mathrm{B}_{\mathrm{II}}\mathrm{to}\mathrm{B}_\mathrm{I}`$ simultaneously with an opposite transition in one of the neighboring steps. Sometimes three consecutive steps are involved and, less often, the opposite strand as well. Many $`\mathrm{B}_\mathrm{I}\mathrm{B}_{\mathrm{II}}`$ transitions are reversed within hundreds of picoseconds, but there are also very long-living conformers and sites where either $`\mathrm{B}_\mathrm{I}\mathrm{or}\mathrm{B}_{\mathrm{II}}`$ states are preferred. A strong preference of $`\mathrm{B}_\mathrm{I}`$ state is observed for all TpT steps, for example. However, it seems to be the only case when the effect repeats at a base pair step level. In a few steps where the $`\mathrm{B}_{\mathrm{II}}`$ conformation is preferred this is apparently determined by a broader sequence context. The corresponding data for TJBa and TJBb were included in our first report and they revealed the same qualitative features . It was very surprising for us, however, that, in spite of the good convergence in terms of the overall bent shape of the molecule, the three trajectories gave rather dissimilar distributions of $`\mathrm{B}_\mathrm{I}\mathrm{and}\mathrm{B}_{\mathrm{II}}`$ conformers along the sequence. Fig. 11b compares these distribution in the final backbone conformations in the three trajectories. There are 14 non TpT steps where the $`\mathrm{B}_\mathrm{I}`$ conformation is found in all three structures. However, since our trajectories started from $`\mathrm{B}_\mathrm{I}`$ states, this number hardly tells us something. On the other hand, the number of $`\mathrm{B}_{\mathrm{II}}`$ conformers found in each structure and in each strand is similar and roughly corresponds to 25% of phosphate groups. Assuming that the $`\mathrm{B}_{\mathrm{II}}`$ states are evenly distributed in the sequence one gets the expectation value of 0.75 for the number of cases when the $`\mathrm{B}_{\mathrm{II}}`$ conformer should be found in the same base pair step in all three structures. The observed number of such sites is three. Note, however, that they all are found in A-strands of A-tracts where, as noted above, the $`\mathrm{B}_{\mathrm{II}}`$ conformers tend to alternate. This, together with the strong preference of TpT steps for $`\mathrm{B}_\mathrm{I}`$, increases the probability of matching. These results suggest that the relationship between the bending of the DNA double helix and the $`\mathrm{B}_\mathrm{I}\mathrm{B}_{\mathrm{II}}`$ backbone transitions, if any, is loose in the sense that a given bent shape does not impose a fixed $`\mathrm{B}_{\mathrm{II}}`$ distribution upon the backbone. ## Discussion This study gives the first example of a successful implementation of the general strategy outlined in Theoretical Background. Namely, we showed that the minimal model of B-DNA, which is biased only by the nucleotide sequence, in dynamics, reproducibly converges to a single state characterized by an ensemble of similar statically bent conformations. The effect has been demonstrated here for one sequence only. Moreover, this sequence was specifically constructed rather than taken from experimental studies. Nevertheless, the sequence motive A<sub>n</sub>TAG used in construction was found in the center of the first bent DNA fragment studied experimentally . In addition, the character of bending in the computed conformations, notably, its direction with respect to the A-tracts, and modulations of the groove width, qualitatively agree with the rules derived from experiments. These observations validate an attempt to make the next step of the above strategy, namely, below we try to disclose the mechanism responsible for the bending within the framework of the minimal model. We believe that, in spite of the obvious limitations of this model, its main features responsible for the bending correspond to reality. At the same time, the real situation is certainly more complex. ### Results Poorly Agree with Earlier Theories of Bending All theories proposed during the last 20 years to explain intrinsic bends in DNA double helices agree with some experimental observations and disagree with the other and, probably, each of them continues to attract some proponents. Here we compare our results with the most popular models of bending regardless of their experimental validation. Comparisons with experimental data have been the subject of many reviews . The wedge model of DNA bending resulted from merging of ideas developed in seventies to explain the ability of a double helix to wrap around nucleosome particles. The first idea was that this can occur due to kinks of the helical axis phased with the helical screw , with kinks implying destacking of base pairs in fable points in order to maintain perfect stacking elsewhere. The second idea was that the double helix can be smoothly bent, without destacking, by small deformations in every base pair step . The wedge model merges the two by postulating that, in every specific dinucleotide, the preferred stacking of bases is slightly non-parallel and this causes bending in the same way as kinks do. It can be further developed by increasing the number of wedge degrees of freedom, by considering triplets, tetraplets, and so forth instead of dinucleotides, and by assuming that the non-zero average wedges result from random sampling from asymmetrical energy valleys around local energy minima, rather than from minimum energy configurations . Depending upon the specific wedge parameters, this model can place the curvature inside A-tracts or between them and also explain bending in non A-tract sequences . For the present discussion, it is convenient to unite all such mechanisms in one group characterized by the tacit emphasis upon the specific base pair stacking preferences as the source of the DNA bending. The results shown in Figs. 9 and 10 obviously disagree with these views. There is little similarity between matching dinucleotides in the same structure and, moreover, base pair steps put in the same sequence context in three closely similar bent conformations exhibit broadly different helical parameters. The last observation means that even a generalized wedge model with dinucleotide blocks replaced by triplets, tetraplets, and so forth, would disagree with our results. The junction model postulates that there is a distinct specific A-tract form of the double helix characterized by a stronger inclination of base pairs with respect to the helical axis than in the normal B-DNA. In this case, planar stacking at the junction between the two structures would result in a kink of the helical axis. Formally, such geometry can also be obtained with the generalized wedge model above, but the junction theory puts an emphasis upon the specific A-tract form of DNA as the principal physical cause of bending. Its structure can be due to cooperative interactions in long DNA stretches and its environment. Within the framework of the junction model, particular roles were sometimes attributed to bifurcated hydrogen bonds , the water spine in the minor groove , or the NH<sub>2</sub> groups of adenines . It is evident that the junction model also poorly agrees with our results. In dynamics, conformations, both smoothly bent and kinked at the two insertions between the A-tracts, are observed periodically. The kinks, however, are not centered at the boundaries between A-tracts and the flanking sequences, and they are not sharp. In such conformations, A-tracts are less bent than regions between them, that is the bend is localized but still smooth. In Fig. 10 the inclination shows smooth oscillations, even without window averaging, with no kinks. It decreases from 5’ to 3’ ends of A-tracts, and since the 3’ ends of A-tracts are dephased and positioned differently with respect to the bending plane, no evident relationship to bending can be readily seen. All helical parameters vary along the sequence so that there is no A-tract fragment where they repeat at two consecutive steps. Thus, although the structures are bent, the specific regular A-tract structure is not seen, as well as the “random B-DNA”, though, which are the two key components of the junction model. The third model, which was first mentioned in the context of the junction theory , but became popular only in the recent years , attributes the cause of bending to solvent counterions. If they are specifically bound by minor grooves of A-tracts, in a phased sequence, phosphate groups would appear partially neutralized at one side of the double helix, and the repulsion between the opposite phosphates would bend DNA towards minor grooves of A-tracts. The very fact that the minimal model of B-DNA, without explicit counterions, produces static bends, in good agreement with experiments, strongly suggests that the counterions hardly play a key role in the A-tract induced bending. At the same time, our results do not contradict less specific non-local theories of A-tract bending. The modified junction model assumes that the deformations at the boundary between the two conformations can propagate for several base pair steps. The A-tracts in the sequence studied here may be too short for their ingenious structure to establish. The second such theory proposed that the bending is caused by the modulations of the minor groove. Really, the double helix is usually bent towards the major groove at the minor groove widenings, and in the opposite direction at its narrowings. In TJA, for instance, this relationship is maintained during the rotation of the bending plane. Sometimes, however, the double helix straightens and remains unbent during nanoseconds, while the minor groove profile does not change . It is understood, however, that the last two non-local models are incomplete. Actually, they cannot be verified or disproved because the issue of the physical origin of bending is tacitly dropped. Simple geometrical considerations dictate that the grooves must be narrower at the inside edge of the bend . One may postulate, therefore, that groove modulations cause bending or, vice versa, that it is bending that causes groove modulations, but the physical origin of the phenomenon remains obscure. Similarly, the modified junction model essentially discards the essence of the original theory, which considers the specific poly-dA structure as the source of the bend. If the “boundary deformations” can exist without the structures and boundaries themselves one should look for another force that maintains these deformations. Generally speaking, the results presented here are best interpreted if one assumes that there is an external force that imposes a bent shape upon the double helix as a series of mechanical constraints. The double helix is allowed to move, but so that these constraints would remain fulfilled. Thus, the bases can change their mutual orientation and the backbone can switch between $`\mathrm{B}_\mathrm{I}\mathrm{and}\mathrm{B}_{\mathrm{II}}`$ conformations, but the overall proportion of the $`\mathrm{B}_{\mathrm{II}}`$ conformers remains constant, and fluctuations of helical parameters in the neighboring base pair steps tend to compensate. ### Possible Physical Origin of Spontaneous Static Bends in Double Helices The hypothesis outlined below is based upon our computational results as well as upon analysis of well-known experimental data. Although it does not answer all unclear questions concerning DNA bending we consider it most likely and describe it here for discussion and further investigation. Let us first ask the following simple geometric question: “What is the shortest line that joins two points on a surface of a cylinder?” To answer it, one should first cut the cylinder parallel to its axis, unfold its surface onto a plane, join the two points by a straight line and then fold the surface back upon the cylinder. The resultant curve represents an interval of a spiral trace with a constant inclination to the cylinder axis. Now consider an ideal canonical B-DNA model of a double helix. The stacked base pairs form the core of a cylinder and the sugar-phosphate backbone forms an ideal spiral trace on its surface, that is the shortest line that joins the “surface” nitrogens $`\mathrm{N}_1/\mathrm{N}_9`$ of the bases. If we now assume that the backbone is a stiff elastic that can be characterized by a certain specific length, we are obliged to conclude that this model implicitly assumes that the backbone is stretched and tends to reduce its length on the surface of this cylinder. Our last question is: “What would happen if the preferred backbone length appeared to be longer that allowed by the canonical model?” A simple answer is: it would try to extend by pushing bases. They can accommodate this extention within the framework of a regular helical structure by increasing the helical twist and changing other helical parameters. This option, however, is opposed by the loss in the stacking energy and, when it becomes difficult to extend in this way, the backbone will try to deviate from the ideal spiral trace. In this case the the grooves on the surface of the Watson-Crick double helix can no longer maintain a constant width. It seems possible to assume that, in physiological conditions, the equilibrium specific length of the DNA backbone is actually larger than that allowed by a regular B-DNA structure with the average helical twist of 34.5. Its further extention is opposed by the limit of the tolerance of the pase pair stacking and, as a result, the backbone appears “frustrated” and is forced to wander on the cylindrical surface formed by base pair stack. The ideal parallel stacking has to be perturbed and we believe that it is this effect that eventually bends the double helix. Modulations of the DNA grooves, which is a well-known ubiquitous feature of the single crystal B-DNA structures, is a natural indicator of this particular state of the backbone. It is observed for very different sequences as an apparently general attribute of the B-DNA form in physiological conditions. Thus, if we had to decide whether the DNA backbone in stretched, relaxed or compressed by looking only at the single crystal B-DNA structures, we would be obliged to conclude that the first option looks unlikely, the second is possible, while the third is the most probable. A compressed backbone is more likely to cause smooth groove modulations found in experimental structures than a relaxed one, which would rather be controlled by the local sequence effects. Let us consider an ideal B-DNA model, with planar base pairs perpendicular to the helical axis, and try to imagine how wandering of the backbone traces can emerge. For simplicity, we first consider the helical twist as the only variable parameter. Obviously, by smoothly increasing and decreasing the twist we obtain, respectively, narrowings and widenings of the minor groove. The desired backbone waving emerges, and a larger its length can be accommodated on the same cylindrical surface. It is easy to see, however, that, if the parallel stacking is maintained, the backbone must be compressed when the twist is reduced and stretched in the opposite phase. In reality, however, the backbone is stiff and it cannot be compressed significantly, therefore, it is the stacking that suffers when the twist is reduced. Although this description is simplistic, and other base pair degrees of freedom in addition to the twist can contribute to the wandering, it captures the essence of the underlying mechanics. In the widenings of the minor groove, where the twist is reduced, the backbone pushes off the neighboring base pairs at C1’ atoms, causing various perturbations of the parallel stacking. On average, they are likely to deviate the helical axis towards the major groove because C1’ atoms are at the minor groove side. These perturbations are delocalized and involve rolling, tilting, as well as other relative motions of base pairs, and there is an ensemble of orientations that fulfill the constraints imposed by the backbone lengths, rather than a single preferred bent conformation. The modulations of the minor groove width and the bending of the double helix appear related, as was suggested earlier , because they represent two consequences of a single cause. They are related as brothers rather than as a parent and a child and, probably, are not bound to always appear together. The major component of the backbone stiffness is the electrostatic repulsion between the charged phosphate groups. Even though this repulsion is shielded by water and counterions, the experiments where bending in B-DNA was induced by specific neutralization of phosphates proved that they are not shielded even when separated by two helical turns. Complete neutralization, therefore, is hardly imaginable. The local electric field around a pair of neighboring phosphates in the same strand is created by all surrounding charges, including the phosphates of the opposite strand, and it favors maximal possible separation between $`\mathrm{P}_\mathrm{n}\mathrm{and}\mathrm{P}_{\mathrm{n}+1}`$. In B-DNA, this distance is close to the absolute maximum, which is achieved by putting all relevant backbone torsions except one in the trans configuration . The maximal extention gives the ground energy state with the $`\mathrm{P}_\mathrm{n}\mathrm{P}_{\mathrm{n}+1}`$ distance around 7.7 Å. The corresponding thermodynamic average for a free backbone in solution is $`D(T,ϵ)<7.7`$ Å, where $`T`$ is the temperature, and $`ϵ`$ is the effective dielectric constant that depends upon the concentration of counterions. In the single crystal structures, the largest $`\mathrm{P}_\mathrm{n}\mathrm{P}_{\mathrm{n}+1}`$ distances observed are in the range of 7.3 – 7.6 Å suggesting that there are no prohibitive steric obstacles for a completely extended backbone. At the same time, the distances most commonly found are around 6.7 Å while in the fiber canonical structure it is 6.5 Å. Apparently, with normal temperature in physiological conditions $`D(T,ϵ)6.7`$, and the backbone is forced to wander. $`D(T,ϵ)`$ should be a decreasing function of both arguments, therefore, the backbone stiffness and, accordingly, the curvature should decrease as the temperature grows and/or as the phosphate shielding is improved by increasing the ionic force. These two non sequence-specific effects have been found in experiments . With $`D(T,ϵ)6.5`$ the backbone relaxes and the phenomenon of DNA bending vanishes. According to this hypothesis, the transition of A-tracts in a specific DNA form is not an indispensable prerequisite of the bending. Moreover, it suggests that the regular poly-dA structure may not exist in solution because, with the average twist increased to 36, the backbone is, possibly, still compressed and continues to wander, although with a longer characteristic wave lengths. The structures of short A-tracts computed here are rather variable and it is not clear how they can be extended to longer poly-dA double helices. We believe that the A-tracts rather label the regions where higher twist values are allowed by the base stacking interactions. The backbone prefers to compress the minor groove here, thus fixing the phase of its modulations along the double helix. In random and homopolymer sequences the minor groove probably also narrows and widens, but the corresponding maxima and minima can migrate along the double helix in a way similar to that observed here during the rotation of the bending plane in TJA. Our model considers the bending of a DNA double helix as a deformation imposed upon the stem of the stacked base pairs by interactions external to it. The bases are forced to “forget” their preferred stacking orientations and look for a possibility to maintain the overall structure by sampling the orientations at the limit of destacking. At the same time, it is the broad “tolerance” of the base pair stacking that makes all this game possible. If true, this theory gives a slightly different overall view of the DNA molecule in physiological conditions and entails important biological consequences. Notably, it suggests that the local DNA structure is not simply determined by the stacking preferences of base pairs in dinucleotides, trinucleotide, and so forth. The two waving backbone profiles on the surface of the helix impose a medium range context upon the local base pair stacking because the phases of these modulations can well correlate over several DNA turns. This makes possible mutual dependence of local conformations in sites separated by considerable DNA stretches. Fine tuning of phases of these modulations may be the function of single short A-tracts as well as of some regulatory proteins. The degree of the backbone compression is connected with supercoiling and can be controlled in this way, which gives yet another possible instrument of structural regulation. One may note also that this theory offers a unified model which explains static bends in A-tract and non A-tract sequences as well as the bending induced by the negative supercoiling in circular DNA. ### Conclusions The static curvature spontaneously emerges in free MD simulations of an atom level molecular model of B-DNA double helix, with the nucleotide sequence as a single structural bias. Convergence to similar statically bent states have been demonstrated in three independent MD trajectories of 10-20 ns. The bending direction and its magnitude are in good agreement with experimental observations. Unexpectedly, the three computed MD structures exhibit a striking microscopic heterogeneity as regards variations of helical and conformational parameters along the molecule. Based upon the computational results as well as the literature experimental data a new possible mechanism of bending in a double helical DNA is proposed. It postulates that, in physiological conditions, the equilibrium specific length of the DNA backbone is larger than is admissible in the regular B-DNA form, which forces it to fold in a wavy trace on the cylindrical surface of the double helix. This results in modulations of the minor groove width, slight asymmetrical destacking of base pairs at the groove widenings and, eventually, in bending of the DNA molecule. ### Methods Molecular dynamics simulations have been performed with the internal coordinate method (ICMD) including special technique for flexible sugar rings . The so-called minimal B-DNA model was used which consists of a double helix with the minor groove filled with explicit water. It does not involve explicit counterions and damps long range electrostatic interactions in a semi-empirical way by using linear distance scaling of the electrostatic constant and reduction of phosphate charges. The DNA model was same as in earlier reports, namely, all torsions were free as well as bond angles centered at sugar atoms, while other bonds and angles were fixed, and the bases held rigid. Molecular dynamics calculations were carried out with a time step of 10 fsec and the effective inertia of planar sugar angles increased by 140 amu$``$Å<sup>2</sup> as explained elsewhere . The coordinates were saved once in 2.5 ps. AMBER94 force field and atom parameters were used with TIP3P water and no cut off schemes. The initial conformation for TJBa was prepared by vacuum energy minimization starting from the fiber B-DNA model constructed from the published atom coordinates . The subsequent hydration protocol to fill up the minor groove normally adds around 16 water molecules per base pair. The starting state for TJBb was obtained by energy minimizing an equilibrated straight structure taken from the initial phase of TJBa. The initial conformation for TJA was prepared by hydrating the minor groove of the corresponding A-DNA model without the preliminary energy minimization. In TJA, the necessary number of water molecules was added after equilibration to make it equal to that in TJBa and TJBb. The heating and equilibration protocols were same as before . During the runs, after every 200 ps, water positions were checked in order to identify those penetrating into the major groove and those completely separated. These molecules, if found, were removed and next re-introduced in simulations by putting them with zero velocities at random positions around the hydrated duplex, so that they could readily re-join the core system. This procedure assures stable conditions, notably, a constant number of molecules in the minor groove hydration cloud and the absence of water in the major groove, which is necessary for fast sampling . The interval of 200 ps between the checks is small enough to assure that on average less then one molecule is repositioned and, therefore, the perturbation introduced is considered negligible. ## Appendix This section contains comments from anonymous referees of a peer-review journal where this and a closely related paper entitled “Molecular Dynamics Studies of Sequence-directed Curvature in Bending Locus of Trypanosome Kinetoplast DNA” has been considered for publication, but rejected. ### A Journal of Molecular Biology #### 1 First referee These companion manuscripts describe a series of molecular dynamics trajectories obtained for DNA sequences containing arrangements of oligo dA - oligo dT motifs implicated in intrinsic DNA bending. Unlike previous MD studies of intrinsically bent DNA sequences, these calculations omit explicit consideration of the role of counterions. Because recent crystallographic studies of A-tract-like DNA sequences have attributed intrinsic bending to the localization of counterions in the minor groove, a detailed understanding of the underlying basis of A-tract-dependent bending and its relationship to DNA-counterion interactions would be an important contribution. Although the MD calculations seem to have been carried out with close attention to detail, both manuscripts suffer from some troubling problems, specifically: The DNA sequence in question is a 25-bp deoxyoligonucleotide that contains 3 A/T tracts. Two of these are arranged in phase with the helix screw with the third tract inverted with respect to the other two. Extrapolating from available experimental data, this sequence is expected to confer some degree of intrinsic bending. The main focus of this manuscript is the comparison of data obtained for an MD trajectory computed from an A-form starting conformation (TJA) with two other trajectories that begin with B-form structures (TJBa and TJBb). Significant differences in behavior and in time-averaged helical parameters are observed for the TJA trajectory compared with both TJBa and TJBb, suggesting that the structures are not fully equilibrated. This is particularly evident in the computed bending direction, which varies dramatically during early times in the TJA trajectory. Even after 15 ns, when the orientations of bending planes appear to have approached asymptotic values, the TJA plane is displaced by between 30 and 60 degrees from those of TJBa and TJBb, which are quite similar to one another. This fact strongly suggests that the MD-simulation results depend nontrivially on initial conditions, even after 15-20 ns, which calls into question most of the results obtained from the computed trajectories. #### 2 Second referee Dr. Mazur reports the results of MD simulations of DNA 25-mers sequences that contain three phased A-tracts. He believes that he has obtained the first model system in which properly directed static curvature emerges spontaneously in conditions excluding any initial bias except for base pair sequence. He observes that the ensemble of curved conformations reveals significant microscopic heterogeneity, which he believes is in contradiction to existing theoretical models of DNA bending. In CAM110/00 he performs a series of simulations on a DNA fragment that has not been shown experimentally to bend in solution. In this case the DNA sequence was chosen based on its propensity to adopt a characteristic structure during simulations. In CAM167/00 he performs a similar investigation on B-DNA fragment composed of a sequence that has been shown experimentally to bend. My view is these two papers should be combined as one, and the review will treat them as one. I found this paper to be interesting and possibly worthy of publication in JMB, even as I took issue with a substantial portion of it. The basic premise of the paper is that a model lacking realistic electrostatics can provide meaningful information about long range DNA conformation. In Mazur’s model, long range electrostatic interactions are dampened and phosphate charges are attenuated. I had some concerns about the basic rationale for the non-electrostatic model. I initially assumed that it’s greater simplicity would allow longer trajectories, etc. But the trajectories of Dr. Mazur are not substantially longer than those described by Beveridge, Pettitte, etc. And in fact that seems not to be the rationale. Dr. Mazur believes that full atom force fields, with explicit ions, give less realistic results than his electrostatic-attenuated model. In particular he says that full atom force fields give slightly overwound DNA, which camouflages DNA bending. What is the cause of this? A problem in the force field? Why not fix that instead of going the non-electrostatic route? I am not comfortable enough with the world of MD pass judgment on this issue, but think someone who is should evaluate that prior to publication. I just went back and re-read Diekmann’s classic 1985 JMB paper \[Diekmann, S., & Wang, J. C. ”On the sequence determinants and flexibility of the kinetoplast DNA fragment with abnormal gel electrophoretic mobilities” (1985) J. Mol. Biol. 186, 1-11.\] Diekmann shows clearly that the electrophoretic anomaly of kinetoplast DNA decreases with increasing Na, and increases very dramatically with increasing Mg. His experiments seem well-conceived , well-conducted and well-analyzed. For example he implants a temperature sensor within his gels to insure constant, fixed temperature. One has to believe Diekmann’s results. My fundamental problem with Mazur’s model is that it cannot account for experimental data. How can bending be cation-dependent, but the mechanism not be electrostatic in nature? Mazur does concede that experimentally ”curvature is reduced in high salt, but for some sequences it is increased in the presence of divalent metal ions” (cites Diekmann). \[page 2 MS CODE CAM110/00\]. But the implication here is that the observation of Diekmann is not general to all A-tracts. The next sentence of the manuscript may be read as confirming that the cation effect is not general, but is length and composition dependent (the text is a little confusing here). However the Woo and Crothers citation, used as support for that, does not discuss the cation effect. If there are data somewhere suggesting that the cation effects are not general, they should be cited and discussed. That would really increase the strength of Mazur’s argument. If not, the text should be clarified. An additional issue that is not illuminated much here is the comparison of Mazur’s model with the results of x-ray crystallography. In crystals of oligonucleotides, A-tracts are straight (”less prone to bending than other sequences” is rather understating it). To accommodate this observation, Dickerson (JMB 1994) proposed a model in which A-tract DNA curvature results from roll-bending of non-A-tracts, and linear A-tracts. Crothers (JMB 1994) is contemptuous of that model, and believes that the linear conformations of A-tracts observed thus far in crystals are not those associated with the curvature ’observed’ in gel mobility experiments. In fact such a discrepancy between dilute solution (where intramolecular forces would dominate) and condensed states (such as crystals, where intermolecular forces dominate) is expected if long-range electrostatics play a key role in curvature. Those long range forces are turned off in Mazur’s model. (He does seem to allude to the crystallography/solution discrepancy on page 19). So again his model does not account for experimental data. As an aside: Mazur believes that groove narrowing and bending are coupled. How does one then explain the observation that A-tracts in crystals have narrow minor grooves, yet are not bent? Finally, some aspects of Mazur’s (combined) model seems to be inconsistent and self-contradictory. In his model (as I understand it), (1) electrostatic repulsion between adjacent phosphate groups drives helical twisting, (2) A-tracts are regions where higher helical twist is facilitated by lower stacking energies in comparison to those of G-C base pairs, (3) higher helical twist narrows the minor groove, and (4) groove narrowing is somehow related to axial curvature (this is a little unclear; the description ”as brothers rather than a parent and a child” did not enlighten me). This model has certain attractive features, \[the idea that electrostatic repulsion between adjacent phosphate groups drives helical twisting while stacking opposes it was previously presented by Alex Rich in 1992 in a chapter of Structure & Function, Volume I: Nucleic Acids pp. 107-125 (from a Sarma meeting)\] but some deficiencies also. If electrostatic repulsion between adjacent phosphate groups drives helical twisting, then how can correct values of helical twist be obtained with attenuated phosphate charges? Or restated: Does this model not ascribe electrostatic forces as the ultimate cause of static bends, contradicting the non-electrostatic assumption? And I just checked in one crystal structure and found a place where OP to OP (phosphate oxygens, where the negative charge resides) across the minor groove are less than those between adjacent phosphates. How can electrostatic repulsion between adjacent phosphate groups drive other phosphate groups together like that, especially if stacking forces are working in opposition? How can one understand such phenomena without explicitly considering electrostatic interactions? Although the bulk of this review might appear rather critical, a model can be useful even if it does not account for all data. And that may be the case here. If Mazur has indeed obtained the first model system where properly directed static curvature emerges spontaneously, then his model clearly has utility. If a reviewer who specializes in MD simulations (not this reviewer) would confirm that, and support the utility of the approach, then publication may be in order. However I would like the paper more if it were reformulated as an exploration of possible models rather than the last word on the physical origin of intrinsic bends. Re: measurement of the groove width: Is the some reason that an old version of Curves was used? The newer versions fit a surface to the groove, rather than just measure phosphate-phosphate distances, and provide a much finer view of groove width.
warning/0004/hep-th0004105.html	ar5iv	text	# On gauge fields - strings duality as an integrable system. ## 1 Introduction and summary. According to , , there is a semiclassical evidence, that the expectation value of a QCD Wilson loop in dimension 4 is given by a partition function of a string theory with the 5 dimensional constant negative curvature target space and with a boundary condition on the loop at the absolute, as the partition function satisfy the loop equation with respect to variation of the boundary loop. Semiclassically, the equations of motion are the equations of the minimal surface in the constant negative curvature space with prescribed boundary; and the action is the area of such minimal surface (we sum over minimal surfaces when the solution is not unique). In the paper, we present the equations of motion for such string theory in a geometrical format, which reveals, that those equations are integrable, since they can be written via a Lax pair with a spectral parameter, and therefore, the methods of boundary integrable models might be applicable, either modern ones like boundary inverse scattering , ; boundary r-matrices, the bootstrap, etc ; or rather more geometrical ones, as they used more then 100 years ago to solve nonlinear partial differential equations such as Liouville, Monge Ampere, and sine Gordon. We also show that the issue of the loop equation being satisfied rests solely on the properly posing the boundary condition, since the terms in the second variation containing $`\delta `$ function is given by a line integral over the boundary. Some other geometric and integrability aspects of the methods used in the paper will be published elswhere, , . ## 2 Moving Frames; surfaces in $`H_{(n)}`$; minimal surfaces. ### 2.1 Frames in $`M^{n+1}`$ Let $`M^{n+1}`$ be Minksowsky space: a real vector space with the metric $`\eta =(1,1,1,\mathrm{}1)`$. A frame in $`M^{n+1}`$ is an $`(n+1)`$-tuple of vectors $$\begin{array}{c}F=(f_{(0)},f_{(1)},f_{(2)},\mathrm{}f_{(n)}),\text{such that }\hfill \\ <f_{(0)},f_{(0)}>=1,<f_{(\alpha )},f_{(\alpha )}>=1,\alpha =1,2,\mathrm{}n.\hfill \end{array}$$ (1) A standart frame is the following $`(n+1)`$ tuple of vectors: $$\begin{array}{c}e_{(0)}=(1,0,0,\mathrm{}0)\hfill \\ e_{(1)}=(0,1,0,\mathrm{}0)\hfill \\ \mathrm{}\hfill \\ e_{(n)}=(0,0,0,\mathrm{}1)\hfill \end{array}$$ Any frame is obtained from the standart frame by an action of the group $`SO(1,n)`$, $`f_{(i)}=g_{ij}e_{(j)}`$, or in components, $$\left(\begin{array}{ccccc}f_{(0)}^0& f_{(0)}^1& f_{(0)}^2& \mathrm{}& f_{(0)}^n\\ f_{(1)}^0& f_{(1)}^1& f_{(1)}^2& \mathrm{}& f_{(1)}^n\\ \mathrm{}& & & & \\ f_{(n)}^0& f_{(n)}^1& f_{(n)}^2& \mathrm{}& f_{(n)}^n\end{array}\right)=\text{G}\left(\begin{array}{ccccc}1& 0& 0& \mathrm{}& 0\\ 0& 1& 0& \mathrm{}& 0\\ \mathrm{}& & & & \\ 0& 0& 0& \mathrm{}& 1\end{array}\right),$$ (2) since ( 2) defines a matrix G with $`\mathrm{G}\eta \mathrm{G}^T=\eta `$, $`\eta =Diag(1,1,1,\mathrm{}1)`$, and therefore, $`𝐆\text{ SO }(1,n)`$. ### 2.2 The hyperboloid $`H_{(n)}`$ Since $`<f_{(0)},f_{(0)}>f_{(0)}^{0}{}_{}{}^{2}+f_{(0)}^{1}{}_{}{}^{2}+\mathrm{}f_{(0)}^{n}{}_{}{}^{2}=1`$, the vector $`f_{(0)}`$ can be identified with a point on a hyperboloid $`H_{(n)}`$ in $`M^{n+1}`$, $$(x_0)^2+(x_1)^2+(x_2)^2+\mathrm{}(x_n)^2=1$$ (3) and for a fixed $`f_{(0)}`$, the tangent space to the hyperboloid at the point $`f_{(0)}`$ is spanned by $`\{f_{(1)},f_{(2)},\mathrm{}f_{(n)}\}`$; indeed, for any $`\alpha =1,2,\mathrm{}n`$, $$<f_{(0)}+ϵf_{(\alpha )},f_{(0)}+ϵf_{(\alpha )}>=<f_{(0)},f_{(0)}>+2ϵ<f_{(0)},f_{(\alpha )}>+ϵ^2<f_{(\alpha )},f_{(\alpha )}>=1+ϵ^2,$$ as $`<f_{(0)},f_{(\alpha )}>=0`$. Suppose at each point of a hyperboloid a frame $`F`$ is choosen, with $`f_{(0)}`$ corresponding to the point itself, and other $`f_{(\alpha )}`$ choosen arbitrary in the tangent space to the hyperboloid; obviuosly it is always possible to do. Since $`𝐟=𝐆𝐞`$, ( which is components is $`f_{(i)}={\displaystyle \underset{j}{}}g_{ij}e_{(j)}`$ ), $$d𝐟=d𝐆𝐞=\left(d𝐆𝐆^1\right)𝐟\omega 𝐟,$$ or $`df_{(i)}=\omega _{ij}f_{(j)}`$, where $`\omega _{ij}=\left(d𝐆𝐆^1\right)_{ij}`$; therefore, $`d\omega =d\left(d𝐆𝐆^1\right)=d𝐆𝐆^1d𝐆𝐆^1=\omega \omega `$, and we arrived at the Maurer-Cartan equations: $$\begin{array}{c}d𝐟=\omega 𝐟\hfill \end{array}$$ (4) $$\begin{array}{c}d\omega =\omega \omega \hfill \end{array}$$ (5) In fact, if ( 5) is satisfied, then equations ( 4) are comptaible. It’s easy to see that ( 5) is necessary, as if $`d𝐟=\omega 𝐟`$, it follows that $`0=d(\omega 𝐟)=(d\omega \omega \omega )𝐟)`$. Therefore, the crucial part in what follows will be to construct one forms satisfying ( 5), as then the frame can be found by integrating the compatible first order equations ( 4). From ( 1), it follows that $$\begin{array}{c}\omega _{0,0}=0\hfill \\ \omega _{0,\alpha }=\omega _{\alpha ,0}\hfill \\ \omega _{\alpha ,\beta }=\omega _{\beta ,\alpha },\alpha ,\beta =1,2,\mathrm{}n.\hfill \end{array}$$ (6) The induced metric on the hyperboloid is just $$h:=<df_{(0)},df_{(0)}>=\omega _{0,i}\omega _{0,j}<f_{(i)},f_{(j)}=\omega _{0,i}^{}{}_{}{}^{2}$$ (7) Since the curvature 2-form on the hyperboloid is $$\mathrm{\Omega }_{\alpha ,\beta }=d\omega _{\alpha ,\beta }\omega _{\alpha ,\gamma }\omega _{\gamma ,\beta },\alpha ,\beta ,\gamma =1,2,\mathrm{}n,$$ and from Maurer Cartan ( 5) it follows that $$\mathrm{\Omega }_{\alpha ,\beta }=\omega _{0,\alpha }\omega _{0,\beta }$$ It is easy to see that the hyperboloid has constant negative curvature in this language: choose the basis of tangent vectors $`X_i`$, $`i=1,2,\mathrm{}`$, such that $`\omega _{0,i}(X_j)=\delta _{ij}`$. Then such vectors are ortonormal in the induced metric, $`h(X_i,X_j)=\delta _{ij}`$. The Riemann tensor in this basis is $`R_{ijkl}=\mathrm{\Omega }_{ij}(X_k,X_l)`$. Contracting two indices, $`R_{kikj}=\mathrm{\Omega }_{ki}(X_k,X_j)=\omega _{0,k}\omega _{0,i}(X_k,X_j)=(n1)\delta _{i,j}=(n1)h(X_i,X_j)`$. ### 2.3 A surface in $`H_{(n)}`$ Since $`f_{(0)}`$ is identified with a point on a hyperboloid, a vector- valued function of 2 real varibles $`f_{(0)}(u,v)`$ defines a surface on the hyperboloid. At each point on the surface, we choose the frame $`F(u,v)`$ , ( 1), in such a way, that $`f_{(1)}`$ and $`f_{(2)}`$ will span the tangent space of the surface (of course, there are many ways to do it); as before, all $`\{f_{(i)}(u,v)\},i=1,2,\mathrm{}n`$ span the tangent space at $`f_{(0)}`$, $`T_{f_{(0)}}`$ in the hyperboloid $`H_{(n)}`$. With this choise, $$df_{(0)}(u,v)=\left(\omega _{01}f_{(1)}+\omega _{02}f_{(2)}\right)(u,v)$$ (8) and $$\omega _{0\mu }=0,\mu =3,4,\mathrm{}n;$$ (9) From Maurer-Cartan also $`d\omega _{0\mu }=\omega _{01}\omega _{1\mu }+\omega _{02}\omega _{2\mu }=0`$ From this follows that $$\begin{array}{c}\omega _{1\mu }=b_{(\mu ),1}\omega _{01}+c_{(\mu )}\omega _{02}\hfill \\ \omega _{2\mu }=c_{(\mu )}\omega _{01}+b_{(\mu ),2}\omega _{02},\hfill \end{array}$$ (10) where $`b_{(\mu ),\alpha }`$ and $`c_{(\mu )}`$ are some functions on the surface. The first fundamental form on the surface is $$I=<df_{(0)},df_{(0)}>=\omega _{01}^2+\omega _{02}^2$$ (11) To each normal (in $`M^{n+1}`$)direction $`\mu =3,4,\mathrm{}`$ there correspond a second fundamental form, $$II_\mu =\omega _{1\mu }\omega _{01}+\omega _{2\mu }\omega _{02}$$ (12) It’s convinient to introduce notations $$\omega _{ij}=\alpha _{ij}du+\beta _{ij}du$$ (13) We choose the conformal coordinates on the surfaces, such that $$I=e^{2\varphi }\left(du^2+dv^2\right)$$ (14) In this coordinates, $`<\frac{}{u}f_{(0)},\frac{}{v}f_{(0)}>=0`$, and therefore we can choose $`f_{(1)}`$ and $`f_{(2)}`$ in such a way that $$df_{(0)}(u,v)=\left(\omega _{01}f_{(1)}+\omega _{02}f_{(2)}\right)(u,v),$$ with $`\alpha _{02}=0,\beta _{01}=0`$ in notations ( 13). In conformal coordinates ( 14), we also have $`(\alpha _{01})^2=e^{2\varphi }`$ and $`(\beta _{02})^2=e^{2\varphi }`$, thus we can make a choise $$\omega _{01}=e^\varphi du;\omega _{02}=e^\varphi dv$$ (15) From ( 5), ( 8) we have $$\begin{array}{c}d\omega _{01}=\omega _{02}\omega _{12}\hfill \\ d\omega _{02}=\omega _{01}\omega _{12},\hfill \end{array}$$ and from ( 13), ( 15) it follows $$\begin{array}{cc}\hfill \alpha _{01}=e^\varphi ,& \beta _{01}=0\hfill \\ \hfill \alpha _{02}=0,& \beta _{02}=e^\varphi \hfill \\ \hfill \alpha _{12}=\frac{\varphi }{v},& \beta _{12}=\frac{\varphi }{u}\hfill \\ \hfill \alpha _{1\mu }=b_{(\mu ),1}e^\varphi ,& \beta _{1\mu }=c_{(\mu )}e^\varphi \hfill \\ \hfill \alpha _{2\mu }=c_{(\mu )}e^\varphi ,& \beta _{2\mu }=b_{(\mu ),2}e^\varphi \hfill \end{array}$$ (16) From ( 5), ( 8) $$d\omega _{12}=\omega _{01}\omega _{02}\underset{\mu =3,4,\mathrm{}}{}\omega _{1\mu }\omega _{2\mu },$$ and therefore, using ( 10), $$(\frac{^2}{u^2}+\frac{^2}{v^2})\varphi =e^{2\varphi }(1+\underset{\mu =3,4,\mathrm{}}{}det\left(\begin{array}{cc}b_{(\mu ),1}& c_{(\mu )}\\ c_{(\mu )}& b_{(\mu ),2}\end{array}\right).)$$ (17) $$\begin{array}{c}d\omega _{1\mu }=\omega _{12}\omega _{2\mu }+\underset{\nu =3,4,\mathrm{}}{}\omega _{1\nu }\omega _{\nu \mu }\hfill \\ d\omega _{2\mu }=\omega _{12}\omega _{1\mu }+\underset{\nu =3,4,\mathrm{}}{}\omega _{2\nu }\omega _{\nu \mu }\hfill \end{array}$$ (18) $$\begin{array}{c}2c_{(\mu )}\frac{\varphi }{u}\left(b_{(\mu ),1}b_{(\mu ),2}\right)\frac{\varphi }{v}+\frac{c_{(\mu )}}{u}\varphi \frac{b_{(\mu ),1}}{v}\varphi =\underset{\nu =3,4,\mathrm{}}{}\left(b_{(\nu ),1}\beta _{\nu \mu }c_{(\nu )}\alpha _{\nu \mu }\right)\hfill \\ \left(b_{(\mu ),2}b_{(\mu ),1}\right)\frac{\varphi }{u}2c_{(\mu )}\frac{\varphi }{v}+\frac{b_{(\mu ),2}}{u}\varphi \frac{c_{(\mu )}}{v}\varphi =\underset{\nu =3,4,\mathrm{}}{}\left(c_{(\nu )}\beta _{\nu \mu }b_{(\nu ),2}\alpha _{\nu \mu }\right)\hfill \\ \mu =3,4,\mathrm{}\hfill \end{array}$$ (19) $$d\omega _{\mu \nu }=\omega _{1\mu }\omega _{1\nu }\omega _{2\mu }\omega _{2\nu }\underset{\eta =3,4,\mathrm{}}{}\omega _{\mu \eta }\omega _{\eta \nu }$$ (20) $$\frac{\beta _{\mu \nu }}{u}\frac{\alpha _{\mu \nu }}{v}=e^{2\varphi }\left((b_{(\mu ),1}b_{(\mu ),2})c_{(\nu )}(b_{(\nu ),1}b_{(\nu ),2})c_{(\mu )}\right)+\underset{\eta =3,4,\mathrm{}}{}(\alpha _{\mu \eta }\beta _{\eta \nu }\beta _{\mu \eta }\alpha _{\eta \nu })$$ (21) ### 2.4 Lagrangians, Variational derivative, minimal Surface The equations of motion are just the equations for the minimal surface. They can be obtained from the condition that the variation of the Lagrangian $$L=\sqrt{Det_{\alpha \beta }(<\frac{}{u_\alpha }f_{(0)},\frac{}{u_\beta }f_{(0)}>)}𝑑u_1du_2;$$ (22) is zero, $`{\displaystyle \frac{\delta L}{\delta f_{(0)}}}=0`$, subject to $`<\delta f_{(0)},f_{(0)}>=0`$; the last ensures that we stay on the hyperboloid. In the conformal coordinates ( 14), those equations are $$\begin{array}{c}0=<\delta f_{(0)},\frac{}{u}\frac{f_{(0)}}{u}+\frac{}{v}\frac{f_{(0)}}{v}>=\hfill \\ <\delta f_{(0)},\frac{}{u}\left(e^\varphi f_{(1)}\right)+\frac{}{v}\left(e^\varphi f_{(2)}\right)>=<\delta f_{(0)},e^{2\varphi }(\mathrm{\#}f_{(0)}+\underset{\mu =3,4,\mathrm{}}{}(b_{(\mu )1}+b_{(\mu )2})f_{(\mu )})>;\hfill \end{array}$$ we made the computation in the conformal basis, ( 15), and used the Maurer Cartan equations, ( 16). Since $`<\delta f_{(0)},f_{(0)}>=0`$, and otherwise arbitrary, it follows that $$(b_{(\mu )1}+b_{(\mu )2})=0,\mu =3,4,\mathrm{}$$ (23) ### 2.5 Minimal surface in $`H_{(3)}`$ as an integrable system The Maurer-Cartan equations for the minimal surface ( 23) simplify, and for the surface in $`H_{(3)}`$ they are $$\begin{array}{c}\left(\frac{^2}{u^2}+\frac{^2}{v^2}\right)\varphi =e^{2\varphi }\left(1+b^2+c^2\right),\hfill \end{array}$$ (24) $$\begin{array}{c}2c\varphi _u2b\varphi _v+c_ub_v=0,\hfill \\ 2b\varphi _u+2c\varphi _v+b_u+c_v=0,\hfill \end{array}$$ (25) where $`bb_{(3)1}=b_{(3)2},cc_{(3)}`$; and $`\alpha _{12},\beta _{12}`$ are determined by $`\varphi `$, $`\alpha _{12}=\varphi _v,\beta _{12}=\varphi _u`$. The system ( 25) is integrable; it has a Lax pair, with a spectral parameter $`\lambda \text{}`$, for example this one (there is in fact a much better one for purposes of inverse scattering; but the one below is more geometric): $$\begin{array}{c}\frac{}{u}\mathrm{\Phi }=\left[\begin{array}{cccc}0& \frac{\lambda ^2+1}{2\lambda }e^\varphi & \frac{i\lambda ^2+i}{2\lambda }e^\varphi & 0\\ \frac{\lambda ^2+1}{2\lambda }e^\varphi & 0& \varphi _v& \frac{\lambda ^2(bic)+(b+ic)}{2\lambda }e^\varphi \\ \frac{i\lambda ^2+i}{2\lambda }e^\varphi & \varphi _v& 0& \frac{\lambda ^2(c+ib)+(cib)}{2\lambda }e^\varphi \\ 0& \frac{\lambda ^2(bic)+(b+ic)}{2\lambda }e^\varphi & \frac{\lambda ^2(c+ib)+(cib)}{2\lambda }e^\varphi & 0\end{array}\right]\mathrm{\Phi }\hfill \\ \frac{}{v}\mathrm{\Phi }=\left[\begin{array}{cccc}0& \frac{i\lambda ^2i}{2\lambda }e^\varphi & \frac{\lambda ^2+1}{2\lambda }e^\varphi & 0\\ \frac{i\lambda ^2i}{2\lambda }e^\varphi & 0& \varphi _u& \frac{\lambda ^2(c+ib)+(cib)}{2\lambda }e^\varphi \\ \frac{\lambda ^2+1}{2\lambda }e^\varphi & \varphi _u& 0& \frac{\lambda ^2(b+ic)+(bic)}{2\lambda }e^\varphi \\ 0& \frac{\lambda ^2(c+ib)+(cib)}{2\lambda }e^\varphi & \frac{\lambda ^2(b+ic)+(bic)}{2\lambda }e^\varphi & 0\end{array}\right]\mathrm{\Phi }\hfill \end{array}$$ (26) In fact, the integrable system here is something quite familiar. It follows from (25), that $$\begin{array}{c}\varphi _u=\frac{bb_u+cc_u+bc_vcb_v}{2(b^2+c^2)}\hfill \\ \varphi _v=\frac{bb_v+cc_vbc_u+cb_u}{2(b^2+c^2)};\hfill \end{array}$$ (27) here subscripts $`u`$ and $`v`$ denote derivatives with respect to $`u,v`$. Let’s introduce $`\rho `$ and $`\mathrm{\Theta }`$, such that $`b=\rho \mathrm{cos}\mathrm{\Theta }`$, $`c=\rho \mathrm{sin}\mathrm{\Theta }`$. Then it follows from ( 27 ) that $$\begin{array}{c}\left(\varphi +\frac{1}{2}\mathrm{log}\rho \right)_u=\mathrm{\Theta }_v,\hfill \\ \left(\varphi +\frac{1}{2}\mathrm{log}\rho \right)_v=\mathrm{\Theta }_u.\hfill \end{array}$$ Therefore, $`\mathrm{\Theta }`$ must be a harmonic function of $`(u,v)`$, ( as well as $`\left(\varphi +\frac{1}{2}\mathrm{log}\rho \right)`$), so whenever the only harmonic functions are constants; say if a surface is an imbedding of a sphere; then $`\left(\varphi +\frac{1}{2}\mathrm{log}\rho \right)`$ is some constant $`\kappa `$ as well; and so the equation ( 24) is in fact a $`cosh`$ -Gordon, $$\left(\frac{^2}{u^2}+\frac{^2}{v^2}\right)\varphi =e^{2\varphi }+\kappa e^{2\varphi }.$$ (28) ### 2.6 A remark on minimal surfaces in $`\text{}^3`$ and the Liouville equation. It is well known that a minimal surface in $`\text{}^3`$ is a surface with the mean curvature equal to zero. In the Maurer Cartan format, the equations for a surface in $`\text{}^3`$ are $$\begin{array}{c}dx=\omega _if_{(i)}\hfill \\ df_{(i)}=\omega _{ij}f_{(j)}\hfill \\ d\omega _i=\omega _j\omega _{ji}\hfill \\ d\omega _{ij}=\omega _{ik}\omega _{kj},\hfill \end{array}$$ (29) where $`x\text{}^3`$, and $`\{f_{(\mu )}\}`$ is an ortonormal frame; the group is the group of Eucledean motions of $`\text{}^3`$, instead of Lorents group which we work with. For a minimal surface, choosing the conformal coordinates ( 14), writing the Maurer-Cartan equations, and taking into account that the mean curvature is zero, similar to what we did in constant negative curvature space above, we would arrive at the Liouville equation, $$\left(\frac{^2}{u^2}+\frac{^2}{v^2}\right)\varphi =e^{2\varphi },$$ for which a solution can be written explicitly, as it is well known for a very long time; but the corresponding quantum field theory is regarded to be notoriously difficult . I find quite amusing the following set of facts: a) Maurer Cartan plays a major role in the geometry of frames on surfaces, and in particular it is responsible for the Liouville equation; b) some Maurer Cartan shows up in the celebrated deformation quantization construction of associative algebras, c) the way they approach quantum Liouville in is via associativity of the operator product algebra, and d) they seem to be using the same software to draw their pictures in their texts in b and c, and if you look at those pictures from far away, they look alike; but I do not know what exactly to make of those observations. For a surface of constant mean curvature h, we would obtain the $`sinh`$ Gordon, $`(\frac{^2}{u^2}+\frac{^2}{v^2})\varphi =h^2e^{2\varphi }+e^{2\varphi })`$. This and other integrable surfaces in $`\text{}^3`$ were studied in . ### 2.7 Minimal surface in $`H_{(5)}`$ as an integrable system The Maurer-Cartan equations for the minimal surface ( 23) in $`H_{(5)}`$ are $$\begin{array}{c}\left(\frac{^2}{u^2}+\frac{^2}{v^2}\right)\varphi =e^{2\varphi }\left(1+b_{3}^{}{}_{}{}^{2}+c_{3}^{}{}_{}{}^{2}+b_{4}^{}{}_{}{}^{2}+c_{4}^{}{}_{}{}^{2}+b_{5}^{}{}_{}{}^{2}+c_{5}^{}{}_{}{}^{2}\right),\hfill \\ L_{1\mu }\stackrel{(def.)}{=}2c_{(\mu )}\frac{\varphi }{u}2b_{(\mu )}\frac{\varphi }{v}+\frac{c_{(\mu )}}{u}\varphi \frac{b_{(\mu )}}{v}\varphi =\underset{\nu =3,4,5,\nu \mu }{}\left(b_{(\nu )}\beta _{\nu \mu }c_{(\nu )}\alpha _{\nu \mu }\right)\hfill \\ L_{2\mu }\stackrel{(def.)}{=}2b_{(\mu )}\frac{\varphi }{u}2c_{(\mu )}\frac{\varphi }{v}\frac{b_{(\mu )}}{u}\varphi \frac{c_{(\mu )}}{v}\varphi =\underset{\nu =3,4,5\nu \mu }{}\left(c_{(\nu )}\beta _{\nu \mu }+b_{(\nu )}\alpha _{\nu \mu }\right),\mu =3,4,5;\hfill \end{array}$$ (30) and $$\begin{array}{c}\frac{\beta _{\mu \nu }}{u}\frac{\alpha _{\mu \nu }}{v}=2e^{2\varphi }\left(b_{(\mu )}c_{(\nu )}b_{(\nu )}c_{(\mu )}\right)+\underset{\eta =3,4,5\eta \mu ,\nu }{}(\alpha _{\mu \eta }\beta _{\eta \nu }\beta _{\mu \eta }\alpha _{\eta \nu })\hfill \\ \mu ,\nu =3,4,5.\hfill \end{array}$$ (31) This system of equations appear integrable, and posess a Lax pair with spectral parameter, as follows. We assume that say $`\alpha _{45},\beta _{45}`$ can be represented in the form $$\alpha _{45}=\psi _u+\chi _v,\beta _{45}=\psi _v\chi _u$$ with certain functions $`\psi (u,v),\chi (u,v)`$; which doesnot seem to be terrribly restrictive. There is a Lax pair, reproducing the Maurer Cartan equations; it involves a spectral parameter $`\lambda \text{fontr C}`$, and the unknowns: $`\psi `$, $`\chi `$, the conformal factor $`\varphi (u,v)`$, as well as $`\{c_m(u,v),b_m(u,v)\|m=3,4,5\}`$, see ( 10), ( 12), (where $`b_m(u,v)\stackrel{(def)}{=}b_{(m),1}(u,v)=b_{(m),2}(u,v)`$, as the surface is minimal); that’s all we need to know to be able to find the Maurer-Cartan 1-forms, ( 5), and then the surface itself is obtained by solving linear compatible first order equations ( 4). Possibly, there are better, for purposes of boundary inverse scattering, Lax pairs; this is under investigation; but at least, there is some Lax pair: $$\begin{array}{c}\frac{}{u}\mathrm{\Phi }=\hfill \\ \left[\begin{array}{ccccc}0& \frac{\lambda ^2+1}{2\lambda }e^\varphi & \frac{i\lambda ^2+i}{2\lambda }e^\varphi & & (\mathrm{0\; 0\; 0})\\ \frac{\lambda ^2+1}{2\lambda }e^\varphi & 0& \varphi _v& & \left(\frac{\lambda ^2(b_\mu ic_\mu )+(b_\mu +ic_\mu )}{2\lambda }e^\varphi \|_{\mu =3,4,5}\right)\\ \frac{i\lambda ^2+i}{2\lambda }e^\varphi & \varphi _v& 0& & \left(\frac{\lambda ^2(c_\mu +ib_\mu )+(c_\mu ib_\mu )}{2\lambda }e^\varphi \|_{\mu =3,4,5}\right)\\ 0& \frac{\lambda ^2(b_3ic_3)+(b_3+ic_3)}{2\lambda }e^\varphi & \frac{\lambda ^2(c_3+ib_3)+(c_3ib_3)}{2\lambda }e^\varphi & & \\ 0& \frac{\lambda ^2(b_4ic_4)+(b_4+ic_4)}{2\lambda }e^\varphi & \frac{\lambda ^2(c_4+ib_4)+(c_4ib_4)}{2\lambda }e^\varphi & & \text{}\\ 0& \frac{\lambda ^2(b_5ic_5)+(b_5+ic_5)}{2\lambda }e^\varphi & \frac{\lambda ^2(c_5+ib_5)+(c_5ib_5)}{2\lambda }e^\varphi & & \end{array}\right]\mathrm{\Phi }\hfill \\ \text{}=\left[\begin{array}{ccc}0& a[3,4]& a[3,5]\\ a[3,4]& 0& (\psi _u+\chi _v)\\ a[3,5]& (\psi _u+\chi _v)& 0\end{array}\right]\hfill \\ a[3,4]=\frac{1}{(b_{3}^{}{}_{}{}^{2}+c_{3}^{}{}_{}{}^{2})}(c_3(L_{14}c_5(\psi _u+\chi _v)+b_5(\psi _v\chi _u))+\hfill \\ b_3(L_{24}+b_5(\psi _u+\chi _v)+c_5(\psi _v\chi _u)))\hfill \\ a[3,5]=\frac{1}{(b_{5}^{}{}_{}{}^{2}+c_{5}^{}{}_{}{}^{2})}\left(c_5L_{13}b_5L_{23}\right)+\hfill \\ \frac{1}{(b_{3}^{}{}_{}{}^{2}+c_{3}^{}{}_{}{}^{2})(b_{5}^{}{}_{}{}^{2}+c_{5}^{}{}_{}{}^{2})}((c_3(b_4b_5+c_4c_5)+b_3(b_4c_5b_5c_4))(L_{14}c_5(\psi _u+\chi _v)+b_5(\psi _v\chi _u))\hfill \\ +(c_3(b_4c_5b_5c_4)b_3(c_4c_5+b_4b_5))(L_{24}+b_5(\psi _u+\chi _v)+c_5(\psi _v\chi _u)))\hfill \end{array}$$ (32) $$\begin{array}{c}\frac{}{v}\mathrm{\Phi }=\hfill \\ \left[\begin{array}{ccccc}0& \frac{i\lambda ^2i}{2\lambda }e^\varphi & \frac{\lambda ^2+1}{2\lambda }e^\varphi & & (\mathrm{0\; 0\; 0})\\ \frac{i\lambda ^2i}{2\lambda }e^\varphi & 0& \varphi _u& & \left(\frac{\lambda ^2(c_\mu +ib_\mu )+(c_\mu ib_\mu )}{2\lambda }e^\varphi \right)\\ \frac{\lambda ^2+1}{2\lambda }e^\varphi & \varphi _u& 0& & \left(\frac{\lambda ^2(b_\mu +ic_\mu )+(b_\mu ic_\mu )}{2\lambda }e^\varphi \right)\\ 0& \frac{\lambda ^2(c_3+ib_3)+(c_3ib_3)}{2\lambda }e^\varphi & \frac{\lambda ^2(b_3+ic_3)+(b_3ic_3)}{2\lambda }e^\varphi & & \\ 0& \frac{\lambda ^2(c_4+ib_4)+(c_4ib_4)}{2\lambda }e^\varphi & \frac{\lambda ^2(b_4+ic_4)+(b_4ic_4)}{2\lambda }e^\varphi & & \text{}\\ 0& \frac{\lambda ^2(c_5+ib_5)+(c_5ib_5)}{2\lambda }e^\varphi & \frac{\lambda ^2(b_5+ic_5)+(b_5ic_5)}{2\lambda }e^\varphi & & \end{array}\right]\mathrm{\Phi }\hfill \\ \text{}=\left[\begin{array}{ccc}0& b[3,4]& b[3,5]\\ b[3,4]& 0& (\psi _v\chi _u)\\ b[3,5]& (\psi _v\chi _u)& 0\end{array}\right]\hfill \\ b[3,4]=\frac{1}{(b_{3}^{}{}_{}{}^{2}+c_{3}^{}{}_{}{}^{2})}(b_3(L_{14}c_5(\psi _u+\chi _v)+b_5(\psi _v\chi _u))+\hfill \\ c_3(L_{24}+b_5(\psi _u+\chi _v)+c_5(\psi _v\chi _u)))\hfill \\ b[3,5]=\frac{1}{(b_{5}^{}{}_{}{}^{2}+c_{5}^{}{}_{}{}^{2})}\left(b_5L_{13}c_5L_{23}\right)+\hfill \\ \frac{1}{(b_{3}^{}{}_{}{}^{2}+c_{3}^{}{}_{}{}^{2})(b_{5}^{}{}_{}{}^{2}+c_{5}^{}{}_{}{}^{2})}((b_5(c_3c_4b_3b_4)+c_5(b_4c_3b_3c_4))(L_{14}c_5(\psi _u+\chi _v)+b_5(\psi _v\chi _u))\hfill \\ +(b_5(b_3c_4b_4c_3)c_5(c_3c_4+b_3b_4))(L_{24}+b_5(\psi _u+\chi _v)+c_5(\psi _v\chi _u)))\hfill \end{array}$$ (33) where $`L_{ij}`$ are defined in ( 30). ## 3 Some thoughts on the loop equation, in the context of zero mean curvature surfaces. We currently do not know how to pose an inverse scattering problem for the equations we got; however, an experience with an inverse scattering on an arbitrary domain for integrable equations which have a linear limit, suggests that there exist a $`\overline{}`$ problem of a shape $$\frac{}{\overline{\lambda }}\mathrm{\Phi }(u,v,\lambda )=S_\gamma (u,v,\lambda )\mathrm{\Phi }(u,v,\lambda ),$$ (34) where $`S_\gamma (u,v,\lambda )`$ is determined from (a yet to be formulated) boundary condition. We do not know how exactly to get this $`\overline{}`$ problem here, but since all examples known so far come in this shape, we conjecture it exist here as well. Our Lagrangian is $$L=_\text{D}\sqrt{Det_{\alpha \beta }(<\frac{}{u_\alpha }f_{(0)},\frac{}{u_\beta }f_{(0)}>)}𝑑u_1du_2=_\text{D}\omega _{01}\omega _{02}$$ (35) with our choice of frame. We assume we have a family of boundary conditions, for which we can resolve the ( 34) problem, and therefore we have a family of solutions of ( 34) depending from the boundary condition $`\mathrm{\Phi }_\gamma (u,v,\lambda )`$, which give rise to a family of one forms $`\omega _{ij}=\left((d\mathrm{\Phi }_\gamma (u,v,1))\mathrm{\Phi }_{\gamma }^{}{}_{}{}^{1}(u,v,1)\right)_{ij}`$, depending from the boundary condition. We would like to compute in the second variation of the lagrangian with respect to change of boundary conditions, $`\delta _1\delta _2L\delta _{f_o(t_1)}\delta _{f_o(t_2)}`$ the term containing a delta function $`\delta (t_1t_2)`$. We will do it formally, assuming that there exist variation $`\delta `$ commuting with the differential. Since $$\begin{array}{c}\delta \omega _{01}=\delta ((d\mathrm{\Phi })\mathrm{\Phi }^1)_{01}=((d\delta \mathrm{\Phi })\mathrm{\Phi }^1)_{01}((d\mathrm{\Phi })\mathrm{\Phi }^1\delta \mathrm{\Phi }\mathrm{\Phi }^1)_{01}=\hfill \\ =((d\delta \mathrm{\Phi })\mathrm{\Phi }^1)_{01}\omega _{02}(\delta \mathrm{\Phi }\mathrm{\Phi }^1)_{21};\hfill \end{array}$$ here $`\mathrm{\Phi }\mathrm{\Phi }(u,v,1)SO(1,n)`$, and therefore the symmetry conditions are $`(\delta \mathrm{\Phi }\mathrm{\Phi }^1)_{oo}=0,(\delta \mathrm{\Phi }\mathrm{\Phi }^1)_{o\alpha }=(\delta \mathrm{\Phi }\mathrm{\Phi }^1)_{\alpha o},(\delta \mathrm{\Phi }\mathrm{\Phi }^1)_{\alpha \beta }=(\delta \mathrm{\Phi }\mathrm{\Phi }^1)_{\beta \alpha };`$ we used also our choice of the frame. We remark that $`\delta \mathrm{\Phi }\mathrm{\Phi }^1`$ are zero forms on the tangent space, and $`d\mathrm{\Phi }\mathrm{\Phi }^1`$ are one forms. Then $$\delta (\omega _{01}\omega _{02})=((d\delta \mathrm{\Phi })\mathrm{\Phi }^1)_{01}\omega _{02}+\omega _{01}((d\delta \mathrm{\Phi })\mathrm{\Phi }^1)_{02}.$$ The only terms in the second variation which would contain a $`\delta `$ function would come only from terms with the second derivative; as products of the first order derivatives cannot produce a delta function. Therefore (terms with a delta function possible in the second variation ) are $$((d\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1)_{01}\omega _{02}+\omega _{01}((d\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1)_{02}$$ Here $`\mathrm{\Delta }=\delta _1\delta _2`$ Proposition For a mean curvature zero surface $`b_{(\mu ),1}+b_{(\mu ),2}=0,\mu =3,4,5,`$ it follows from the Maurer Cartan equations, that the terms in the second variation which contain a $`\delta `$ function depend only from the boundary condition, and given by $$_\text{D}((\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1))_{01}\omega _{02}((\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1)_{02}\omega _{01},$$ (36) since $$\begin{array}{c}((d\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1)_{01}\omega _{02}+\omega _{01}((d\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1)_{02}=\hfill \\ =d\left(((\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1)\right)_{01}\omega _{02}\left((\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1\right)_{02}\omega _{01}).\hfill \end{array}$$ (37) Proof: The difference between the right hand side and the left hand side in ( 37) is $$\begin{array}{c}((\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1)_{o\mu }(\omega _{1\mu }\omega _{02}+\omega _{o1}\omega _{2\mu })+\hfill \\ ((\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1)_{o1}(d\omega _{02}\omega _{o1}\omega _{02})\hfill \\ ((\mathrm{\Delta }\mathrm{\Phi })\mathrm{\Phi }^1)_{o1}(d\omega _{01}\omega _{o2}\omega _{21});\hfill \end{array}$$ the first term is zero since $`(\omega _{1\mu }\omega _{02}+\omega _{o1}\omega _{2\mu })=(b_{(\mu ),1}+(b_{(\mu ),1})\omega _{01}\omega _{02}`$, and we have a $`(b_{(\mu ),1}+(b_{(\mu ),1})=0`$ surface; the other terms are zero due to the Maurer Cartan, as it looks in our choice of the frame. ## 4 Acknowledgments Hospitality and financial support of IHES and BRIMS institutes is appreciated. I was introduced to some of the ideas and methods used here during a conversation with Prof. I.M. Gelfand and Dr. Juan Carlos Alvarez Paiva, which occured at IHES in summer, 1998. I am grateful to Prof. Gelfand and Prof. Fokas for discussions. I am also grateful to X. Y. for sticking my nose into just a month after it appeared and then leaving me alone.
warning/0004/cond-mat0004031.html	ar5iv	text	# Nucleation on top of islands in epitaxial growth ## I Introduction A fundamental problem in the theory of thin film growth is the question, under which condition flat, two-dimensional films form on the substrate surface in contrast to mutually-separated, three-dimensional clusters. For films growing under equilibrium conditions, this question has been answered many years ago: If the interfacial tension between the substrate and adsorbate is larger than the difference of the respective surface free energies, then cluster formation is preferred (“Volmer-Weber growth” ), while a smaller (or equal) interfacial tension leads to the formation of flat films (“Van-der Merwe growth”). An intermediate case is the “Stranski-Krastanov growth” mode, where cluster formation sets in after the thickness of an initially smooth film exceeds a critical height. This case may be understood from an interfacial tension that varies with the film thickness. More recently, the influence of strain effects on equilibrium film morphologies has been investigated by various authors. Films developing in the process of molecular beam epitaxy (MBE) are usually not in thermodynamic equilibrium.. Rapid growth of films is achieved by a high supersaturation of the vapor, and growth kinetics is governed by evaporation, diffusion and aggregation processes far from equilibrium. Determining the film morphology in this situation is a problem of stochastic dynamics. (For a recent review on both kinetically and thermodynamically induced instabilities in MBE see ref. .) During the MBE experiment, two-dimensional islands composed of adsorbate atoms form on the substrate. If these islands coalesce before stable clusters nucleate on top of the islands in the second layer, a flat two-dimensional film results. By contrast, if the onset of second layer nucleation precedes island coalescence, then three-dimensional cluster formation is obtained. The term “second layer nucleation” should not be taken literally here but rather should apply to the formation of stable nuclei composed of $`(i+1)`$ atoms on top of islands in general. As for the equilibrium structures, it might be possible that cluster formation sets in above a certain film thickness, when the relevant parameters governing the nucleation of stable clusters on top of islands (see below) depend sensitively on the film thickness. However, despite this similarity of the possible growth processes with the equilibrium growth modes, it should be noted that the dynamic problem is very different. In MBE flat films can be produced even if the adsorbate does not wet the substrate. The first theory for second layer nucleation in MBE was set up by J. Tersoff, A. W. Denier van der Gon, and R. M. Tromp, which will be referred to as “TDT approach” in the following. Solving the stationary diffusion equation in the presence of an incoming flux and employing classical nucleation theory, these authors succeeded in deriving an explicit expression for the rate of nucleation $`\mathrm{\Omega }(R)`$ on top of circular shaped islands of radius $`R`$. Assuming all island radii to evolve approximately as the mean island radius $`R(t)t^{1/2}`$ at time $`t`$ (this situation will be referred to as the “single-island model” in the following), they calculated the fraction $`f_0(t)`$ of “covered islands” (i.e. on top of which a stable cluster has nucleated) from $`\mathrm{\Omega }(R)`$. It turned out that $`f_0(t)`$ rises from zero to one in the vicinity of a “critical time” $`t_c`$, which allows one to define a critical island radius $`R_cR(t_c)`$ for second layer nucleation. A simple criterion for the occurrence of “rough multilayer” as opposed to smooth “layer-by-layer growth” is that $`R_c`$ is smaller than the mean distance $`l`$ between islands in the first layer. An important factor controlling the film morphology is the additional step edge barrier $`\mathrm{\Delta }E_\mathrm{S}=E_\mathrm{S}E_0`$ (Ehrlich-Schwoebel barrier) that has to be surmounted by an adatom in addition to the bare surface diffusion barrier $`E_0`$, when an adatom crosses an island edge. For larger $`\mathrm{\Delta }E_\mathrm{S}`$ one expects adatoms to remain longer on islands and therefore to accumulate more easily, which would lead to an increased second layer nucleation rate $`\mathrm{\Omega }(R)`$ and a smaller $`R_c`$. In fact, the theory predicts that only for sufficiently large $`\mathrm{\Delta }E_\mathrm{S}`$ three-dimensional clusters can occur on the substrate (for an alternative possibility see however ref. ). Using the TDT approach, $`\mathrm{\Delta }E_\mathrm{S}`$ was estimated for a variety of different systems. An alternative approach for treating the problem of second layer nucleation within a stochastic description based on scaling arguments was developed recently by us. It was shown that for $`i=1`$ the TDT approach is not applicable, but the detailed treatment of fluctuations with only two atoms on top of the island yields a correct description of the process (see also ref.; for an earlier approach focusing on one dimension see ref. ). In this work we will extend our former study of second layer nucleation by means of both kinetic Monte Carlo simulations and scaling analysis. In particular, we will show that the mean-field assumptions underlying the TDT approach are valid for large critical nuclei $`i3`$, while for small critical nuclei $`i=1,2`$ second layer nucleation is dominated by fluctuations. Furthermore, we develop a novel rate equation approach, which allows one to calculate the time-development of cluster configurations on compact two-dimensional islands under quite general conditions. From the outset, one should distinguish between nucleation on top of islands with compact shape as opposed to nucleation on islands with strongly ramified shape. In the latter situation, diffusion of adatoms on the islands becomes a rather complex phenomenon due to the confined motion along branches of various lengths. We restrict our discussion to nucleation on compact islands here. Intuitively, one would expect second layer nucleation on ramified islands to be unlikely, so that this restriction is of minor importance. Moreover, it should be noted that even for compact shapes, the island boundaries may have a fractal or, more precisely, self-affine structure (this is the case e.g. for Eden clusters ). Then the microscopic step edge barrier can vary strongly along the island boundary. In any case we will always understand $`\mathrm{\Delta }E_\mathrm{S}`$ as an effective barrier (see below and ref. ). Parameters governing the second layer nucleation are the incoming atom flux $`F`$, the jump rate $`D/a^2`$ of adatoms, the step edge barrier $`\mathrm{\Delta }E_\mathrm{S}`$, and various dissociation rates of unstable clusters of size $`si`$. If the bond energies of the unstable clusters (i.e. of clusters of size $`si`$) are negligibly small, then the nucleation rate $`\mathrm{\Omega }(R)`$ and critical radius $`R_c`$ depend only on two dimensionless parameters. These are the ratio $$\mathrm{\Gamma }\frac{D}{Fa^4}$$ (1) (with $`a`$ being the lattice spacing in the substrate plane) and the edge crossing probability $$\alpha \mathrm{exp}\left(\frac{\mathrm{\Delta }E_\mathrm{S}}{k_\mathrm{B}T}\right).$$ (2) Based on a simplified stochastic description (see also ref. ) we will argue that for small critical nuclei $`i=1,2`$ the mean number of atoms on top of the island is smaller than one and the stable nucleus is formed due to fluctuations. This gives rise to four scaling regimes in an $`\alpha \mathrm{\Gamma }`$ diagram, where $`R_c\mathrm{\Gamma }^\gamma \alpha ^\mu `$ with different exponents $`\gamma `$ and $`\mu `$. For $`i3`$ by contrast, nucleation starts out from a situation with many atoms present on the island. Under these circumstances, three different scaling regimes can be identified, and two of them correspond to the ones predicted by the TDT approach. By comparing $`R_c`$ with the mean distance $`l`$ between islands on the substrate surface, the transition line separating rough multilayer from smooth layer-by-layer growth is identified in the $`\alpha \mathrm{\Gamma }`$ diagram. When the bond energies of unstable clusters become appreciable, the corresponding dissociation rates enter the problem as additional relevant parameters. It becomes difficult then to separate scaling regimes in practice, and the simplified stochastic description becomes of limited value. However, by employing the novel rate equation approach it is still possible to determine $`f_0(t)`$ and $`R_c`$ in a simple manner. Moreover, we will discuss how to derive the fraction of $`f(t)`$ of covered islands, when one relaxes the assumption that all island radii evolve as the mean radius $`R(t)`$. In the time regime of almost constant island density (“saturation regime” preceeding island coalescence) we can define an effective “capture area” for adatoms by the Voronoi cell for each island. In order to calculate $`f(t)`$ for the “multi-island model” from $`f_0(t)`$, one needs to know the probability distribution of islands with a certain size and capture area, when the saturation regime is reached. The paper is organized as follows. In Sec. II we give a short review on basic concepts used in the description of submonolayer growth and discuss in part II A important quantities and equations underlying the physical processes involved. In part II B we summarise the results of the TDT approach. We then continue with a detailed description of the simulation techniques in Sec. III and test the equivalence of the multi-island and single-island model. In Sec. IV a simplified stochastic description of second layer nucleation is first presented in its general methodology and subsequently employed to small and large critical nuclei. It is instructive to observe the different physical conditions that lead to the nucleation event in the two cases. After showing that the treatment of metastable clusters is rather complicated within the simplified description, we discuss in Sec. V a general approach for second layer nucleation on the basis of novel rate equations. Sec. VI concludes the paper with a summary and discussion of the most important results as well as an outlook to further research. ## II Basic quantities and concepts ### A Atomistic processes in thin film growth In MBE atoms are deposited on a substrate surface with a rate $`Fa^2`$ per unit cell. At very high temperatures the adatoms reevaporate but under ordinary conditions this reevaporation can be neglected or effectively taken into account by a reduced deposition rate. Once an adatom is deposited, it starts a thermally activated diffusive motion with jump rate $`D/a^2\mathrm{exp}(E_0/k_\mathrm{B}T)`$. Adatoms come into contact as time progresses, and form islands that are held together by some bonding energy. While unstable islands of small size $`si`$ dissociate, islands with size $`s>i`$ are stable (on all relevant time scales of the experiment). An island of size $`s=i`$ is called a critical nucleus. Islands of larger size are formed by aggregation of adatoms (or small mobile islands) to existing immobile islands, and by coalescence. These processes lead to an island size distribution which becomes broader with increasing time. In the submonolayer regime, the most important physical questions are: (i) How large is the density $`\rho _x(t)`$ of stable islands on the substrate surface at time $`t`$? (ii) What is the form of the distribution $`\chi _s(s,t)`$ of island sizes $`s`$ at time $`t`$ (with $`s`$ being the number of atoms forming the island)? (iii) What do the stable islands look like? These questions have been extensively studied in the past, both by experiment and by theory. We will briefly summarise those results, which are relevant for the following analysis. The typical behavior of $`\rho _x(t)`$ is depicted in Fig. 1. Also shown is the adatom density $`\rho _1(t)`$. As suggested by Amar and Family, one may distinguish between four different time regimes: The low-coverage regime $`L`$ where $`\rho _x(t)`$ increases with $`t`$ and $`\rho _x(t)<\rho _1(t)`$, the intermediate coverage regime I, where $`\rho _x(t)`$ increases with $`t`$ and $`\rho _x(t)>\rho _1(t)`$, the saturation regime S (called aggregation regime A in ref. ), where $`\rho _x(t)`$ stays approximately constant, and the coalescence regime C, where $`\rho _x(t)`$ strongly decreases due to coalescence of stable clusters. At high coverages $`\theta Fa^2t`$, the monomer density $`\rho _1(t)`$ becomes small (see Fig. 1), and the standard rate equations for submonolayer growth predict $`\rho _x`$ to evolve as $`\rho _x(t)`$ $``$ $`\left({\displaystyle \frac{D}{Fa^4}}\right)^{\frac{i}{i+2}}(Fa^2t)^{\frac{1}{i+2}}e^{E_i/(i+2)k_\mathrm{B}T}`$ (4) $`=\mathrm{\Gamma }^{\frac{i}{i+2}}\theta ^{\frac{1}{i+2}}e^{E_i/(i+2)k_\mathrm{B}T},`$ where $`E_i`$ is the bonding energy of the critical nucleus in its preferred atomic configuration. It should be noted that eq. (4) is not valid in the saturation regime S, where the densities of islands with subcritical and critical size are very small (unless there are metastable subcritical nuclei). In this regime almost all adatoms being deposited attach to preexisting stable islands, so that $`\rho _x(t)`$ stays constant, $`\rho _x(t)=\rho _x`$. Within the standard rate equation approach, this effect may be accounted for by a proper dependence of the “capture numbers” on the adatom density $`\rho _1(t)`$ (for a detailed discussion of this point in relation to experiments see ref. ). The scaling of $`\rho _x`$ with $`\mathrm{\Gamma }`$, however, is still correct in regime S, $$\rho _x\mathrm{\Gamma }^{\frac{i}{i+2}}.$$ (5) One of the most detailed studies of the island size distribution has been performed by Amar and Family based on the scaling ansatz $$\chi _s(s,t)=\frac{1}{s(t)}f\left(\frac{s}{s(t)}\right).$$ (6) Here $`\chi _s(s,t)=\rho _s(t)/\rho _{\mathrm{tot}}(t)`$ is the probability that an island has size $`s`$; ($`\rho _s(t)`$ is the density of islands with size $`s`$ and $`\rho _{\mathrm{tot}}(t)`$ is the total island density). $`\mathrm{}`$ denotes an average over $`s`$ with respect to $`\chi _s(s,t)`$. From (6) follows, for $`s^m<\mathrm{}`$, $`s^m=s^m_0^{\mathrm{}}𝑑xx^mf(x)`$, and taking $`m=0,1`$ one obtains $`𝑑xf(x)=𝑑xxf(x)=1`$. The mean island size $`s(t)`$ is given by $`s(t)=_{s=1}^{\mathrm{}}s\rho _s(t)/\rho _{\mathrm{tot}}(t)=\theta /\rho _{\mathrm{tot}}(t)a^2`$. In the saturation regime S, in particular, where $`\rho _{\mathrm{tot}}\rho _x`$, since the density of islands with subcritical and critical size is small, it follows from eq. (5) $$s(t)\frac{\theta }{\rho _xa^2}\theta \mathrm{\Gamma }^{\frac{i}{i+2}}.$$ (7) The relation $`s(t)theta/\rho _xa^2=Ft/\rho _x`$ can also be understood more directly, since the increase of $`s(t)`$ with $`t`$ is given by the flux times the mean capture area $`\rho _x^1`$ of adatoms. The scaling function $`f(u)`$ was suggested to have the form $$f(u)=C_iu^i\mathrm{exp}(ia_iu^{1/a_i})$$ (8) in regime S. This function has a maximum at $`u=1`$ and the two conditions $`𝑑xf(x)=𝑑xxf(x)=1`$ determine the parameters $`C_i`$ and $`a_i`$. Equations (6-8) have been shown to give a fairly good approximation of some simulations and experiments. The problem of the island shapes is not yet well understood, but one may roughly answer the third question posed above as follows. For strictly irreversible attachment, where local relaxation of atoms due to fast edge diffusion is suppressed, one obtains dendritic or random fractal structures. Dendritic growth is preferred at low $`T`$ or small $`F`$, and a shape transition from dendritic to random fractal structures has been found e.g. for Ag/Pt(111) upon lowering the deposition flux. An example for a random fractal structure obtained in a computer simulation is shown in Fig. 2a. At high temperatures, edge diffusion becomes relevant, and polygonal or “irregular” compact island morphologies develop (see e.g. ref. ). An example for an irregular structure is shown in Fig. 2b. As mentioned in the Introduction, for the study of second layer nucleation we will focus on compact island shapes. Moreover, second layer nucleation in the intermediate regime I is unlikely to occur, since the island radii in this regime are typically smaller than the critical radius $`R_c`$. We therefore consider the second layer nucleation in the saturation regime S, where eqs. (5-8) apply. ### B TDT Approach In the TDT approach, one starts by calculating the adatom density $`\rho _1^{\mathrm{st}}`$ on a circular island with radius $`R`$ in the stationary state. The stationary diffusion equation with the incoming atom flux acting as a source term reads $$D\left[\frac{^2}{r^2}+\frac{1}{r}\frac{}{r}\right]\rho _1^{\mathrm{st}}+F=0,$$ (9) and it is supplemented by the boundary conditions ($`\alpha =\mathrm{exp}(\mathrm{\Delta }E_\mathrm{S}/k_\mathrm{B}T)`$) $$\frac{\rho _1^{\mathrm{st}}}{r}\|_{r=0}=0,\frac{\rho _1^{\mathrm{st}}}{r}\|_{r=R}=\frac{\alpha }{a}\rho _1^{\mathrm{st}}\|_{r=R},$$ (10) where $`a/\alpha `$ is commonly referred to as the “Schwoebel length”. The boundary conditions express the fact that the current density $`D\rho _1^{\mathrm{st}}/r`$ must vanish at the origin and that at the edge it is given by the density $`\rho _1`$ times the “velocity” (rate times lattice spacing) $`(D\alpha /a^2)a`$ to cross the step edge barrier. The solution of eqs. (9,10) is $$\rho _1^{\mathrm{st}}(r)=\rho _1^{\mathrm{st}}(0)\frac{Fr^2}{4D},\rho _1^{\mathrm{st}}(0)=\frac{FR^2}{4D}\left(1+\frac{2a}{\alpha R}\right).$$ (11) According to standard rate equation theory the local nucleation rate is proportional to $`D\rho _1^{i+1}`$, so we obtain from eq. (11) $`\mathrm{\Omega }(R)`$ $`=`$ $`\kappa {\displaystyle \frac{D}{a^2}}{\displaystyle _0^R}{\displaystyle \frac{2\pi rdr}{a^2}}[\rho _1^{\mathrm{st}}(r)a^2]^{i+1}`$ (12) $`=`$ $`{\displaystyle \frac{4\pi \kappa \mathrm{\Gamma }^{(i+1)}}{(i+2)\alpha ^{2(i+2)}}}{\displaystyle \frac{D}{a^2}}\left({\displaystyle \frac{\alpha R}{2a}}\right)^{i+2}\left[\left(1+{\displaystyle \frac{\alpha R}{2a}}\right)^{i+2}1\right]`$ (13) $``$ $`\{\begin{array}{cc}4\pi \kappa \mathrm{\Gamma }^{(i+1)}\alpha ^{(i+1)}{\displaystyle \frac{D}{a^2}}\left({\displaystyle \frac{R}{2a}}\right)^{i+3},\hfill & \alpha {\displaystyle \frac{2a}{R}}\hfill \\ {\displaystyle \frac{4\pi \kappa \mathrm{\Gamma }^{(i+1)}}{i+2}}{\displaystyle \frac{D}{a^2}}\left({\displaystyle \frac{R}{2a}}\right)^{2(i+2)},\hfill & \alpha {\displaystyle \frac{2a}{R}}\hfill \end{array}`$ (16) where $`\kappa `$ is a constant. For a given time evolution of the island radius $`R=R(t)`$, one can calculate the probability $`f_0(t)`$ for a stable nucleus to have formed on top of the island up to time $`t`$ as follows: The increase $`f_0(t+\mathrm{\Delta }t)f_0(t)`$ in a small time interval $`\mathrm{\Delta }t`$ is equal to the probability $`[1f_0(t)]`$ that up to time $`t`$ no stable nucleus has formed times the probability $`\mathrm{\Omega }(R(t))\mathrm{\Delta }t`$ that the nucleation takes place in the time interval $`[t,t+\mathrm{\Delta }t]`$. Taking the limit $`\mathrm{\Delta }t0`$ and solving the corresponding differential equation with the initial condition $`f_0(0)=0`$ yields $$f_0(t)=1\mathrm{exp}\left[_0^t𝑑t^{}\mathrm{\Omega }(R(t^{}))\right].$$ (17) For compact island growth during an MBE experiment, we have $`R(t)s(t)^{1/2}`$ and thus from eq. (7) $$\frac{R(t)}{a}=A(Fa^2t)^{1/2}\mathrm{\Gamma }^{i/2(i+2)}$$ (18) with $`A`$ being some constant. Inserting this growth law into (16,17) yields $`f_0(t)`$ $`=`$ $`1\mathrm{exp}\left[{\displaystyle \frac{2\mathrm{\Gamma }^{\frac{i}{i+2}}}{A^2Fa^4}}{\displaystyle _0^{R(t)}}𝑑rr\mathrm{\Omega }(r)\right]`$ (22) $`\{\begin{array}{cc}1\mathrm{exp}\left[C_<\mathrm{\Gamma }^{\frac{i(i+3)}{i+2}}\alpha ^{(i+1)}\left({\displaystyle \frac{R}{a}}\right)^{i+5}\right],\hfill & \alpha {\displaystyle \frac{2a}{R(t)}}\hfill \\ 1\mathrm{exp}\left[C_>\mathrm{\Gamma }^{\frac{i(i+3)}{i+2}}\left({\displaystyle \frac{R}{a}}\right)^{2(i+3)}\right],\hfill & \alpha {\displaystyle \frac{2a}{R(t)}}\hfill \end{array}`$ where $`C_>(2^{2(i+1)}\pi \kappa A^2)/[(i+2)(i+3)]`$ and $`C_<2^i\pi \kappa A^2`$. In going from the first to the second line in (22) we have used that the integral over $`r`$ is dominated by the upper integration bound $`R(t)`$ (for $`R(t)/a1`$). It follows that the critical radius scales as $$R_c\mathrm{\Gamma }^\gamma \alpha ^\mu ,$$ (23) where $$\gamma =\{\begin{array}{cc}\frac{i(i+3)}{(i+2)(i+5)},\hfill & \alpha \mathrm{\Gamma }^{i/[2(i+2)]}\hfill \\ \frac{i}{2(i+2)},\hfill & \alpha \mathrm{\Gamma }^{i/[2(i+2)]}\hfill \end{array}$$ (24) and $$\mu =\{\begin{array}{cc}\frac{(i+1)}{(i+5)},\hfill & \alpha \mathrm{\Gamma }^{i/[2(i+2)]}\hfill \\ 0,\hfill & \alpha \mathrm{\Gamma }^{i/[2(i+2)]}\hfill \end{array}$$ (25) Equations (23-25) predict that for large step edge barriers, $`R_c`$ depends strongly on $`\mathrm{\Delta }E_\mathrm{S}`$, $`R_c\mathrm{exp}[(i+1)\mathrm{\Delta }E_\mathrm{S}/(i+5)k_\mathrm{B}T]`$, while for small barriers, $`R_c`$ becomes independent of $`\mathrm{\Delta }E_\mathrm{S}`$. For $`i=1`$ in particular, one finds $`R_c\mathrm{\Gamma }^{2/9}\alpha ^{1/3}`$ for $`\alpha \mathrm{\Gamma }^{1/6}`$ and $`R_c\mathrm{\Gamma }^{1/6}`$ for $`\alpha \mathrm{\Gamma }^{1/6}`$. ## III Kinetic Monte Carlo Simulations Kinetic Monte Carlo simulations are a well-established technique for modeling MBE experiments. In our investigation of second layer nucleation, we adopt a simulation scheme similar to previous, successful models of surface growth kinetics. We chose a substrate with fcc(111) symmetry, since surfaces of that kind are often studied in metal epitaxy, and commonly exhibit high $`\mathrm{\Delta }E_\mathrm{S}`$. In a full simulation scheme of the growth kinetics (multi-island model), we include all processes of evaporation, diffusion and aggregation occurring in the MBE experiment. By analyzing the set of islands of various size on the substrate, we determine the fraction $`f(t)`$ of covered islands at time $`t`$. On the other hand, we consider, as in the TDT approach, only one island with the mean radius $`R(t)`$ evolving deterministically in time. The fraction $`f_0(t)`$ of covered islands in this single-island model is then determined by calculating the probability for second layer nucleation up to time $`t`$ from a large set of independent simulations. By examining both models we are able to quantify the influence of the cluster size distribution under generic growth conditions. ### A Multi-island model Atoms are randomly deposited with a rate $`Fa^2`$ per unit cell onto a triangular lattice. After instantaneously relaxing to a position, where they are supported by three nearest neighbors in the layer below (“downward funneling”), the atoms change their position by performing thermally activated jumps to a vacant nearest neighbor site in the same layer with a rate $`D/6a^2`$. Only one atom is allowed to occupy a given lattice site. We first consider a “non-interacting particle model” where all binding energies of subcritical clusters of size $`si`$ are neglected. This means that $`(i+1)`$ adatoms have to encounter each other on nearest neighbor sites in order to form a stable nucleus. Within the single-island model we later will also consider finite binding energies of metastable clusters, which causes various dissociation rates to enter the problem as additional parameters. Once a stable cluster of size $`s>i`$ has formed, adatoms can attach to it. Compact island morphologies are known to emerge if a fast diffusion process is present along island edges. Here we model this process similar as in earlier approaches (see e.g. refs. ) by including a local relaxation mechanism. In this method an atom being in contact with at least one nearest neighbor after a jump, is immediately transferred to a nearest neighbor site, if it can increase its coordination number. This procedure is repeated until the atom can no longer increase its local coordination (see Fig. 3). Interlayer diffusion of atoms deposited onto islands is hindered by the Ehrlich-Schwoebel barrier $`\mathrm{\Delta }E_\mathrm{S}`$, which reduces the jump rate $`D/6a^2`$ by the edge crossing probability $`\alpha =\mathrm{exp}(\mathrm{\Delta }E_\mathrm{S}/k_\mathrm{B}T)`$. For computational convenience, we model the crossing by a two-step process in the simulation: First, when an atom passes the boundary, it remains in the same layer but moves to a place, where it is supported by only two atoms underneath. Then the atom immediately drops down to the layer below and moves to the nearest “stable” site according to the local relaxation mechanism introduced above (for similar simulations including $`\mathrm{\Delta }E_\mathrm{S}`$, see, e.g. ). We do not distinguish between crossing of A and B steps and have not attempted to model any more realistic scenarios, as e.g. collective rearrangements of atoms including exchange processes. This is well justified as long as one is interested in the influence of an effective Schwoebel barrier. In the following, we focus on the case $`i=1`$ first. Typical film morphologies resulting from the simulations have been shown in Fig. 2b. Note that the boundaries of the islands are still rough despite the local relaxation mechanism. The fraction $`f(t)`$ of covered islands as a function of the total coverage $`Fa^2t`$ is shown in Fig. 4 for some representative parameters (full symbols). As expected, $`f(t)`$ first is close to zero, then increases strongly in some time interval around a “critical time” $`t_c`$, and finally saturates at one. In the inset of Fig. 4 we show the dependence of the mean island radius $`R(t)(s(t)/\pi )^{1/2}a`$ on $`Fa^2t`$ during the evaporation. In agreement with eq. (18), we find $`R(t)=A(Fa^2t)^{1/2}\mathrm{\Gamma }^{1/6}`$ with $`A0.78`$. To be specific, let us define the critical time $`t_c`$ via the condition $`f(t_c)=1/2`$, and the corresponding critical island radius $`R_c`$ by $`R_c=R(t_c)`$, $$f(t_c)=1/2,R_cR(t_c)=A(Fa^2t_c)^{1/2}\mathrm{\Gamma }^{i/2(i+2)}a.$$ (26) Plots of $`R_c`$ as a function of $`\alpha `$ for various fixed $`\mathrm{\Gamma }`$ are shown in Fig. 5 (full symbols). With increasing step edge barrier, i.e. decreasing $`\alpha `$, adatoms on average remain longer on an island and nucleation of stable dimers occurs at smaller island radii. Accordingly, $`R_c`$ decreases with decreasing $`\alpha `$ (see “regime II” in the figure). For very small $`\alpha `$, however, the step edge barrier is practically never surmounted and thus is in effect infinitely high. Therefore, $`R_c`$ becomes independent of $`\alpha `$ (“regime I” in Fig. 5). The crossover between the two regimes is marked by the thick solid line. The full symbols in Fig. 5 terminate at the dashed line $`\alpha _{}(\mathrm{\Gamma })`$, which marks the onset of island coalescence. For $`\alpha >\alpha _{}(\mathrm{\Gamma })`$, islands in the first layer merge before second layer nucleation takes place and $`R_c`$ can no longer determined from the multi-island model. It is important to note that the dependence of $`R_c`$ on $`\alpha `$ is much weaker than predicted by the TDT approach: The solid lines in regime II have slope 1/7 corresponding to a power law $`R_c\alpha ^{1/7}`$ rather than $`R_c\alpha ^{1/3}`$ as predicted by eqs. (23, 25). Moreover, regime I does not occur in the TDT approach. ### B Single-island model Second layer nucleation can also be addressed in a simpler model, which does not attempt to describe the entire growth dynamics, but focuses on the decisive factors that determine nucleation in the presence of the step-edge barrier. In this model the complicated nucleation and diffusion-mediated growth of the two-dimensional islands, on which the second layer nucleation takes place, is replaced by letting the radius of one circular island expand deterministically in time as $`R(t)/a=A(Fa^2t)^{1/2}\mathrm{\Gamma }^{i/2(i+2)}`$, where $`A`$ is taken from the full simulation of the multi-island model. The island is embedded in a substrate area large enough to accommodate the island at all relevant times. Deposition and diffusion of adatoms take place in the same manner as in the multi-island model. Atoms inside the island boundary can escape by overcoming the step edge barrier. Those atoms that have surmounted the barrier or that have been deposited outside the island boundary are removed from the lattice. Thus the single-island model considers the deposition of random walkers within a time-dependent, circular boundary that is partially reflecting. Due to its greater simplicity, it allows for more specific analysis with a larger parameter space (there is no restriction due to coalescence of distinct islands). In the non-interacting particle model, the “critical event” is to find $`(i+1)`$ atoms on neighboring lattice sites. Analogous to the multi-island model we can define the fraction $`f_0(t)`$ of covered islands up to time $`t`$. The fraction now refers to a set of islands obtained in independent simulation runs. All islands in these runs grow with the same deterministic growth law. Results for $`f_0(t)`$ are shown in Fig. 5 (open symbols) for the same parameters as in the multi-island model. Good agreement with $`f(t)`$ is achieved for small times (corresponding to $`f(t)1/2`$), when the time in the single-island model is rescaled by a constant factor, i.e. $`f(t)f_0(t^{})`$ with $`t^{}=t/1.21`$. The factor is a consequence of the idealized circular island perimeter in the single-island model. In the multi-island model by contrast, the islands are far from being perfectly circular (see Fig. 2). They have rougher edges with more boundary sites, which causes adatoms to escape the islands more easily and second layer nucleation to occur at later times $`t1.21t^{}`$. At larger times (corresponding to $`f(t)1/2`$), however, $`f(t)`$ deviates from $`f_0(t^{})`$ and these deviations become more pronounced for larger $`\mathrm{\Gamma }`$. The reason for this discrepancy is the presence of islands with size much smaller than $`s(t)`$ in the multi-island model. Nucleation of stable clusters on top of these islands occurs at a later time, which causes $`f(t)`$ to be smaller than $`f_0(t^{})=f_0(t/1.21)`$ at large $`t`$. In fact, we will show in Sec. III C that this effect can be accounted for by considering the probability distribution of islands with a certain size and capture area. When $`R_c`$ approaches the mean distance $`l`$, coalescences of larger islands also lead to modifications of $`f(t)`$ for $`tt_c`$. The critical radius $`R_c^{}`$ in the single-island model can be defined as in the many island model by $`R_c^{}=R(t_c^{})`$, where $`f_0(t_c^{})=1/2`$. Due to the fact that $`t_c=1.21t_c^{}`$ we expect $`R_c=1.21^{1/2}R_c^{}=1.1R_c^{}`$. Results for $`1.1R_c^{}`$ as a function of $`\alpha `$ are shown in Fig. 5 (open symbols) for the same parameters as in the multi-island model (full symbols). As can be seen from the figure, there is almost perfect agreement between both data sets. Moreover, the data for $`R_c^{}`$ can be obtained also beyond the dashed line marking the onset of layer-by-layer growth. Let us also note that, as long as one is interested only in $`R_c^{}`$ (or $`R_c=1.1R_c^{}`$), one may obtain it even more simply in the single-island model (without calculating $`f_0(t)`$) by determining the average radius of the island at the time of the nucleation event, $$R_c^{}=R(t_c^{})_0^{\mathrm{}}𝑑t\frac{df_0(t)}{dt}R(t).$$ (27) Note that $`df_0(t)/dt`$ is the probability density of the second layer nucleation times and that the average of $`R(t)`$ with respect to $`df_0(t)/dt`$ is approximately equal to $`R(t_c^{})`$, since $`df_0(t)/dt`$ is sharply peaked around $`t_c^{}`$. ### C Equivalence of the single-island and the multi-island model In order to determine $`f(t)`$ from $`f_0(t)`$ we define by $`\psi (s,\sigma ,t)dsd\sigma `$ the probability for an island to have a size in the interval $`[s,s+ds]`$ and a capture area in the interval $`[\sigma ,\sigma +d\sigma ]`$ at time $`t`$, where the capture area is given by the Voronoi cell associated with an island. Let us consider $`f_0(t)`$ to be a functional of the growth law $`R(t)`$ only, as it is the case, for example, when one approximates the second layer nucleation by a Poisson process with a time dependent nucleation rate $`\mathrm{\Omega }(R(t))`$. Then $`f_0(t)=G_0[R(t)]=1\mathrm{exp}[_0^t𝑑t^{}\mathrm{\Omega }(R(t^{}))]`$ (see eq. 17). In the saturation regime the growth law for an island can be written as $`\pi R^2(t)=s_\times +F\sigma (tt_\times )`$, where $`t_\times `$ is the time when the saturation is reached (see Fig. 1) and $`s_\times `$ is the island size at that time. (We restrict ourselves to film-morphologies far from coalescence here, so that $`\sigma `$ can be regarded as time-independent.) With the specified growth law, the functional $`G_0[R(t)]`$ can be expressed by a function $`g_0=g_0(t;s_\times ,\sigma ,t_\times )`$, and $`f(t)`$ is calculated via $$f(t)=_0^{\mathrm{}}𝑑s_\times _0^{\mathrm{}}𝑑\sigma \psi (s_\times ,\sigma ,t_\times )g_0(t;s_\times ,\sigma ,t_\times ).$$ (28) A detailed investigation of the probability distribution $`\psi (s,\sigma ,t)`$ is certainly of interest but beyond the scope of the present work. A simple idea would be to neglect correlations between the stochastic variables $`s`$ and $`\sigma `$, $`\psi (s,\sigma ,t)\chi _s(s,t)\chi _\sigma (\sigma ,t)`$, and to use previously derived scaling forms for the island size distribution $`\chi _s(s,t)`$ (see e.g. refs. ) and the capture area distribution $`\chi _\sigma (\sigma ,t)`$ (see e.g. refs. ). Here we will follow a simpler approach. Since for typical situations we find both $`s_\times `$ and $`t_\times `$ to be significantly smaller than $`s_c=\pi R_c^2`$ and $`t_c`$, respectively, we use the growth law $`s(t)=F\sigma t`$ for an island with capture area $`\sigma `$ in the full simulation. The on-top nucleation probabilities $`\stackrel{~}{g}_0(t;\sigma )`$ for islands exhibiting different capture areas can then be related by a rescaling of time, i.e. $`\stackrel{~}{g}_0(t;\sigma _1)=\stackrel{~}{g}_0(\sigma _1t/\sigma _2;\sigma _2)`$. Moreover, since for film morphologies far from coalescence ($`\alpha \alpha _{}(\mathrm{\Gamma })`$), $`\chi (\sigma ,t)`$ is approximately independent of time, we have $`f_0(t)=\stackrel{~}{g}_0(t;\overline{\sigma })`$, where $`\overline{\sigma }=𝑑\sigma \chi _\sigma (\sigma ,t)\sigma \rho _x^1`$. Hence, $$f(t)=_0^{\mathrm{}}𝑑\sigma \chi _\sigma (\sigma )\stackrel{~}{g}_0(t;\sigma )=_0^{\mathrm{}}𝑑\sigma \chi _\sigma (\sigma )f_0\left(\frac{\overline{\sigma }}{\sigma }t\right).$$ (29) In this simplified eq. (29) knowledge of the nucleation rate $`\mathrm{\Omega }(R)`$ is not necessary and $`f(t)`$ can be directly obtained from $`f_0(t)`$ when $`\chi _\sigma (\sigma )`$ is known. Writing $`\chi _\sigma (\sigma )=\overline{\sigma }^1h(\sigma /\overline{\sigma })`$, where $`𝑑xh(x)=𝑑xh(x)x=1`$, the transformation (29) becomes $`f(t)=_0^{\mathrm{}}𝑑xh(x)f_0(xt)`$. For a random distribution of point islands, we would have $`h(x)=\mathrm{exp}(x)`$. However, since there is a depletion zone of adatoms near an island, the probability for other islands to nucleate in an area close to an existing one is reduced and not exponential. For an isolated island, dimensional analysis predicts the extension $`\xi `$ of the depletion zone to be of order $`(D/F)^{1/4}`$. By comparing $`\xi `$ with the mean distance $`l\rho _x^{1/2}\mathrm{\Gamma }^{i/2(i+2)}`$ between islands, we expect that $`h(x)`$ exhibits a large $`x`$ regime with $`h(x)\mathrm{exp}(x)`$ only for $`i>2`$. For $`i=1`$ we thus are satisfied with a simple power law ansatz $`h(x)=Cx^\varphi `$ for $`xx_{}`$, where $`C`$ and $`x_{}`$ follow from the two conditions imposed on $`h(x)`$, and $`\varphi `$ is a fitting parameter. To test this ansatz we take $`f_0(t)`$ for $`\mathrm{\Gamma }=10^7`$ from Fig. 4 (open symbols or dotted lines) and compare $`f(t)`$ as calculated from eq. (29) (solid lines in Fig. 4) with the corresponding $`f(t)`$ as obtained in the simulation (full symbols in Fig. 4). As can be seen from Fig. 4, for $`\mathrm{\Gamma }=10^6`$ and $`\mathrm{\Gamma }=10^7`$ a fairly good account of the differences between $`f_0(t)`$ and $`f(t)`$ can be obtained by choosing $`\varphi =2`$. However, for $`\mathrm{\Gamma }=10^5`$ the theoretical curve underestimates the fraction of covered islands at large times (where $`f(t)1/2`$). Better agreement between theory and simulation can only be obtained if one would allow $`\varphi `$ to depend on $`\mathrm{\Gamma }`$. Alternatively, we have tried an ansatz for $`h(x)`$ similar to that used by Amar and Family for the scaling function characterizing the island size distribution (see eq. (8)). This ansatz yields comparable results, but is also not successful in accounting for the changes with $`\mathrm{\Gamma }`$. We finally have to note, that at the time when submitting the paper, a more detailed theoretical account for $`\psi (s,\sigma ,t)`$ was published. The use of this finding in eq.(28) and the comparison of the resulting $`f(t)`$ with Monte Carlo data will be presented elsewhere. Having shown that the single-island and multi-island models are essentially equivalent, except for differences between $`f(t)`$ from $`f_0(t)`$ for large times that can be attributed to the island size distribution, we will focus on the single-island model in the remaining part of the paper. ## IV Second layer nucleation in simple situations In this Section we develop a stochastic description of the nucleation process based on the scaling approach for second layer nucleation presented in ref. (see also ref. ). The procedure focuses on the non-interacting particle model, although formally it is possible to extend scaling concepts also to situations, where the lifetimes of unstable clusters become important. This was shown by Krug et al. and is discussed in a more general context in Sec. IV E. The treatment of the non-interacting particle model outlined in Sec. IV A already captures the salient features of the problem in terms of lifetimes, occupation probabilities and encounter rates. We will show that there exist two possible mechanisms for the formation of a stable cluster: In the first case, there is typically no atom on top of the island and a stable cluster is formed due to fluctuations, in which by chance $`i+1`$ atoms are present on the island. In the second case by contrast, there are on average more than $`i+1`$ atoms on top of the island during the formation of a stable cluster so that the nucleation process can be described in a mean-field type manner. It turns out that the fluctuation-dominated case takes place for $`i2`$, while the mean-field situation occurs for $`i3`$. The TDT approach corresponds to the mean-field case with the notable supplement that for very large step edge barriers one should deal with the time-dependent adatom density $`\rho _1(r,t)`$ (solution of eqs. (9,10)) to calculate the nucleation rate $`\mathrm{\Omega }(R)`$ from eq.(16). In the language of critical phenomena, one may regard $`i=2`$ as the upper critical size of the critical nucleus above which mean-field theory becomes applicable. We have to note that the existence of this upper critical size was not perceived by us in ref. , and accordingly, the extension of the scaling arguments for the fluctuation-dominated situation to $`i=3`$ was not allowed. In the stochastic formulation presented below we will develop many of the necessary ingredients for the general treatment of second layer nucleation in the next Sec. V. Moreover, it is discussed under which conditions mean-field type expressions $`D\rho _1^{i+1}`$ for local nucleation rates can be used. ### A General procedure In order to determine a second layer nucleation rate $`\mathrm{\Omega }(R)`$ we start by considering a time interval $`\mathrm{\Delta }t(R)`$, during which $`R(t)`$ does not change significantly. For example, for the generic growth law (18) we may require $`\mathrm{\Delta }t(R)`$ to correspond to a 10% change of $`R`$, which would give $`\mathrm{\Delta }t(R)=0.21(R/a)^2/[A^2Fa^2\mathrm{\Gamma }^{i/(i+2)}]`$, i.e. $$\mathrm{\Delta }t(R)F^1\mathrm{\Gamma }^{i/(i+2)}R^2$$ (30) The nucleation rate $`\mathrm{\Omega }(R)`$ is the mean number $`n_{\mathrm{nuc}}(R)`$ of nucleation events in time $`\mathrm{\Delta }t(R)`$ divided by $`\mathrm{\Delta }t(R)`$, $$\mathrm{\Omega }(R)=\frac{n_{\mathrm{nuc}}(R)}{\mathrm{\Delta }t(R)}.$$ (31) A nucleation event occurs, if $`i+1`$ atoms encounter each other on nearest neighboring sites. For an island with radius $`R`$ and infinite step edge barrier ($`\alpha =0`$), and in total $`n`$ single atoms on top of it, let us approximate the encounter dynamics by a Poisson process, where $`\omega _n(R)`$ denotes the encounter rate of exactly $`i+1`$ atoms. Within the Poisson approximation this rate can be precisely defined as the inverse average time for $`i+1`$ atoms to encounter each other for the first time, when initially $`n`$ atoms are randomly distributed on top of the island. A simple scaling argument yields $$\omega _n(R)=\kappa _\mathrm{e}\left[\underset{k=0}{\overset{i}{}}(nk)\right]\frac{D}{a^2}\left(\frac{a^2}{\pi R^2}\right)^{i+1}\frac{\pi R^2}{a^2},$$ (32) where $`\kappa _\mathrm{e}`$ is a constant. The term $`(a^2/\pi R^2)^{i+1}`$ is proportional to the probability to find $`i+1`$ atoms on nearest neighbor sites, and the factor $`(\pi R^2/a^2)`$ takes into account that the encounter can occur everywhere on the island. The combinatorial factor $`_{k=0}^i(nk)`$ is slightly more subtle. At first sight, one may think that one should include the number $`\left(\genfrac{}{}{0pt}{}{n}{i+1}\right)`$ of possibilities to choose any $`i+1`$ atoms out of the $`n`$ atoms but this is not correct, since the accumulation of $`i+1`$ atoms does not happen “in parallel” at a certain instant of time but in order: First a dimer forms out of $`n`$ single atoms (combinatorial factor $`n(n1)/2`$) and then some of the remaining $`nk`$ atoms ($`k=2,3,\mathrm{},i`$) have to attach one after another to an intermediate cluster of size $`k`$ before this cluster dissociates (the intermediate cluster is assumed to be much less mobile than single adatoms). The sequential attachment process yields an additional combinatorial factor $`_{k=2}^i(nk)`$. Clearly, the scaling argument gives only a rough approximation for $`\omega _n(R)`$ and a more refined treatment justifying eq. (32) is presented in Appendix A. Determination of $`\omega _n(R)`$ for $`i=1,2`$ and various $`n`$ in our simulations confirms the behavior predicted by eq. (32), see Fig. 6. For $`i=1`$ the scaling law is only valid for large $`R100a`$, because at smaller $`R`$, two atoms typically encounter each other before the delta-functions characterizing the initial occupancy smear out to a uniform distribution (for larger $`i`$ this effect becomes less important). Moreover, we find $`\kappa _\mathrm{e}0.087`$ for $`i=1`$ and $`\kappa _\mathrm{e}0.53`$ for $`i=1`$, i.e. the coefficient $`\kappa _\mathrm{e}`$ is constant for fixed $`i`$, but changes strongly with $`i`$. This dependence is expected, since we neglected the memory effect that, when $`n`$ atoms, $`2ni`$, are already close to each other, they keep close together for a while so that the encounter of $`i+1`$ atoms during this intermediate time becomes more likely. This memory effect is not included in the treatment in Appendix A, where after each “dissociation” of an unstable cluster of size $`ki`$ a configuration is assumed to emerge, where a cluster of size $`k1`$ is left and the remaining $`nk`$ atoms are assumed to be randomly distributed. Accordingly, $`\kappa _\mathrm{e}`$ should increase with increasing $`i`$ as it is the case. Equation (32) has been derived for an infinite step edge barrier. For finite step edge barriers, we have to take into account that a state corresponding to an island with $`n`$ atoms on top of it has a finite lifetime $`\tau _n(R)`$ only. This lifetime is defined by the average time required for the first of the $`n`$ atoms to escape from the island (if any encounter processes are neglected). To a good approximation, $`\tau _n(R)`$ is the $`n`$th fraction of the lifetime $`\tau _1(R)`$ of a single atom, $`\tau _n(R)\tau _1(R)/n`$ (this approximation would become exact, if the escape were a simple Poisson process). In the limit of large $`\alpha `$, $`\tau _1(R)`$ is proportional to the characteristic time $`R^2/D`$ for an atom to reach the boundary, while for small $`\alpha `$ an atom typically returns many times to the boundary before escaping from the island. Thus, in the latter limit, the characteristic escape rate (inverse lifetime $`\tau _1^1`$) is approximately given by the product of the probability $`2\pi Ra/\pi R^2`$ for the atom to be at the boundary and the rate $`\alpha D/6a^2`$ to overcome the step edge barrier. Combining these results gives $$\tau _n(R)=\frac{1}{n}\frac{R^2}{D}\left(\kappa _1\frac{a}{\alpha R}+\kappa _2\right),$$ (33) where $`\kappa _1`$ and $`\kappa _2`$ are constants. Indeed, an exact solution of the corresponding diffusion problem allows one to derive $`\tau _n(R)`$ exactly in the continuum limit, as we have shown in Appendix B. In particular, when the escape is approximated by a Poisson process, one finds $`\kappa _11`$ and $`\kappa _2=1/2`$ after proper renormalization and taking into account the lattice corrections (see Appendix B). Direct determination of $`\tau _1(R)`$ in our simulations confirms this result, see Fig. 7. Knowing $`\tau _n(R)`$ we can calculate the probability $`p_n(R)=p_n(R(t))`$ to find exactly $`n`$ atoms on top of the island at time $`t`$ before onset of second layer nucleation. This is achieved by considering the time evolution of $`p_n(R(t))`$, which is described by the master equation $`{\displaystyle \frac{dp_n}{dt}}`$ $`=`$ $`\pi FR(t)^2\left[(1\delta _{n,0})p_{n1}p_n\right]`$ (35) $`+\left[{\displaystyle \frac{p_{n+1}}{\tau _{n+1}(R(t))}}{\displaystyle \frac{p_n}{\tau _n(R(t))}}\right]`$ with the initial condition $`p_n(0)=\delta _{n,0}`$. Note that we have formally introduced $`p_1`$ and that $`1/\tau _nn`$ so that the last term on the right hand side of (35) does not contribute for $`n=0`$. As can be expected and is explicitly shown in Appendix B, the solution of eq. (36) is the Poisson distribution $$p_n(R)=\frac{\overline{n}(R)^n}{n!}\mathrm{exp}[\overline{n}(R)],$$ (36) where the mean number $`\overline{n}(R)`$ of atoms on top of the island before onset of nucleation is $$\overline{n}(R)=\frac{2\pi }{A^2\mathrm{\Gamma }^{\frac{i}{i+2}}}(1+\stackrel{~}{\alpha }R)^\phi _0^R𝑑xx^3(1+\stackrel{~}{\alpha }x)^\phi .$$ (37) Here $`\stackrel{~}{\alpha }\kappa _2\alpha /\kappa _1`$ and $`\phi 2A^2\mathrm{\Gamma }^{2/(i+2)}/\kappa _2`$. An explicit solution after evaluating the integral in eq. (37) is given in eq. (B19) of Appendix B. For fixed $`\alpha `$ and $`\mathrm{\Gamma }`$, three distinct $`R`$ regimes can be identified from eq. (37): For $`\phi \stackrel{~}{\alpha }R1`$ we can use $`(1+\stackrel{~}{\alpha }x)^\phi 1`$ in (37), while for $`\alpha R1`$ but $`\phi \stackrel{~}{\alpha }R1`$ we can use $`(1+\stackrel{~}{\alpha }x)^\phi \mathrm{exp}(\phi \stackrel{~}{\alpha }x)`$. For $`\stackrel{~}{\alpha }R1`$, we can set $`(1+\stackrel{~}{\alpha }R)\stackrel{~}{\alpha }R`$ in (37), and, since the integral over $`x`$ is dominated by the upper bound, $`(1+\stackrel{~}{\alpha }x)\stackrel{~}{\alpha }x`$ also. We thus obtain $$\overline{n}(R)\{\begin{array}{cc}\mathrm{\Gamma }^{i/(i+2)}R^4,\hfill & R/a\mathrm{\Gamma }^{2/(i+2)}\alpha ^1\hfill \\ \mathrm{\Gamma }^1\alpha ^1R^3,\hfill & \mathrm{\Gamma }^{2/(i+2)}\alpha ^1R/a\alpha ^1\hfill \\ \mathrm{\Gamma }^1R^4,\hfill & \alpha ^1R/a\hfill \end{array}$$ (38) The two regimes for large $`R`$ correspond to a quasi-stationary situation ($`dp_n/dt=0`$ in eq. (35)), where $`p_n(R)`$ from eq. (36) equals the stationary distribution for $`R=R(t)`$ with $`\overline{n}(R)=\pi FR^2\tau _1(R)`$. In these regimes the same result (38) can be obtained also by integrating $`\rho _1^{\mathrm{st}}`$ from eq. (11) over the island area. In fact, we used this connection to renormalize the constants $`\kappa _1`$ and $`\kappa _2`$ in eq. (33), see Appendix B. The small $`R`$ regime in eq. (38) corresponds to a non-stationary situation, where $`p_n`$ in general depends on the function $`R(t^{})`$ at all times $`0t^{}t`$ and not only on its value $`R(t)`$ at time $`t^{}=t`$. This fact, however, which also concerns the crossover value $`R_\times \mathrm{\Gamma }^{2/(i+2)}\alpha ^1`$ to the non-stationary regime, is of minor importance here, since we consider the generic growth law (18) throughout the paper. We thus can use $`R`$ and $`t`$ interchangeably. Note that the crossover from the non-stationary to the quasi-stationary situation occurs when $`\tau _1(R_\times )\mathrm{\Delta }t(R_\times )`$, that means in the non-stationary small $`R`$ regime the changes in the radius occur on a faster scale than the escape of an atom from the island, $`\mathrm{\Delta }t(R)\tau _1(R)`$, while in the two quasi-stationary large $`R`$ regimes $`\mathrm{\Delta }t(R)\tau _1(R)`$. Let us now return to the different scenarios discussed in the introductory part of this Section. When $`\overline{n}(R)i+1`$, nucleation of a stable cluster can take place at any instant of time. The number of nucleations in $`\mathrm{\Delta }t(R)`$ that result from states with exactly $`n`$ atoms on top of the island is proportional to $`\omega _n(R)\mathrm{\Delta }t(R)`$. The total number $`n_{\mathrm{nuc}}(R)`$ is the weighted sum of $`\omega _n(R)\mathrm{\Delta }t(R)`$ over $`n`$, i.e. we find $`n_{\mathrm{nuc}}(R)=_{n=i+1}^{\mathrm{}}p_n(R)\omega _n(R)\mathrm{\Delta }t(R)`$ (we are allowed to extend the sum up to infinity due to the sharp decrease of the Poisson distribution for $`n\overline{n}(R)`$). With eq. (31) we thus obtain for the mean-field nucleation rate $`\mathrm{\Omega }_{\mathrm{mf}}(R)`$ $`=`$ $`{\displaystyle \underset{n=i+1}{\overset{\mathrm{}}{}}}p_n(R)\omega _n(R)`$ (39) $`=`$ $`\kappa _\mathrm{e}{\displaystyle \frac{D}{a^2}}\left({\displaystyle \frac{\overline{n}(R)}{\pi R^2}}a^2\right)^{i+1}\left({\displaystyle \frac{\pi R^2}{a^2}}\right).`$ (40) Equation (40) can be interpreted as resulting from a local nucleation rate $`D\rho _1^{i+1}=D[\overline{n}(R)/\pi R^2]^{i+1}`$ integrated over the island area (factor $`\pi R^2`$). Compared to the TDT approach the radial variation of the diffusion profile $`\rho _1=\rho _1(r)`$ is neglected in the stochastic description, so that $`\mathrm{\Omega }(R)`$ from eq. (16) may be preferred over eq. (40). However, as will be discussed further in Sec. IV D below, for large $`\alpha `$ one should use the non-stationary solution of eqs. (9,10) for calculating $`\mathrm{\Omega }(R)`$ from (16) corresponding to the small $`R`$ regime of $`\overline{n}(R)`$ in eq. (38). More important, eq. (40) (or (16)) can be used only if $`\overline{n}(R_c)i+1`$ in the relevant time interval $`\mathrm{\Delta }t(R_c)`$ at the onset of second layer nucleation. The stochastic description allows us to treat also the fluctuation dominated case, where $`\overline{n}(R_c)i+1`$. In this situation $`i+1`$ adatoms have to be deposited and to encounter each other on the island. We can restrict our consideration to the deposition of exactly $`i+1`$ atoms, since for $`\overline{n}(R_c)i+1`$, fluctuations corresponding to more than $`i+1`$ atoms on the island occur with a probability $`_{n=i+2}^{\mathrm{}}p_n(R)<\mathrm{exp}(1)p_{i+1}(R)\overline{n}(R)/(i+2)p_{i+1}(R)`$. If an atom is deposited on the island already containing $`i`$ atoms, we view this as the start of a nucleation trial. The number $`n_{\mathrm{tr}}(R)`$ of nucleation trials in time $`\mathrm{\Delta }t(R)`$ is $`n_{\mathrm{tr}}(R)=\pi FR^2\mathrm{\Delta }t(R)p_i(R)`$. For a trial to be successful, the $`i+1`$ atoms on the island right after its start have to encounter each other before any of the atoms escapes by passing the step-edge barrier. The probability $`p_{\mathrm{enc}}(R)`$ for this to happen is $$p_{\mathrm{enc}}(R)=1\mathrm{exp}[\omega _{i+1}(R)\tau _{i+1}(R)].$$ (41) Accordingly, the total number $`n_{\mathrm{nuc}}(R)`$ of nucleation events in time $`\mathrm{\Delta }t(R)`$ is now $`n_{\mathrm{nuc}}(R)=n_{\mathrm{tr}}(R)p_{\mathrm{enc}}(R)`$, and using eq. (31) we obtain for the fluctuation-dominated nucleation rate $`\mathrm{\Omega }_{\mathrm{fl}}(R)`$ $`=`$ $`\pi FR^2p_i(R)p_{\mathrm{enc}}(R)`$ (42) $`=`$ $`\pi FR^2{\displaystyle \frac{\overline{n}(R)^i}{i!}}e^{\overline{n}(R)}\left(1\mathrm{exp}[\omega _{i+1}(R)\tau _{i+1}(R)]\right).`$ (43) We note that in both formulae (40,42) the only parameter not known a priori is the coefficient $`\kappa _\mathrm{e}`$, which has to be taken from simple simulations of the encounter process (see Fig. 6 and the discussion above). Hence they do not require more input parameters than the expression (16) resulting from the TDT approach. It remains to clarify, when the mean-field or the fluctuation dominated situation occurs, i.e. when $`\mathrm{\Omega }_{\mathrm{mf}}(R)`$ or $`\mathrm{\Omega }_{\mathrm{fl}}(R)`$ has to be used as second layer nucleation rate. The answer to this question can be found by self-consistency requirements: Suppose first that the fluctuation dominated case takes place. Then, using (42), one can calculate the critical radius $`R_c`$ and check if the condition $`\overline{n}(R_c)i+1`$ is fulfilled. In addition the condition $`\omega _{i+1}(R_c)\mathrm{\Delta }t(R_c)1`$ should be fulfilled too, since the encounter of $`i+1`$ atoms in the characteristic time $`\omega _{i+1}(R_c)^1`$ should happen before $`R_c`$ changes. If these necessary conditions for the fluctuation-dominated case are obeyed, then the mean-field situation is ruled out. This conclusion can be drawn, since $`\overline{n}(R)`$ is monotonously increasing with $`R`$, which implies that $`R_c`$ following from $`\mathrm{\Omega }_{\mathrm{fl}}(R)`$ is always smaller than $`R_c`$ resulting from $`\mathrm{\Omega }_{\mathrm{mf}}(R)`$. Hence, when the fluctuations are likely enough to initiate second layer nucleation, they lead to the formation of stable clusters at an earlier time $`t_c`$ than that expected from the mean-field approach. We will now show that the fluctuation-dominated case occurs for $`i=1,2`$. The detailed analysis is a bit technical and the reader, who is interested in the main findings only, may skip the discussion of the various regimes I-IV in the following subsection and proceed with the summary of the results given right after this discussion. ### B Small critical nuclei ($`i=1,2`$) Using $`\mathrm{\Omega }_{\mathrm{fl}}(R)`$ from eq. (42) we can determine the critical radius $`R_c`$ (or, more precisely, $`R_c^{}`$) by calculating $`f_0(t)`$ as in the TDT approach (see eq. (17)). However, for discussing the scaling of $`R_c`$ with $`\mathrm{\Gamma }`$ and $`\alpha `$, it is easier to obtain $`R_c`$ from the condition $$\mathrm{\Omega }_{\mathrm{fl}}(R_c)\mathrm{\Delta }t(R_c)1,$$ (44) which expresses the fact that the probability of second layer nucleation in $`\mathrm{\Delta }t(R_c)`$ becomes of the order of one. Since we consider the fluctuation-dominated case for small critical nuclei here ($`i=1,2`$), we assume $`\overline{n}(R_c)1`$ and thus set $`\mathrm{exp}[\overline{n}(R_c)]1`$, when inserting $`\mathrm{\Omega }_{\mathrm{fl}}(R_c)`$ from eq. (42) into eq. (44). Four different regimes are then predicted by eq. (44): * In the limit $`\alpha 0`$ we have $`\overline{n}(R)\mathrm{\Gamma }^{i/(i+2)}R^4`$ and $`\tau _{i+1}\mathrm{}`$. Hence we obtain from eqs. (30,42,44) $`FR_c^2\mathrm{\Gamma }^{i^2/(i+2)}R_c^{4i}F^1\mathrm{\Gamma }^{i/(i+2)}R_c^2R_c^{4(i+1)}\mathrm{\Gamma }^{i(i+1)/(i+2)}\mathrm{const}.`$, i.e. $$R_c\mathrm{\Gamma }^{i/[4(i+2)]}$$ (45) From (45) follows $`\overline{n}(R_c)\mathrm{const}.`$, which means that the assumption of a fluctuation-dominated situation is not necessarily justified. In fact, eq. (45) appears here as the result of a rather lengthy calculation, but in the limit $`\alpha 0`$, the same scaling behavior (45) can be obtained very simply by calculating the average time needed for the deposition of $`i+1`$ atoms (see ref. ). Hence, despite $`\overline{n}(R_c)i+1`$, eq. (45) gives the correct scaling behavior. However, eq. (45) predicts $`\omega _{i+1}(R_c)\mathrm{\Delta }t(R_c)DR_c^{2i}F^1\mathrm{\Gamma }^{i/(i+2)}R_c^2\mathrm{\Gamma }^{(i^2i4)/[2(i+2)]}`$ and since $`\mathrm{\Gamma }=D/Fa^41`$, the inequality $`\omega _{i+1}(R)\mathrm{\Delta }t(R_c)1`$ becomes violated for $`i3`$. For $`i3`$ therefore, the condition $`\omega _{i+1}(R_c)\mathrm{\Delta }t(R_c)1`$ should be used for calculating $`R_c`$, and because this yields $`\overline{n}(R_c)>i+1`$, one may alternatively use $`\mathrm{\Omega }_{\mathrm{mf}}(R_c)\mathrm{\Delta }t(R_c)1`$ as the determining relation (see Sec. IV D). * With increasing $`\alpha `$, for $`i2`$, either the non-stationarity condition $`\tau _1(R_c)\mathrm{\Delta }t(R_c)`$ ($`\overline{n}(R_c)\mathrm{\Gamma }^{i/(i+2)}R_c^4`$ in eq. (42)) or the condition $`\omega _{i+1}(R_c)\tau _{i+1}(R_c)1`$ ($`p_{enc}1`$ in eq. (42)) breaks down first. Taking $`R_c`$ from eq. (45), the first condition implies $`\alpha \mathrm{\Gamma }^{(i+8)/[4(i+2)]}`$, while the second implies $`\alpha \mathrm{\Gamma }^{i(2i1)/[4(i+2)]}`$. Since the first condition is more restrictive for $`i2`$, regime I ceases to be valid when $`\alpha `$ becomes larger than $`\mathrm{\Gamma }^{(i+8)/[4(i+2)]}`$ and the quasi-stationarity situation is reached. In eq. (42) we now have to take $`\overline{n}(R_c)=\pi FR^2\tau _1(R)\mathrm{\Gamma }^1\alpha ^1R^3`$ (see eq. (38)) and it follows $`\mathrm{\Omega }_{\mathrm{fl}}(R_c)\mathrm{\Delta }t(R_c)\alpha ^i\mathrm{\Gamma }^{i(i+3)/(i+2)}R_c^{3i+4}\mathrm{const}.`$, i.e. $$R_c\alpha ^{i/(3i+4)}\mathrm{\Gamma }^{i(i+3)/[(i+2)(3i+4)]}.$$ (46) Since $`\overline{n}(R_c)(\mathrm{\Gamma }^{(i+8)/4(i+2)}\alpha ^1)^{4/(3i+4)}1`$ the condition for a fluctuation-dominated situation is fulfilled, and since $`\mathrm{\Delta }t(R_c)\tau _1(R_c)\tau _{i+1}(R_c)`$ and $`\omega _{i+1}(R_c)\tau _{i+1}(R_c)1`$ the condition $`\omega _{i+1}(R_c)\mathrm{\Delta }t(R_c)1`$ is obeyed too. * By further increasing $`\alpha `$ we obtain $`\omega _{i+1}(R_c)\tau _{i+1}(R_c)1`$ for $`\alpha \mathrm{\Gamma }^{i(i+3)(2i1)/[2(i+2)(i^2+i+2)]}`$. Hence we now have to use $`p_{enc}\omega _{i+1}(R_c)\tau _{i+1}(R_c)`$ when inserting eq. (42) into eq. (44) and find $$R_c\alpha ^{(i+1)/(i+5)}\mathrm{\Gamma }^{i(i+3)/[(i+2)(i+5)]}.$$ (47) The condition $`\overline{n}(R_c)\mathrm{\Gamma }^1\alpha ^1R_c^31`$ requires $`\alpha ^{i1}\mathrm{\Gamma }^{(i^2+i5)/(i+2)}`$ and is fulfilled for $`i=1`$. For $`i=2`$, it is valid for $`\alpha \mathrm{\Gamma }^{1/4}a/R_c`$. The second requirement $`\omega _{i+1}(R_c)\mathrm{\Delta }t(R_c)1`$ gives $`\alpha ^{i1}\mathrm{\Gamma }^{(i^3+2i^24i5)/[(i+1)(i+2)]}`$ and again is obeyed for $`i=1`$ and valid for $`i=2`$ as long as $`\alpha \mathrm{\Gamma }^{1/4}a/R_c`$. * In this last regime $`\alpha `$ becomes larger than $`a/R_c`$, that means eq. (47) predicts the regime to occur for $`\alpha \mathrm{\Gamma }^{i/[2(i+2)]}`$. Taking $`\overline{n}(R_c)\mathrm{\Gamma }^1R_c^4`$ from eq. (38) and $`\omega _{i+1}(R_c)\tau _{i+1}(R_c)R_c^{2(i1)}`$ from eqs. (32,33), we find $$R_c\mathrm{\Gamma }^{i/[2(i+2)]}.$$ (48) We used $`\omega _{i+1}(R_c)\tau _{i+1}(R_c)1`$ ($`p_{\mathrm{enc}}(R_c)1`$) to derive (48), which for $`i=2`$ is valid and for $`i=1`$ is obeyed when taking into account the prefactors (for $`i=1`$ $`\omega _{i+1}(R_c)\tau _{i+1}(R_c)=\kappa _\mathrm{e}\kappa _2`$). Moreover, eq. (48) gives $`\overline{n}(R_c)\mathrm{\Gamma }^{(i2)/(i+2)}`$, which is much smaller than one for $`i=1`$. For $`i=2`$, a decision on whether the fluctuation-dominated or the mean-field situation occurs would require a closer inspection of the prefactors. However, since for $`\alpha a/R_c`$ one finds the same scaling (48) in the mean-field situation (see Sec. IV D), eq. (48) is valid in any case. The second condition $`\omega _{i+1}(R_c)\mathrm{\Delta }t(R_c)1`$ is fulfilled for $`i=1`$, and for $`i=2`$ the situation again depends on the prefactors. In summary we have found that the second layer nucleation for $`i=1,2`$ occurs due to various mechanisms in four distinct regimes I-IV: In regime I ($`\alpha \mathrm{\Gamma }^{(i+8)/[4(i+2)]}`$), the nucleation takes place once $`i+1`$ have been deposited on the island, in regime II ($`\mathrm{\Gamma }^{(i+8)/[4(i+2)]}\alpha \mathrm{\Gamma }^{i(i+3)(2i1)/[2(i+2)(i^2+i+2)]}`$) the loss of atoms becomes important and the nucleation takes place once the probability for finding $`i+1`$ atoms on the island at some time instant in $`\mathrm{\Delta }t(R)`$ becomes of the order of one, in regime III ($`\mathrm{\Gamma }^{i(i+3)(2i1)/[2(i+2)(i^2+i+2)]}\alpha \mathrm{\Gamma }^{i/[2(i+2)]}`$) the probability $`p_{\mathrm{enc}}`$ for the encounter of $`i+1`$ atoms during a nucleation trial has to be taken into account in addition to the probability for the occurrence of $`i+1`$ atoms, and in regime IV ($`\alpha \mathrm{\Gamma }^{i/[2(i+2)]}`$) both the occurrence and encounter probability matter but these probabilities no longer depend on the step edge barrier. For convenient reference, we provide the exponents $`\gamma `$ and $`\mu `$ defined in eq. (23) and their corresponding ranges of validity in Table I. When comparing the scaling in the fluctuation-dominated situation with that predicted by eqs. (24,25) of the TDT approach it is remarkable that the same behavior is found in regime III and IV. We believe this to be caused by the fortunate circumstance that local nucleation rates of form $`D\rho _1^{i+1}`$ might be effectively applicable even if the requirements for the mean-field situation are not fulfilled (see however ). Moreover, we have to note that for $`i=1`$ the lower and upper crossovers $`\alpha _2`$ and $`\alpha _3`$ specifying regime III (see table I) both scale as $`\mathrm{\Gamma }^{1/6}`$. This is due to the fact that when $`\omega _{i+1}(R_c)\tau _{i+1}(R_c)`$ becomes less than one, we already obtain $`\alpha a/R_c`$, which is the condition for regime IV. Nevertheless, due to the pronounced small $`R`$ corrections to eq. (32) for $`i=1`$ (see Fig. 6) both conditions $`\omega _{i+1}(R_c)\tau _{i+1}(R_c)1`$ and $`\alpha a/R_c`$ can be fulfilled in a small transient regime III. However, for $`i=1`$ this is no longer a true scaling regime, where a simple power law dependence of $`R_c`$ on $`\mathrm{\Gamma }`$ and $`\alpha `$ can be identified. In this respect, the small $`\alpha `$ regime of the TDT approach does not occur for $`i=1`$, also not at larger $`\alpha `$. ### C Comparison with simulations for i=1,2 Taking $`\mathrm{\Omega }_{\mathrm{fl}}`$ from eq. (42) we can calculate $`f_0(t)`$ according to eq. (17). Representative results for $`i=1`$ are shown in Fig. 4 (solid lines), and the comparison with the Monte Carlo data yields a very good agreement. The $`R_c`$ values derived from $`f_0(t)`$ are plotted as a function of $`\alpha `$ for $`\mathrm{\Gamma }`$ val- ues in the range $`10^510^8`$ in Fig. 8. Note that, compared to the results of the full island model shown in Fig. 5, the data cover the full $`\alpha `$ range from zero to one, since the restrictions imposed by island coalescence in the multi-island model are not present in the single-island model (see also the discussion in Sec. III above). Moreover the simpler single-island model allows one to explore the behavior for larger $`\mathrm{\Gamma }`$ values in the range $`\mathrm{\Gamma }=10^910^{12}`$ also. It is possible to fit the $`R_c`$ curves over the entire range of $`\alpha `$ and $`\mathrm{\Gamma }`$ values (see Sec. V B) but we focus on the scaling behavior of $`R_c`$ in the following in order to demonstrate the various scaling regimes associated with the different physical mechanisms of second layer nucleation. Indeed, the simulated data in Fig. 8 confirm the theoretical predictions. For small $`\alpha \alpha _1`$ (regime I), $`R_c`$ is independent of $`\alpha `$, while for $`\alpha \alpha _1\mathrm{\Gamma }^{3/4}`$ (regime II) we find $`R_c\mathrm{\Gamma }^{4/21}\alpha ^{1/7}`$. Since $`R_c(\alpha _1)\alpha _1^{1/9}`$ at the crossover, the boundary line between regimes I and II has slope (-1/9). The correctness of the scaling of $`R_c`$ with $`\mathrm{\Gamma }`$ in regimes I and II can be deduced from the offset of the curves for various $`\mathrm{\Gamma }`$ in Fig. 8 (alternatively, one can collapse the data onto a common master curve by a proper rescaling as it was shown in ref. ). Regime II is followed by the transient regime III, and the dashed border line separating regime III from regime II was determined numerically from the condition $`\omega _2(R_c)\tau _2(R_c)1`$ by using the results for $`\omega _2(R_c)`$ and $`\tau _2(R_c)`$ displayed in Figs. 6,7. For $`\alpha \mathrm{\Gamma }^{i/[2(i+2)]}`$ (regime IV), we find $`R_c\mathrm{\Gamma }^{i/[2(i+2)]}`$ independent of $`\alpha `$, and the boundary line between regimes III and IV has slope (-1). Figure 9 depicts the various regions characterizing the mechanism of second layer nucleation for $`i=1`$ in an $`\alpha \mathrm{\Gamma }`$ diagram. Varying $`\mathrm{\Gamma }`$ and $`\alpha `$ within one of the regions results in the corresponding behavior of $`R_c`$ according to eqs. (45,46,48). The border line between regions I and II has slope $`(4/3)`$, between regions III and IV slope $`(6)`$, and the dashed line marks the border line between regions II and III. In addition, we have drawn the transition line from rough multilayer to smooth layer-by-layer growth into the diagram. In our simulations island coalescence occurs in regime II (see Fig. 5), where $`R_c\mathrm{\Gamma }^{4/21}\alpha ^{1/7}`$. The criterion $`R_cl\rho _x^{1/2}\mathrm{\Gamma }^{1/6}`$ thus yields $`\alpha _{}(\mathrm{\Gamma })\mathrm{\Gamma }^{1/6}`$. Results for $`R_c`$ obtained from simulations for a critical nucleus of size $`i=2`$ are shown in Fig. 10. Again the results confirm the predictions of the theory. In particular, for large $`\mathrm{\Gamma }`$, the exponents $`\mu =1/5`$ in regime II and $`\mu =3/7`$ in regime III can be clearly identified. In contrast to the behavior for $`i=1`$ shown in Fig. 8, regime III develops into a full scaling regime. ### D Large critical nuclei ($`i3`$) Analogous to the fluctuation-dominated case treated in the previous subsection we can obtain the scaling of $`R_c`$ with $`\mathrm{\Gamma }`$ and $`\alpha `$ from the condition $`\mathrm{\Omega }_{\mathrm{mf}}(R_c)\mathrm{\Delta }t(R_c)1`$ with $`\mathrm{\Omega }_{\mathrm{mf}}(R)`$ and $`\mathrm{\Delta }t(R_c)`$ from eqs. (40,30), respectively. For critical island radii belonging to the two quasi-stationary large $`R`$-regimes in eq. (38) this gives the same behavior (24,25) as predicted by the TDT approach. However, for large step edge barriers corresponding to the non-stationary small $`\alpha `$-regime in eq. (38) we find $`\mathrm{\Omega }_{\mathrm{mf}}(R_c)\mathrm{\Delta }t(R_c)D(\mathrm{\Gamma }^{i/(i+2)}R_c^4)^{i+1}R_c^{2i}F^1\mathrm{\Gamma }^{i/(i+2)}R_c^2R_c^{2(i+3)}\mathrm{\Gamma }^{(i1)}\mathrm{const}.`$, i.e. $$R_c\mathrm{\Gamma }^{(i1)/[2(i+3)]}.$$ (49) With increasing $`\alpha `$ this scaling breaks down when $`R_c`$ enters the quasi-stationary regime in eq. (38) that means for $`\alpha \mathrm{\Gamma }^{2/(i+2)}R_c^1\mathrm{\Gamma }^{(i^2+5i+10)/[2(i+2)(i+3)]}`$. For $`i3`$ we thus have in total three distinct regimes I-III with different mechanisms for second layer nucleation: In regime I ($`\alpha \mathrm{\Gamma }^{(i^2+5i+10)/[2(i+2)(i+3)]}`$) the nucleation takes place once the island radius $`R`$ has grown large enough so that the encounter of $`i+1`$ atoms out of typically $`\overline{n}(R)i+1`$ atoms happens in a time comparable to $`\mathrm{\Delta }t(R)`$, in regime II $`\overline{n}(R)`$ becomes dependent on $`\alpha `$, while in regime III, for large $`\alpha a/R_c`$, $`\tau _1(R)`$ no longer depends on the step edge barrier and $`\overline{n}(R)`$ becomes independent of $`\alpha `$ again. The overall behavior characterized by the scaling exponents $`\gamma `$ and $`\mu `$ is summarized in Table II. Computer simulations for $`i=3`$ are in accordance with these theoretical predictions, see Fig. 11. The predicted scaling $`R_c\alpha ^{1/2}`$ in regime II is not yet fully developed for the $`\mathrm{\Gamma }`$-values in the range $`10^510^8`$ but it can be expected to become more clearly visible for larger $`\mathrm{\Gamma }`$. However, we could not obtain reliable simulation results for larger $`\mathrm{\Gamma }`$ values, since the amount of CPU time for determining the onset of second layer nucleation becomes tremendous due to the increasing number of atoms contributing to the nucleation event. ### E Influence of metastable clusters To demonstrate how the presence of metastable nuclei may be included into the general procedure presented in Sec. IV A, we consider, as in ref. , the simplest case of second layer nucleation of a trimer ($`i=2`$), when a dimer is metastable with characteristic dissociation time $`\tau _{\mathrm{dis}}`$. For $`i=2`$ we have to deal with the fluctuation-dominated situation. We note in passing that this can be true even for larger $`i`$ when metastable clusters can form, since their presence tends to drive second layer nucleation into the fluctuation-dominated situation. In contrast to the discussion leading to eq. (42) for the non-interacting particle model, the formation of the stable trimer is not necessarily the rate limiting process. It is possible that the dissociation time $`\tau _{\mathrm{dis}}`$ becomes so large that the nucleation happens effectively instantaneously once the dimer has formed. To decide whether the formation of the stable trimer or metastable dimer is rate limiting, we have to compare $`p_1(R)p_{\mathrm{enc}}^{\left(2\right)}(R)`$ with $`p_2(R)p_{\mathrm{enc}}^{\left(3\right)}(R)`$, where $`p_{\mathrm{enc}}^{\left(\mathrm{j}\right)}(R)`$ denotes the encounter probability of $`j`$ atoms (in Sec. IV A no superscript (j) was introduced, since only $`j=i+1`$ had to be considered). Hence we write $$\mathrm{\Omega }_{\mathrm{fl}}(R)=\{\begin{array}{cc}\pi FR^2p_2(R)p_{\mathrm{enc}}^{\left(3\right)}(R),\hfill & p_1p_{\mathrm{enc}}^{\left(2\right)}p_2p_{\mathrm{enc}}^{\left(3\right)}\hfill \\ \pi FR^2p_1(R)p_{\mathrm{enc}}^{\left(2\right)}(R),\hfill & p_1p_{\mathrm{enc}}^{\left(2\right)}p_2p_{\mathrm{enc}}^{\left(3\right)}\hfill \end{array}$$ (50) To calculate the occupation probabilities $`p_n(R)`$, $`n3`$, we first need to know the modified lifetimes $`\tau _n^{}(R)`$ of states with exactly $`n`$ atoms on top of the island. Clearly, $`\tau _1^{}(R)=\tau _1(R)`$ with $`\tau _1(R)`$ from eq. (33), since the metastable dimer has no influence on the lifetime of a single atom. The characteristic time $`\tau _2^{}(R)`$, however, will be enlarged in comparison to $`\tau _2(R)`$ from eq. (33) and can be estimated as follows (we disregard any prefactors): As in ref. we consider the first deposited atom as immobile and the second deposited atom as diffusing. Once the second atom has been deposited it needs a time of order $`(R^2/D+\tau _{\mathrm{dis}})`$ to reach the step edge, since one encounter with the first deposited atom typically takes place during one traversal of the island within time $`\tau _{\mathrm{tr}}R^2/D`$. At the boundary the second atom is “reflected” a typical number $`M\alpha ^1`$ of times before leaving the island. Between all reflections, the overall elapsed time is of order $`(MRa/D+m\tau _{\mathrm{dis}})`$, where $`Ra/D`$ is the typical time for a single atom to return to the edge and $`m(MRa/D)/\tau _{\mathrm{tr}}a/\alpha R`$ is the typical number of times the second atom encounters the first atom. Summing up all time contributions we obtain (neglecting the prefactors belonging to the four individual terms) $$\tau _2^{}(R)(\tau _{\mathrm{tr}}+\tau _{\mathrm{dis}})+\frac{a}{\alpha R}(\tau _{\mathrm{tr}}+\tau _{\mathrm{dis}}).$$ (51) Note that for $`\tau _{\mathrm{dis}}\tau _{\mathrm{tr}}=R^2/D`$, $`\tau _2^{}(R)`$ reduces to $`\tau _2(R)`$ from eq. (33) (without prefactors). To estimate $`\tau _3^{}(R)`$ we note that if the dimer state is the prevalent one, $`\tau _3^{}(R)\tau _1(R)`$, whereas, if all three atoms are likely to be separated, $`\tau _3^{}(R)\tau _3(R)`$. Since $`\tau _3(R)\tau _1(R)`$, we find $`\tau _3^{}(R)\tau _1(R)`$ in either case. In the strong barrier limit $`\alpha a/R`$, in particular, the first two terms on the right hand side of eq. (51) can be neglected, and, since $`\tau _1Ra/D\alpha \tau _{\mathrm{tr}}a/\alpha R`$, we can simply write $`\tau _2^{}\tau _1(1+\tau _{\mathrm{dis}}/\tau _{\mathrm{tr}})`$ which agrees with the result derived in ref. . This finding for strong step edge barriers implies that for $`\tau _{\mathrm{dis}}/\tau _{\mathrm{tr}}1`$ the two atoms are effectively always in the dimer state and $`\tau _2^{}\tau _2\tau _{\mathrm{dis}}/\tau _{\mathrm{tr}}`$, while for $`\tau _{\mathrm{dis}}/\tau _{\mathrm{tr}}1`$ they are effectively always separated and $`\tau _2^{}\tau _2`$. When inserting the modified lifetimes into eq. (35) and neglecting states with $`n>i+1=3`$ ($`p_n0`$ for $`n>3`$ until onset of nucleation in the fluctuation-dominated situation), we can calculate the occupation probabilities $`p_n(t)`$. In the quasi-stationary limit ($`dp_n/dt=0`$ but $`R=R(t)`$), in particular, we obtain (for $`0ni+1=3`$) $$p_n=\frac{_{j=1}^nq_j}{_{j=0}^{i+1}_{k=1}^jq_k},q_j\pi FR^2\tau _j^{}.$$ (52) To calculate the encounter probabilities $`p_{\mathrm{enc}}^{\left(\mathrm{n}\right)}(R)=1\mathrm{exp}[\omega _n^{}(R)\tau _n^{}(R)]`$, $`n=2,3`$, we furthermore need to know the modified encounter rates $`\omega _n^{}(R)`$. From eq. (A8) in Appendix A we find $`\omega _2^{}=w_1`$ (eq. (A8) for $`i=1`$) and $`\omega _3^{}w_1w_2/v_2^{}`$ (eq. (A8) for $`i=2`$), where $`w_1w_2\tau _{\mathrm{tr}}^1`$ from eq. (A) and $`v_2^{}=\tau _{\mathrm{dis}}^1`$ (modification of eq. A4), i.e. $`\omega _2^{}1/\tau _{\mathrm{tr}}`$ and $`\omega _3^{}\tau _{\mathrm{dis}}/\tau _{\mathrm{tr}}^2`$. To discuss eq. (50) we may now distinguish various cases depending on whether we have to consider (i) the non-stationary or quasi-stationary situation, (ii) the strong ($`\alpha a/R`$) or weak barrier ($`\alpha a/R`$) limit, (iii) the formation of the metastable dimer or stable trimer as rate limiting, (iv) the encounter processes to be faster or slower than the escape process ($`\omega _n^{}\tau _n^{}1`$ or not for $`n=2,3`$), and (v) $`\tau _2^{}`$ to be dominated by the metastable dimer state ($`\tau _2^{}\tau _{\mathrm{dis}}`$ in the weak barrier limit and $`\tau _2^{}\tau _1\tau _{\mathrm{dis}}/\tau _{\mathrm{tr}}`$ in the strong barrier limit) or to be dominated by the state of separated atoms ($`\tau _2^{}\tau _2`$). Rather than treating all these possible cases (and analyzing their possible occurrence for the generic growth law (18) by employing self-consistency requirements) we only remark that the results obtained by Krug et al. are entailed in our description. In this work, certain regimes corresponding to the quasi-stationary case in the strong barrier limit are considered for both $`q_1q_31`$ and $`q_1q_21`$, where we obtain $`p_1q_1`$ and $`p_2q_1q_2`$ from eq. (52). Since $`\omega _1^{}\tau _1^{}=\tau _1/\tau _{\mathrm{tr}}a/\alpha R1`$ in the strong barrier limit, we can always set $`p_{\mathrm{enc}}^{\left(2\right)}1`$ in eq. (50). The following regimes are then discussed in ref. with increasing $`\tau _{\mathrm{dis}}`$: * For $`\tau _{\mathrm{dis}}\tau _{\mathrm{tr}}^2/\tau _1`$ we have $`p_{\mathrm{enc}}^{\left(3\right)}\omega _3^{}\tau _3^{}\tau _{\mathrm{dis}}\tau _1/\tau _{\mathrm{tr}}^2`$ and $`\tau _2^{}\tau _1`$, i.e. $`p_2q_1^2`$ and $`p_1p_2p_{\mathrm{enc}}^{\left(3\right)}`$. Accordingly, we obtain $`\mathrm{\Omega }_{\mathrm{fl}}(\pi FR^2\tau _1)^3\tau _{\mathrm{dis}}/\tau _{\mathrm{tr}}^2`$ corresponding to eq. (14) (regime I) in ref. . * For $`\tau _{\mathrm{tr}}^2/\tau _1\tau _{\mathrm{dis}}\tau _{\mathrm{tr}}`$, we find $`p_{\mathrm{enc}}^{\left(3\right)}1`$ and $`\tau _2^{}\tau _1`$, i.e. $`p_2q_1^2`$ and $`p_1p_2p_{\mathrm{enc}}^{\left(3\right)}`$ as in (i), and hence $`\mathrm{\Omega }_{\mathrm{fl}}(\pi FR^2)^3\tau _1^2`$ corresponding to eq. (15) (regime II) in ref. . In the following cases, where $`\tau _{\mathrm{dis}}`$ becomes even larger (and $`\tau _1`$, $`\tau _{\mathrm{tr}}`$ do not change), we still have $`p_{\mathrm{enc}}^{\left(3\right)}1`$. * For $`\tau _{\mathrm{tr}}\tau _{\mathrm{dis}}\tau _{\mathrm{tr}}/\pi FR^2\tau _1`$, $`\tau _2^{}\tau _1\tau _{\mathrm{dis}}/\tau _{\mathrm{tr}}`$, i.e. $`p_2q_1q_2(\pi FR^2)^2\tau _1^2\tau _{\mathrm{dis}}/\tau _{\mathrm{tr}}`$. The condition $`\tau _{\mathrm{dis}}\tau _{\mathrm{tr}}/\pi FR^2\tau _1`$ is equivalent to $`p_1p_2p_{\mathrm{enc}}^{\left(3\right)}p_2`$, and we thus find $`\mathrm{\Omega }_{\mathrm{fl}}\pi FR^2p_2(\pi FR^2)^3\tau _1^2\tau _{\mathrm{dis}}/\tau _{\mathrm{tr}}`$ corresponding to eq. (16) (regime III) in ref. . * For $`\tau _{\mathrm{tr}}/\pi FR^2\tau _1\tau _{\mathrm{dis}}`$ finally, $`p_1p_2`$ and the formation of the stable trimer is no longer the rate limiting process. From eq. (50) we then obtain $`\mathrm{\Omega }_{\mathrm{fl}}\pi FR^2p_1(\pi FR^2)^2\tau _1`$ corresponding to eq. (17) (regime IV) in ref. . As expected, the scaling behavior of $`\mathrm{\Omega }_{\mathrm{fl}}`$ in this limit reduces to the case $`i=1`$ (see also the discussion in Sec. V C). It is clear that the above analysis is difficult to extend to even more complicated situations. Moreover, due to the growing number of characteristic time scales, we found it increasingly difficult to discern pronounced scaling regimes in practice, see e.g. Fig. 18. We therefore prefer to treat the problem of second layer nucleation in the presence of metastable clusters within the more general framework outlined in the following Section. ## V Second layer nucleation in general situations In a more general approach to the problem of second layer nucleation we distinguish between different individual states of the island during its growth with respect to the number of atoms that are on top of the island and the way a given number of atoms is decomposed into clusters of various sizes. Employing a Poisson approximation, the transition processes between the states exhibit no (intrinsic) memory and can be characterized by elementary rates. For the non-interacting particle model these elementary transition rates are the deposition rate $`\pi FR^2`$, the rate for the attachments of a single atom to an intermediate cluster of size $`k`$, and the loss rate of adatoms. The latter is given by the inverse lifetime $`\tau _1^1`$ of a single atom (see eq. (33)). Dissociation rates enter the problem as additional parameters, when the lifetimes of intermediate metastable clusters can not be neglected. The consequences of such dissociation rates will be discussed in Sec. V C. First, however, we will present the general procedure in part V A and show in part V B how the results of the simplified stochastic description in Sec. IV can be recovered. In addition to these previously derived results, it is also discussed how the general treatment allows one to gain detailed insight into the dominant microscopic pathways that are followed to form a stable nucleus on top of the island. ### A General Procedure Let us introduce a common notation for the elementary transition rates: $`W_\mathrm{F}`$ for the deposition rate, $`W_\mathrm{l}`$ for the loss rate, $`W_{\mathrm{a},j}^{(n)}`$ for the attachment rate for a single atom to an intermediate cluster of size $`j`$ if in total $`n`$ atoms are present on top of the island (see eq. A in Appendix A; we formally include the case $`j=1`$), and $`W_{\mathrm{d},j}`$ for the dissociation rate of an unstable cluster composed of $`ji`$ atoms (again we do not distinguish between different cluster configurations for the same cluster size, see also the remark in ref. ). According to the results derived in Sec. IV and Appendix A the transition rates are $$W_\mathrm{F}=\pi FR^2,$$ (54) $$W_\mathrm{l}=\frac{D}{R^2}\left(\kappa _1\frac{a}{\alpha R}+\kappa _2\right)^1,$$ (55) $$W_{\mathrm{a},j}^{(n)}=\kappa _{\mathrm{a},\mathrm{j}}^{(n)}\{\begin{array}{cc}\frac{n(n1)}{2}\frac{2D}{\pi R^2},\hfill & j=1\\ (nj)\frac{D}{\pi R^2},\hfill & 2jn\end{array}$$ (56) $$W_{\mathrm{d},j}=\kappa _{\mathrm{d},\mathrm{j}}\frac{D}{a^2}\mathrm{exp}\left(\frac{\mathrm{\Delta }E_j^{\mathrm{dis}}}{k_\mathrm{B}T}\right),2ji,$$ (57) where $`\mathrm{\Delta }E_j^{\mathrm{dis}}=E_j^{\mathrm{dis}}E_0`$ is the dissociation energy of a single atom from an unstable cluster of size $`ji`$. The prefactors $`\kappa _{\mathrm{a},\mathrm{j}}^{(n)}`$ and $`\kappa _{\mathrm{d},\mathrm{j}}`$ contain the effective sizes of cluster perimeters on one hand (see the discussion in Appendix A), and various corrections involved in the overall approximation scheme (Poisson approximation, cutoff $`n_{}`$ introduced below, etc.); they are considered to be independent of $`D`$, $`\alpha `$ and $`R`$. In principle one should also take into account the possibility that a subcluster composed of more than one atom can dissociate from an unstable cluster. In fact, it has been argued that such cluster dissociations are sometimes more likely to occur than the dissociation of single atoms, as e.g. for dimer dissociation from a tetramer on a (100) surface by a kind of “shearing mode”. For simplicity we will take into account only the dissociation of single atoms here, although conceptually the inclusion of cluster dissociation processes into the general treatment poses no difficulty. Also, we do not consider the influence of cluster mobilities. If one would allow for a small jump rate $`D_j/a^2`$ of a cluster of size $`j2`$, the relative diffusion of a $`j`$ cluster and a single atom would be larger by a factor $`1+D_j/D`$ and accordingly we had to multiply $`W_{\mathrm{a},j}^{(n)}`$ in eq. (56) by this factor for $`j2`$. The method is best introduced by an example. To this end consider Fig. 12 that illustrates the situation for a critical nucleus of size $`i=2`$. Various states of the island are shown, which are distinguished according to the total number $`n`$ of atoms on top of the island, and the possible configurations that can be assumed for a given $`n`$. Between the states the possible transitions are marked by arrows that are labeled by the corresponding rates. Note that the loss from a state with $`n`$ single atoms is $`n`$ times larger than the loss from the state with one atom. It is clear that Fig. 12 shows only a small part of the possible states and in principle can be extended by including larger numbers $`n`$. However, as will be pointed out below, these states with larger $`n`$ do not contribute much to the onset of second layer nucleation. Moreover, we have not included states containing stable clusters of size $`j>i+1`$ and transitions between different states containing a stable nucleus of size $`i+1`$. These are irrelevant for the fraction $`f_0(t)`$ of covered islands at time $`t`$. We denote by $`p_{n,\nu }`$ the probability for the island to be in state $`(n,\nu )`$, where $`n`$ refers to the number of atoms on top of the island and $`\nu `$ to a specific configuration for a given $`n`$. A complete description of the stochastic process amounts to specifying the set $`\{p_{n,\nu }(t)\}`$ of state probabilities at all times $`t`$. The time evolution of the $`\{p_{n,\nu }(t)\}`$ is described by the master equation $`{\displaystyle \frac{dp_{n,\nu }}{dt}}`$ $`=`$ $`{\displaystyle \underset{n^{},\nu ^{}}{}}[W(n^{},\nu ^{}n,\nu )p_{n^{},\nu ^{}}`$ (59) $`W(n,\nu n^{},\nu ^{})p_{n,\nu }],`$ where for the rates $`W(n,\nu n^{},\nu ^{})`$ the appropriate expressions from eq. (V A) have to be substituted (see Fig. 12). Note that transitions are possible only between a limited number of states. In the situation considered here, where only single atoms can leave the island, we have $`W(n,\nu n^{},\nu ^{})=0`$ for $`\|nn^{}\|2`$. To treat the problem of second layer nucleation under generic growth conditions one has to solve the set of eqs. (59) for $`R=R(t)`$ with $`R(t)`$ from eq. (18) subject to the initial condition $`p_{n,\nu }=\delta _{n,0}`$. To this end it is convenient to solve (59) using $`R`$ as the independent variable. The integration of the differential equations (59) using standard solvers takes very little CPU time on ordinary workstations, so that results for $`f_0(t)`$ and $`R_c(\alpha )`$ can be obtained almost immediately. Numerical results are discussed in the following. ### B Negligible Lifetimes of Unstable Clusters In this subsection we consider the case $`\mathrm{\Delta }E_j^{\mathrm{dis}}=0`$ that was treated extensively in Sec. IV. The fraction $`f_0(t)`$ of covered islands within our more general framework is given by $$f_0(t)=\underset{n=i+1}{\overset{\mathrm{}}{}}p_{n,\nu _n}(t),$$ (60) where $`\nu _n`$ is the configuration containing a stable nucleus for a given $`n`$ (for example, $`\nu _3=3`$ and $`\nu _4=4`$ in Fig. 12). In practice, states corresponding to large $`n`$ contribute a negligible amount up to times $`t_c`$, so that one needs to consider a finite maximum number of atoms $`n_{}`$ only ($`n_{}=4`$ in Fig. 12 turned out to be sufficient). Figure 13a shows $`f_0(t)`$ and the probabilities $`p_{n,\nu }(t)`$ (labeled according to Fig. 12) as a function of the coverage $`Fa^2t`$ for $`\alpha =10^4`$, and $`\mathrm{\Gamma }=10^6`$. Also shown is the mean total number $$N(t)\underset{n=1}{\overset{n_{}}{}}\underset{\nu \nu _n}{}p_{n,\nu }(t)n$$ (61) of atoms that are not in states possessing a stable nucleus. In accordance with the predictions of the simplified stochastic description, this number is less than one up to time $`t_c`$. Accordingly, the pathway followed by the system to form a stable nucleus is dominated by fluctuations as discussed in Sec. IV. The important role of the fluctuations can even more clearly be recognized by looking at the state probabilities $`p_{n,\nu }`$, $`\nu \nu _n`$, and the currents $$j_3(t)W_{\mathrm{a},2}^{(3)}p_{3,2}(t),j_4(t)W_{\mathrm{a},2}^{(4)}p_{4,2}(t)$$ (62) into the states containing a stable nucleus. As can be seen from Fig. 13a, only the probabilities $`p_{n,1}(t)`$ are significant, while the other state probabilities $`p_{n,2}(t)`$ and $`p_{4,3}(t)`$ cannot be discerned on the scale used in the figure. On the other hand, we find that the current $`j_3(t)`$ from the state $`(n=3,\nu =2)`$ (which has a very small probability $`p_{3,2}(t)`$) contributes most to the growth of $`f_0(t)`$, see Fig. 13b. The fact that $`j_4(t)`$ gives only a subdominant contribution to second layer nucleation, indicates that the incorporation of states with $`n>5`$ will not significantly change the behavior of $`f_0(t)`$. The results for $`f_0(t)`$ compare well with the data obtained from kinetic Monte Carlo simulations, the quality of agreement between theory and simulation being as good as in Fig. 4. The values of the optimal prefactors $`\kappa _{\mathrm{a},\mathrm{j}}^{(n)}`$ and $`\kappa _{\mathrm{d},2}`$ are given in ref. . To exemplify the good agreement between theory and simulations, we have re-plotted in Fig. 14 the critical radius $`R_c`$ as a function of $`\alpha `$ for various $`\mathrm{\Gamma }`$ for $`i=1,2`$ from Figs. 8,10. The solid lines referring to the numerical results give an excellent fit to the Monte Carlo data. For $`i=1`$, only the states with $`n2`$ in Fig. 12 had to be included to achieve this almost perfect agreement. For $`i>2`$ we expect that a very large number $`n_{}`$ has to be chosen in order to obtain a correct description of second layer nucleation within the rate equation approach. Diagrams corresponding to that shown in Fig. 12 then become very complicated and not easily tractable from the practical point of view. It is thus helpful to introduce the “renormalized” encounter rates $`\omega _n`$ defined in eq. (32) and to consider simplified diagrams as shown in Fig. 15 for $`i=3`$. For a given number $`ni+1`$ of atoms on top of the island we have only included two states $`\nu =1,2`$: One of these refers to a state where all $`n`$ atoms are separated ($`\nu =1`$), and the other one to a state, where exactly $`i+1`$ atoms form a stable nucleus, while the remaining $`n(i+1)`$ atoms are not bound to other atoms in the same layer $`(\nu =2)`$. Plots of $`f_0(t)`$, $`p_{n,1}(t)`$, $`N(t)`$, and $`j_n(t)p_{n,1}(t)\omega _n(R(t))`$ for $`i=3`$ analogous to Fig. 13 are shown in Fig. 16. As expected from the discussion in Sec. IV, we now had to take into account states with $`n`$ up to $`n_{}=50i(i+1)/2=6`$ before reaching the limit, where $`f_0(t)`$ as calculated from eq. (60) did not change much by incorporation of states with larger $`n`$. Near $`t_c`$, $`N(t)`$ is significantly larger than $`i+1=4`$ (at $`t_c`$ we find $`N(t_c)10`$), and the dominant currents $`j_n(t)`$ initiating second layer nucleation are those for $`n=15204`$, see Fig. 16. In order to see how the preferred paths for second layer nucleation change with the step crossing probability $`\alpha `$, we define the integrated current $`j_n(t)`$ up to $`t_c`$ by $$J_n_0^{t_c}𝑑tj_n(t)=_0^{t_c}𝑑tp_{n,1}(t)\omega _n(R(t)),ni+1.$$ (63) This quantity equals the fraction of covered islands at time $`t_c`$ for which the stable nucleus originates from a state possessing exactly $`n`$ adatoms. Figure 17 shows $`J_n`$ as a function of $`n`$ for fixed $`\mathrm{\Gamma }=10^8`$ and various $`\alpha `$. We see that for all $`\alpha `$ states with $`n4`$ dominate the onset of the nucleation. The number of particles $`n_{\mathrm{peak}}`$ in the state where $`J_n`$ has a maximum strongly increases with increasing $`\alpha `$. For $`\alpha =10^3`$, the second layer nucleation is typically initiated by $`n_{\mathrm{peak}}18`$ adatoms on top of the island. ### C Influence of Metastable Clusters The general procedure outlined in Sec. V A allows us also to describe situations, where the binding energies of unstable clusters of size $`ji`$ are not small compared to $`k_\mathrm{B}T`$. To demonstrate this we again consider the case $`i=2`$ and the corresponding diagram in Fig. 12. The dimer in the intermediate states possessing no stable nucleus is now considered to be metastable, and we introduce the parameter $$\beta \mathrm{exp}(\mathrm{\Delta }E_j^{\mathrm{dis}}/k_\mathrm{B}T)$$ (64) as “dissociation probability” (analogous to the step edge crossing probability $`\alpha =\mathrm{exp}(\mathrm{\Delta }E_\mathrm{S}/k_\mathrm{B}T)`$). For $`\beta =1`$ we recover the non-interacting particle model. From the outset it is clear that second layer nucleation will proceed faster for smaller $`\beta `$, since the state probabilities $`p_{3,2}(t)`$ and $`p_{4,2}(t)`$ in Fig. 12 and accordingly the currents $`j_3(t)`$ and $`j_4(t)`$ defined in eq. 62 will become strongly enhanced. Figure 18 shows $`R_c`$ as a function of $`\alpha `$ for fixed $`\mathrm{\Gamma }=10^8`$ and various $`\beta `$ obtained from Monte Carlo simulations (open symbols). As expected, the critical radius decreases with decreasing $`\beta `$. In fact, for $`\beta =10^6`$ one can regard the dimer as effectively being stable on the relevant time scale $`t_c`$, so that the changes with $`\beta `$ correspond to a continuous transition from $`i=2`$ ($`\beta =1`$) to $`i=1`$ ($`\beta =10^6`$). The comparison with the numerical solution of eqs. 59 (solid lines) yields very good agreement. To achieve this agreement, we used the same set of prefactors $`\kappa _{\mathrm{a},\mathrm{j}}^{(n)}`$, $`\kappa _{\mathrm{d},2}`$ as in Fig. 14. This highlights the power of the rate equation approach to treat second layer nucleation in general situations. ## VI Summary and Discussion In summary, we have presented a detailed theoretical investigation of the nucleation on top of islands in epitaxial growth. In the non-interacting particle model, where the lifetimes of unstable clusters can be neglected, it was possible to tackle the problem within a simplified stochastic description based on scaling arguments. An important result for the non-interacting case is that the nucleation for critical nuclei of size $`i2`$ is dominated by fluctuations, while for larger critical nuclei it can be treated in a mean-field type manner. The second layer nucleation rate for both cases was derived in compact form (see eqs. (42,40)). When metastable clusters can form with appreciable lifetimes, the simplified description can in principle be extended (see Sec. IV E), but becomes of limited value due to the fact that many elementary processes get mutually coupled both sequentially and in parallel. In such situations it is better to employ the more general framework outlined in Sec. V that is based on our derivation for the transition rates of the elementary processes. Results obtained from both theoretical approaches were shown to agree with Monte Carlo data. Throughout the paper, we have used the generic growth law (18) for the mean island radius, but it is straightforward to treat other growth laws also (as e.g. an exponential behavior), which may be realized by special preparation techniques. Neither the general expressions (40,42) for the second layer nucleation rates in simple situations nor the master equation (59) depend on the specific form of the growth law (the expressions for $`\overline{n}`$, $`p_n`$, etc. in the quasi-stationary case however get modified, see the discussion in Sec. IV A). Moreover, it is straightforward to rewrite all formulae for the case of heteroepitaxy by replacing the jump rate $`D/a^2`$ of adatoms on top of the islands by a modified jump rate $`D^{}/a^2`$. The theoretical understanding of second layer nucleation is not only of basic importance but has numerous applications. One of these is the determination of the effective step edge barrier $`\mathrm{\Delta }E_\mathrm{S}`$ for systems, where the more direct and simpler method via the measurement of adatom lifetimes by field ion microscopy cannot be applied. As pointed out in ref. , the breakdown of the TDT approach in the fluctuation-dominated situation calls for a reexamination of some experimental data for estimating $`\mathrm{\Delta }E_\mathrm{S}`$. In fact, such reexamination has been carried out recently by Krug et al. with notable results: By reanalyzing the fraction of covered islands $`f(t)`$ measured for Ag/Ag(111) they corrected the previously reported estimate $`\mathrm{\Delta }E_s0.12\mathrm{eV}`$ to $`\mathrm{\Delta }E_s0.32\mathrm{eV}`$ (they also reported another estimate yielding $`\mathrm{\Delta }E_s0.20\mathrm{eV}`$ based on a modified data analysis, see the comment in ref. ). Krug et al. moreover studied the influence of step decoration by CO molecules on $`R_c`$ (and hence $`\mathrm{\Delta }E_\mathrm{S}`$) for Pt/Pt(111). They found a strong increase of $`\mathrm{\Delta }E_\mathrm{S}`$ with CO partial pressures, when analyzing the data corresponding to regime II (for $`i=1`$) of the fluctuation-dominated situation. Hence contamination by CO is expected to favor multilayer growth. On the other hand, surfactants may promote smooth layer-by-layer growth. For example, the presence of only small amounts of Sb for growth of Ag on Ag(111) were shown to convert rough multilayer to layer-by-layer growth. It was suggested that Sb reduces $`\mathrm{\Delta }E_\mathrm{S}`$, but, since it was observed that Sb increases the island density in the first layer, it is also possible that the induced layer-by-layer growth results from a decrease of the mean island distance. Even in the absence of surfactants, a change of the effective step edge barrier may go along with a shape transition of the islands with varying temperature (see refs. and the comment in ref. ), and this can induce changes in the film morphology as well. With respect to the transition from the fluctuation-dominated to the mean-field type situation with varying $`i`$ predicted in this work, it would also be interesting to conduct proper experiments for metal epitaxy on (100) surfaces, where a change from $`i=1`$ to $`i=3`$ is often observed with increasing temperature. A further application pertaining to the design of self-organized nanostructures is the possibility to create pyramidal mounds on a substrate, which are called “wedding cakes”. As suggested by Michely et al., the size $`L_{\mathrm{top}}`$ of the top terrace of the pyramid should be roughly given by $`\mathrm{\Omega }(L_{\mathrm{top}})F`$, where $`\mathrm{\Omega }`$ is the second layer nucleation rate. Recently, an expression for the distribution of $`L_{\mathrm{top}}`$ has been suggested within a self-consistent analysis of a model for the dynamics of the top terrace. In recent developments of nanostructure formation also larger clusters of atoms are considered as basic building blocks in epitaxial growth. The underlying processes seem to be very similar to the case of deposition of single atoms or simple molecules (for a recent review, see Ref. ), so that it could well be that also for cluster deposition an effective step edge barrier has a decisive influence on the film topography. In light of the basic importance and the manifold applications, there is certainly need for further improvement of our understanding of second layer nucleation. Topics worthy of further study are in particular the influence of strain effects and longer-range interactions between the adatoms. The latter may be attributed to direct forces (e.g. induced dipole-dipole forces in the case of magnetic adsorbates), or they can be mediated by perturbations of the electron structure of the substrate. By extending the approach presented in this work, these issues may be tackled in the near future. ###### Acknowledgements. We thank R. J. Behm, H. Brune, W. Dieterich, and H. J. Ernst for very interesting discussions. P.M. thanks the Deutsche Forschungsgemeinschaft for financial support (SFB 513, Ma 1636/2). ## A We want to calculate the characteristic rate $`\omega _n(R)`$ for an encounter of $`i+1`$ atoms, if initially $`ni+1`$ atoms are randomly placed on top of an island with radius $`R`$ and infinite step edge barrier ($`\alpha =0`$). For this purpose let us consider the encounter as a sequential process as depicted in Fig. 19 (for $`i=3`$ and $`n=5`$): First a dimer forms, then one of the remaining atoms attaches to the dimer and a trimer is created, and so on until a stable cluster composed of $`i+1`$ atoms has been formed. Denoting the rate for the formation of the dimer by $`w_1`$, and the rate for the attachment of an atom to an already existing cluster composed of $`k`$ atoms (“$`k`$-cluster”) by $`w_k`$, we may write $$w_1=\frac{n(n1)}{2}\frac{2D}{a^2}\frac{b_1a^2}{\pi R^2},$$ (A2) $$w_k=(nk)\frac{D}{a^2}\frac{b_ka^2}{\pi R^2},2ki.$$ (A3) The factors $`b_k`$ can be viewed as the effective number of perimeter sites of a $`k`$-cluster. Similarly, we may writefor the rate of dissociation $`v_k`$ of a single atom from a $`k`$-cluster (in the case of negligible binding energies of unstable clusters) $$v_k=d_k\frac{D}{a^2},k2,$$ (A4) where again $`d_k`$ has the meaning of an effective number of perimeter sites. (In principle one may also take into account the possibility that a subcluster composed of more than one atom can dissociate from an unstable cluster and other states with various intermediate unstable clusters of size $`2ki`$.) The idea now is to renormalize the process depicted in Fig.19 by replacing it by an effective transition rate $`w_{\mathrm{eff}}`$ between the initial state composed of $`n`$ isolated atoms and the final state containing the stable cluster. Clearly, such a replacement is only approximately valid. After the replacement, the encounter rate $`\omega _n(R)`$ in eq. (32) can be identified with $`w_{\mathrm{eff}}`$. In order to derive $`w_{\mathrm{eff}}`$, we consider a stationary situation, where the probability $`p_1`$ of the initial state is kept fixed and a constant current $`J`$ flows between neighboring states containing a $`k`$\- and $`k+1`$-cluster. We thus write $$J=w_kp_kv_{k+1}p_{k+1},1ki1,$$ (A5) where $`p_k`$ denotes the probability of the state containing a $`k`$-cluster. The set of equations (A5) can be readily solved for $`p_i`$ yielding $$J=w_ip_i=w_1p_1\underset{k=2}{\overset{i}{}}\frac{w_k}{v_k}J\underset{k=2}{\overset{i}{}}\underset{j=k}{\overset{i}{}}\frac{w_j}{v_j}.$$ (A6) On the other hand we have $$J=w_{\mathrm{eff}}p_1.$$ (A7) Eliminating $`J`$ from eqs. (A6,A7), we obtain $$w_{\mathrm{eff}}=\frac{w_1_{k=2}^iw_k/v_k}{1+_{k=2}^i_{j=k}^iw_j/v_j}.$$ (A8) For large radii $`Ra`$, it holds $`w_j/v_j1`$ so that we can neglect the sum over $`k`$ in the denominator on the right hand side of (A8). Hence we find $$\omega _n(R)w_{\mathrm{eff}}\kappa _\mathrm{e}\left[\underset{k=0}{\overset{i}{}}(nk)\right]\frac{D}{a^2}\left(\frac{a^2}{\pi R^2}\right)^i,$$ (A9) where $`\kappa _\mathrm{e}=b_1_{k=2}^ib_k/d_k`$. ## B The solution of the diffusion problem $$\frac{\rho }{t}=D\mathrm{\Delta }\rho ,$$ (B2) $$\frac{\rho }{r}\|_{r=0}=0,[\frac{\rho }{r}+\frac{\alpha }{a}\rho ]_{r=R}=0$$ (B3) with the initial condition $`\rho (𝐫,t=0)=1/(\pi R^2)`$ has been derived by Harris: $$\rho (𝐫,t)=\underset{k=1}{\overset{\mathrm{}}{}}\frac{c_k\lambda _k^2}{2\pi R^2\frac{\alpha R}{a}}\frac{J_0(\frac{\lambda _kr}{R})}{J_0(\lambda _k)}\mathrm{exp}\left(\lambda _k^2\frac{D}{R^2}t\right).$$ (B4) Here $`J_\nu (.)`$ is the Bessel function of $`\nu `$th order, $`c_k4(\alpha R/a)^2/(\lambda _k^2[\lambda _k^2+(\alpha R/a)^2])`$ and $`\lambda _k`$ is the $`k`$th root ($`\lambda _1<\lambda _2<\mathrm{}`$) of $$\left(\frac{\alpha R}{a}\right)J_0(\lambda )=\lambda J_1(\lambda ).$$ (B5) The solution (B4,B5) describes the probability density for a single diffusing atom that at time $`t=0`$ is randomly deposited on top of a circular island with a partially reflecting boundary. The probability that the atom has not escaped from the island up to time $`\tau `$ is $`\mathrm{\Psi }(\tau )=2\pi _0^R𝑑rr\rho (𝐫,t)`$, which yields $$\mathrm{\Psi }(\tau )=\underset{k=1}{\overset{\mathrm{}}{}}c_k\mathrm{exp}\left(\lambda _k^2\frac{D}{R^2}\tau \right)$$ (B6) Note that, since $`\mathrm{\Psi }(0)=1`$, it must hold $`_{k=1}^{\mathrm{}}c_k=1`$. The probability that none of $`n`$ independent atoms has escaped from the island up to time $`\tau `$ is $`\mathrm{\Psi }(\tau )^n`$. Accordingly, the probability $`\varphi (\tau )d\tau `$ that the first atom leaves the island in the time interval $`[\tau ,\tau +d\tau ]`$ is $`\varphi (\tau )`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Psi }(\tau )^n}{d\tau }}=n\mathrm{\Psi }(\tau )^{n1}{\displaystyle \frac{d\mathrm{\Psi }(\tau )}{d\tau }}`$ (B7) $`=`$ $`n{\displaystyle \frac{D}{R^2}}{\displaystyle \underset{j_1,\mathrm{},j_n=1}{\overset{\mathrm{}}{}}}c_{j_1}\mathrm{}c_{j_n}\times `$ (B9) $`\times \lambda _{j_1}^2\mathrm{exp}\left[(\lambda _{j_1}^2+\mathrm{}+\lambda _{j_n}^2){\displaystyle \frac{D}{R^2}}\tau \right],`$ from which for the average time $`\tau _n(R)_0^{\mathrm{}}𝑑\tau \varphi (\tau )\tau `$ follows: $$\tau _n(R)=n\frac{R^2}{D}\underset{j_1,\mathrm{},j_n=1}{\overset{\mathrm{}}{}}c_{j_1}\mathrm{}c_{j_n}\frac{\lambda _{j_1}^2}{(\lambda _{j_1}^2+\mathrm{}+\lambda _{j_n}^2)^2}.$$ (B10) It is easy to show that $`j_{1,k}<\lambda _k<j_{0,k}`$, where $`j_{\nu ,k}`$ is the $`k`$th zero of $`J_\nu (.)`$. Since $`j_{\nu ,k}(k+\nu /21/4)\pi `$ for $`k\nu `$, the terms in the series of (B10) rapidly decrease with increasing $`j_k`$, $`k=1,\mathrm{},n`$ (note that $`c_j`$ depends on $`\lambda _j`$). The leading term can be obtained by setting $`c_j=\delta _{j,1}`$ in eq. (B6), which amounts to a Poisson approximation of the escape process, $`\mathrm{\Psi }(\tau )\mathrm{exp}(\lambda _1D\tau /R^2)`$. Within this approximation we obtain $$\tau _n(R)=\frac{1}{n}\frac{R^2}{D}\frac{1}{\lambda _1^2},$$ (B11) where $`\lambda _1`$ follows from eq. (B5). In the limit of small $`\alpha R/a1`$ one finds $`\lambda _1^22\alpha R/a`$, while in the limit of large $`\alpha R/a1`$, $`\lambda _1^2j_{0,1}^2`$. Combining these two limits yields the interpolation formula $$\tau _n(R)\frac{1}{n}\frac{R^2}{D}\left(\kappa _1\frac{a}{\alpha R}+\kappa _2\right)$$ (B12) with $`\kappa _1=1/2`$ and $`\kappa _2=1/j_{0,1}^20.173`$. Knowing $`\tau _n(R)`$ we can set up the master equation for the probabilities $`p_n(t)`$ to find exactly $`n`$ atoms on top of the island at time $`t`$ in the presence of an incoming flux $`F`$, see eq. (35). Introducing the generating function $`Q(z,t)_{n=0}^{\mathrm{}}p_n(t)z^n`$ we obtain from eq. (35) $$\frac{Q}{t}=(z1)\left[\pi FR^2Q\frac{1}{\tau _1(R)}\frac{Q}{z}\right],$$ (B13) where $`R=R(t)`$ from eq. (18). Transforming variables from $`t`$ to $`R`$ and defining $`\stackrel{~}{Q}(z,R)`$ by $`\stackrel{~}{Q}(z,R=R(t))Q(z,t)`$, eq. (B13) gives ($`a1`$ here) $$\frac{\stackrel{~}{Q}}{R}+\frac{\phi \stackrel{~}{\alpha }}{1+\stackrel{~}{\alpha }R}(z1)\frac{\stackrel{~}{Q}}{z}=\frac{2\pi }{A^2\mathrm{\Gamma }^{i/(i+2)}}R^3(z1)\stackrel{~}{Q},$$ (B14) where $`\stackrel{~}{\alpha }=\kappa _2\alpha /\kappa _1`$ and $`\phi =2A^2\mathrm{\Gamma }^{2/(i+2)}/\kappa _2`$ Equation (B14) is a semi-linear partial differential equation of first order that can be solved by the method of characteristics. For the initial condition $`\stackrel{~}{Q}(z,R=0)=1`$ we obtain $$\stackrel{~}{Q}(z,R)=\mathrm{exp}\left[(1z)\overline{n}(R)\right],$$ (B15) which for $`p_n(R)=p_n(R(t))=[_z^n\stackrel{~}{Q}(z,R)/n!]_{z=0}`$ yields the Poisson distribution (36) with $`\overline{n}(R)`$ $`=`$ $`{\displaystyle \frac{2\pi }{A^2\mathrm{\Gamma }^{\frac{i}{i+2}}}}(1+\stackrel{~}{\alpha }R)^\phi {\displaystyle _0^R}𝑑xx^3(1+\stackrel{~}{\alpha }x)^\phi `$ (B16) $`=`$ $`{\displaystyle \frac{2\pi }{A^2\mathrm{\Gamma }^{\frac{i}{i+2}}\stackrel{~}{\alpha }^4}}(1+\stackrel{~}{\alpha }R)^\phi `$ (B19) $`\times [{\displaystyle \frac{(1+\stackrel{~}{\alpha }R)^{\phi +4}1}{\phi +4}}3{\displaystyle \frac{(1+\stackrel{~}{\alpha }R)^{\phi +3}1}{\phi +3}}`$ $`+3{\displaystyle \frac{(1+\stackrel{~}{\alpha }R)^{\phi +2}1}{\phi +2}}{\displaystyle \frac{(1+\stackrel{~}{\alpha }R)^{\phi +1}1}{\phi +1}}].`$ In the quasi-stationary case ($`_tQ=0`$ in eq. (B13)) one obtains $$\overline{n}^{\mathrm{st}}(R)=\pi FR^2\tau _1(R)=\pi \frac{FR^4}{D}\left(\kappa _1\frac{a}{\alpha R}+\kappa _2\right),$$ (B20) to which (B19) simplifies for $`\phi \stackrel{~}{\alpha }R1`$ (see the discussion right after eq. (37)). Alternatively, we can determine $`\overline{n}^{\mathrm{st}}(R)`$ by integrating $`\rho _1^{\mathrm{st}}(r)`$ over the island area, which yields $$\overline{n}^{\mathrm{st}}(R)=2\pi _0^R𝑑rr\rho _1^{\mathrm{st}}(r)=\pi \frac{FR^4}{D}\left(\frac{1}{2}\frac{a}{\alpha R}+\frac{1}{8}\right).$$ (B21) Hence we can improve the Poisson approximation by renormalizing the bare coefficient $`\kappa _2=1/j_{0,1}^20.173`$ to $`\kappa _2=1/8`$ (note that $`\kappa _1=1/2`$ does not change). Finally, in order to obtain the constants $`\kappa _1`$ and $`\kappa _2`$ in eq. (33), one has to take into account the “lattice corrections”. Let us denote by $`𝐥`$ the position of a lattice site and by $`𝜹_j`$ the nearest-neighbor bond vectors, i.e. for a triangular lattice $`𝜹_j=(\mathrm{cos}[2\pi j/6],\mathrm{sin}[2\pi j/6])a`$, $`j=0,\mathrm{},5`$. The master equation describing the diffusion of a single adatoms on the island reads $$\frac{w(𝐥,t)}{t}=\frac{D}{6a^2}\underset{j=0}{\overset{5}{}}\left[w(𝐥+𝜹_j,t)w(𝐥,t)\right],$$ (B22) where $`w(𝐥,t)`$ is the probability to find the atom at lattice site $`𝐥`$. Equation (B22) is valid as long as $`𝐥`$ is not a boundary site. In the continuum limit we can write $`_{j=0}^5[w(𝐥+𝜹_j,t)w(𝐥,t)]=(1/2)_{j=0}^5(𝜹_j\mathbf{})^2w(𝐥,t)+𝒪(a^4)`$, which yields eq. (B2) when $`D`$ is replaced by $`D_{\mathrm{cont}}=D/4`$ (see also the remark in ref. ). As a consequence, one has to substitute $`D`$ by $`D_{\mathrm{cont}}`$ in all continuum equations, in particular in eq. (B21), which means that in eq. (33) (referring to the lattice simulations) one should take $`\kappa _1=4(1/2)=2`$ and $`\kappa _2=4(1/8)=1/2`$. The value of $`\kappa _1`$ is still not correct, since we have not taken into account the lattice correction to the parameter $`\alpha `$. To derive this correction we consider a lattice site $`𝐥`$ at the boundary. For example, one may encounter the situation sketched in Fig. 20, where $`𝜹_0`$ and $`𝜹_1`$ lead to sites outside the island (where $`w(𝐥+𝜹_j,t)0`$) and the remaining nearest neighbor sites $`𝐥+𝜹_j`$, $`j=2,\mathrm{},5`$ are on the island. The equation corresponding to (B22) then reads $`{\displaystyle \frac{w(𝐥,t)}{t}}`$ $`=`$ $`{\displaystyle \frac{\alpha D}{6a^2}}{\displaystyle \underset{j=0}{\overset{1}{}}}\left[w(𝐥+𝜹_j,t)w(𝐥,t)\right]`$ (B24) $`+{\displaystyle \frac{D}{6a^2}}{\displaystyle \underset{j=2}{\overset{5}{}}}\left[w(𝐥+𝜹_j,t)w(𝐥,t)\right]`$ $`=`$ $`{\displaystyle \frac{2\alpha D}{6a^2}}w(𝐥,t)+{\displaystyle \frac{D}{6a^2}}{\displaystyle \underset{j=0}{\overset{5}{}}}\left[w(𝐥+𝜹_j,t)w(𝐥,t)\right]`$ (B26) $`{\displaystyle \frac{D}{6a^2}}{\displaystyle \underset{j=0}{\overset{1}{}}}\left[w(𝐥+𝜹_j,t)w(𝐥,t)\right].`$ In a discretization of the second boundary condition in eq. (B3) on a triangular lattice one has to eliminate the outer boundary points $`𝐥+𝜹_0`$ and $`𝐥+𝜹_1`$ via the discretized version of the “bulk equation” (B3). This amounts to a cancellation of the term on the left hand side and the second term on the right hand side of (B26) in the continuum limit, and the replacement $`_{j=0}^1[w(𝐥+𝜹_j,t)w(𝐥,t)]=_{j=0}^1(𝜹_j\mathbf{})w(𝐥,t)+𝒪(a^2)\sqrt{3}aw/r`$. Hence eq. (B26) corresponds to the second boundary condition in eq. (B3), when $`\alpha `$ is replaced by $`\alpha _{\mathrm{cont}}=2\alpha /\sqrt{3}`$ in eq. (B3). In general, $`k`$ nearest neighbor sites of a boundary site $`𝐥`$ can lie outside the island ($`k=1,\mathrm{},4`$). The weights how often such $`𝐥`$ occur and the way the normal direction is oriented with respect to the nearest neighbor bond vectors leading to the sites outside the island depend sensitively on the shape of the island edge. Hence, the factor $`2/\sqrt{3}`$ is only an estimate, which gives an impression on the influence of the lattice correction to the coefficient $`\kappa _1`$. Our comparison with the simulation results in Fig. 7 yields $`\alpha _{\mathrm{cont}}2\alpha `$, i.e. $`\kappa _14(1/2)^2=1`$. We note that in general lattice corrections have always to be included in a continuum description after the effective Schwoebel barrier (for the lattice) has been calculated from the microscopic barriers (see the comment in ref. ).
warning/0004/nucl-th0004055.html	ar5iv	text	# SNUTP-00-009nucl-th/0004055 ELECTROMAGNETIC PRODUCTION OF VECTOR MESONS AT LOW ENERGIES11footnote 1Talk at the NSTAR2000 Workshop, The Physics of Excited Nucleons, JLab, Newport News, Feb. 16–19, 2000 ## Acknowledgments We are grateful to V. Burkert, W.-C. Chang, F. J. Klein, N. I. Kochelev and T. Nakano for fruitful discussions and informations. Y.O. thanks the Theory Division of Argonne National Laboratory and the organizer of this workshop for the warm hospitality and financial support. This work was supported in part by KOSEF of Korea, NSC of Republic of China, Russian Foundation for Basic Research under Grant No 96-15-96426 and U.S. DOE Nuclear Physics Division Contract No. W-31-109-ENG-38. ## References
warning/0004/astro-ph0004207.html	ar5iv	text	# A robust method for fitting peculiar velocity field models ## 1 Introduction The study of the large-scale motions of galaxies in the Universe may provide valuable information concerning the dynamics of large-scale structures and the nature of the underlying dark matter. According to the gravitational instability scenario, the peculiar velocity field (i.e. the deviation from the smooth Hubble flow) may be used to infer the power spectrum of the mass fluctuations on intermediate scales and to constrain the cosmological density parameter $`\mathrm{\Omega }`$ (see for example Dekel 1994). Since the discovery of the Great Attractor (Lynden-Bell et al. 1988), the field has proven to be particularly active. Various observationnal programs have been completed, providing large and accurate datasets: e.g. the W91CL and W91PP samples (Willick 1990); MAT sample (Mathewson et al. 1992); HM sample (Han & Mould 1992); CF sample (Courteau et al. 1993); Abell BCG sample (Lauer & Postman 1994); SCI sample (Giovanelli et al. 1997a); KLUN sample (Theureau et al. 1997); nearby SNIa sample (Riess et al. 1997); MarkIII dataset (Willick et al. 1997b); SBF survey (Tonry et al. 1997); SFI sample (Giovanelli et al. 1998); SMAC sample (Hudson et al. 1999); EFAR and ENEAR samples (Colless et al. 1999 and Wegner et al. 1999); SCII sample (Dale et al. 1999); Shellflow survey (Courteau et al. 1999); LP10k survey (Willick 1999a). Another significant advance during the past decade has been the improved understanding of the statistical formalism underlying the use of galaxy distance indicators – and in particular the principles and practical methods of correcting for Malmquist bias. (See for example Hendry & Simmons 1990, Teerikorpi 1990, Bicknell 1992, Landy & Szalay 1992, Triay et al. 1994, Willick 1994, Hendry & Simmons 1994, Sandage 1994, Willick et al. 1995, Freudling et al. 1995, Willick et al. 1996, Rauzy & Triay 1996, Ekholm 1996, Triay et al. 1996, Rauzy 1997, Willick et al. 1997b, Giovanelli et al. 1997b, Theureau et al. 1998, Teerikorpi et al. 1999). Several methods for extracting dynamical and kinematical information from distance indicator datasets have been proposed, i.e. the POTENT method (Bertschinger & Dekel 1989, Dekel et al. 1990, Bertschinger et al. 1990, Dekel et al. 1999 and references therein) and its variants (Rauzy et al. 1993 and 1995, Newsam et al. 1995); the ITF method (Nusser & Davis 1995, Davis et al. 1996, Da Costa et al. 1998); the VELMOD method (Willick et al. 1997a, Willick & Strauss 1998). The comparison between the peculiar velocity or density fields inferred from distance indicator data with their corresponding fields derived from whole-sky redshift surveys has been one of the major issues addressed throughout the last decade. The question here is whether the spatial distribution of luminous matter e.g. the galaxies, traces the underlying mass fluctuations and if not, what are the properties of the “biasing” between the two fields? Up to now, the point has not received any consensual answer. Indeed, the application of POTENT to various distance indicator datasets (Sigad et al. 1998 and references therein) favours a value of $`\beta _I=\mathrm{\Omega }_0^{0.6}/b_I1`$ for the linear “biasing” density parameter, while the VELMOD and ITF fitting methods lead to a value of $`\beta _I0.5`$ (Davis et al. 1996, Willick et al. 1997a, Riess et al. 1997, Da Costa et al. 1998, Willick & Strauss 1998). The origins of this significant discrepancy have not yet been elucidated (see for example Strauss 1999, Willick 1999b). At least one of these methods is suffering from some systematic effects, i.e. some statistical biases plaguing the estimate of the parameter $`\beta `$ and not accounted for in the error analysis. This remark leads us to the object of the present paper. The POTENT, VELMOD and ITF methods all require, at some stage of the analysis, to assume some a priori working hypotheses concerning the characteristics of the distance indicator dataset. For Tully-Fisher data for example, it will be assumed that the Tully-Fisher law is well described by a linear relation. These methods apply moreover under the hypothesis that the observational selection effects obey particular conditions. How are the results affected if one or more of the working assumptions fails to be satisfied by the dataset is generally a question not addressed in the error analysis. The philosophy of the method we present herein is to reduce as far as possible the number of a priori hypotheses concerning the distance indicator sample. A direct consequence is that the range of application of the method will be considerably broadened. The statistical background of the method is presented section 2. Its potential is illustrated by testing the predicted IRAS peculiar velocity model on two samples. In section 3, we perform the analysis using the fluxes of the IRAS 1.2 Jy survey as the distance indicator. We have deliberately chosen this sample looking to demonstrate the wide range of application of the method. Where the POTENT, VELMOD or ITF methods would not have been successful in extracting kinematical information from this dataset, our method does. We also treat a more classical case, the Tully-Fisher MarkIII MAT sample, in section 4. Finally, in section 5 we summarise our conclusions. ## 2 The Method ### 2.1 The statistical model It is assumed herein that the distribution function of the absolute magnitudes $`M`$ of the population, i.e. the luminosity function $`f(M)`$, does not depend on the spatial position $`𝐫=(r,l,b)`$ of the galaxies. The probability density describing the sample splits in this case as $$dPdP_𝐫\times dP_M=\rho (r,l,b)r^2\mathrm{cos}bdldbdr\times f(M)dM$$ (1) where $`\rho (r,l,b)`$ is the spatial distribution function of the sources. The application of the method will be restricted to samples strictly complete up to a given magnitude limit $`m_{\mathrm{lim}}`$, or in other words where the selection function in apparent magnitude is well described by a sharp cut-off, i.e. $`\psi (m)=\theta (m_{\mathrm{lim}}m)`$ with $`\theta (x)`$ the Heaviside function. Accounting for selection effects, the probability density of the sample may be rewritten as $$dP=\frac{1}{A}h(\mu ,l,b)\mathrm{cos}bdldbd\mu f(M)dM\theta (m_{\mathrm{lim}}m)$$ (2) where $`\mu =mM=5\mathrm{log}_{10}r+25`$ is the distance modulus, $`h(\mu ,l,b)`$ the line-of-sight distribution of $`\mu `$ and $`A`$ is the normalisation factor warranting $`𝑑P=1`$. For convenience in notation the angular dependence in $`l`$ and $`b`$ will be hereafter implicit. Observational selection effects in apparent magnitude then introduce a correlation between $`M`$ and $`\mu `$. The milestone of the method consists in defining the random variable $`\zeta `$ as follows $$\zeta =\frac{F(M)}{F(M_{\mathrm{lim}})}$$ (3) where $`F(M)=_{\mathrm{}}^Mf(x)𝑑x`$ stands for the cumulative distribution function in $`M`$ and $`M_{\mathrm{lim}}M_{\mathrm{lim}}(\mu )`$ is the maximum absolute magnitude for which a galaxy at distance $`\mu `$ would be visible in the sample (e.g. $`M_{\mathrm{lim}}(\mu )=m_{\mathrm{lim}}\mu `$ if the k-correction is neglected). The volume element may be rewritten as $$d\mu d\zeta =\frac{f(M)}{F(M_{\mathrm{lim}}(\mu ))}d\mu dM$$ (4) and by definition the random variable $`\zeta `$ for a sampled galaxy belongs to the interval $`[0,1]`$. The probability density of Eq. (2) reduces thus to $$dP=\frac{1}{A}h(\mu )F(M_{\mathrm{lim}}(\mu ))d\mu \times \theta (\zeta )\theta (1\zeta )d\zeta $$ (5) with $`A=h(\mu )F(M_{\mathrm{lim}}(\mu ))𝑑\mu `$. Note that the probability density $`dP_\mu =\frac{1}{A}h(\mu )F(M_{\mathrm{lim}}(\mu ))d\mu `$ describes the observed spatial distribution function of the sources. It follows from Eq. (5) that: * P1: $`\zeta `$ is uniformly distributed between $`0`$ and $`1`$. * P2: $`\zeta `$ and $`\mu `$ are statistically independent, i.e. the distribution of $`\zeta `$ does not depend on the spatial position of the galaxies. Property P1 may be used to construct a test for assessing the completeness of the sample in apparent magnitude. The details of this test are presented in a separate paper (Rauzy, in preparation), although we apply it to the Mark III MAT sample later in this paper. The new method we propose hereafter for fitting peculiar velocity field models is based on property P2. ### 2.2 Estimate of the random variable $`\zeta `$ The random variable $`\zeta `$ can be estimated without any prior knowledge of the cumulative luminosity function $`F(M)`$. To each data point with coordinates $`(M_i,\mu _i)`$ is associated the region $`S_i=S_1S_2`$ defined as * $`S_1=\{(M,\mu )\mathrm{such}\mathrm{that}MM_i\mathrm{and}\mu \mu _i\}`$ * $`S_2=\{(M,\mu )\mathrm{such}\mathrm{that}M_i<MM_{\mathrm{lim}}^i\mathrm{and}\mu \mu _i\}`$ The random variables $`M`$ and $`\mu `$ are independent in each subsample $`S_i`$ since by construction the cut-off in apparent magnitude is superseded by the constraints $`MM_{\mathrm{lim}}^i`$ and $`\mu \mu _i`$ (see figure 1). This implies that the number of points $`r_i`$ belonging to $`S_1`$ is proportional to $`_{\mathrm{}}^{M_i}f(M)𝑑M=F(M_i)`$, the number of points $`n_i`$ in $`S_i=S_1S_2`$ is proportional to $`F(M_{\mathrm{lim}}^i)`$ and that the quantity $$\widehat{\zeta }_i=\frac{r_i}{n_i+1}$$ (6) is an unbiased estimate of the random variable $`\zeta `$. Equivalently the estimator $`\widehat{\zeta }_i`$ may be defined as the normalised rank of the point $`M_i`$ when the $`M`$’s are sorted by increasing order within the subsample $`S_i`$ (see Efron & Petrosian 1992). ### 2.3 Radial peculiar velocity field models Let us first assume that the true radial peculiar velocity field $`v(𝐫)`$ can be described by a one-parameter velocity model $`v_\beta (𝐫)`$, i.e. there exists a solution $`\beta ^{}`$ satisfying $`v_\beta ^{}(𝐫)v(𝐫)`$. For a given value of the parameter $`\beta `$, the model dependent variables $`\mu _\beta `$ and $`M_\beta `$ can be computed (modulo the value of the Hubble constant $`H_0`$) from the observed redshift $`z`$ and apparent magnitude $`m`$ following $$\mu _\beta =5\mathrm{log}_{10}\frac{cz}{H_0}+25u_\beta ;M_\beta =m\mu _\beta $$ (7) where the quantity $`u_\beta `$ is defined as $$u_\beta =5\mathrm{log}_{10}\left(1\frac{v_\beta }{cz}\right)$$ (8) The quantities $`\mu _\beta `$ and $`M_\beta `$ are related to the true absolute magnitude $`M`$ and distance modulus $`\mu `$ via $$\mu _\beta =\mu +u_\beta ^{}u_\beta ;M_\beta =Mu_\beta ^{}+u_\beta $$ (9) Computing $`\zeta _\beta `$ from $`\mu _\beta `$ and $`M_\beta `$ as proposed in Eq. (3) gives for the probability density of Eq. (5), $$dP=\frac{1}{A}h(\mu )F(M_{\mathrm{lim}}(\mu _\beta ))d\mu C_\beta \theta (\zeta _\beta )\theta (1\zeta _\beta )d\zeta _\beta $$ (10) where $`C_\beta `$ takes the following form when $`(u_\beta ^{}u_\beta )1`$ (or equivalently $`(v_\beta ^{}v_\beta )cz`$), $$C_\beta =\frac{f(M)}{f(M_\beta )}1+(u_\beta u_\beta ^{})\left(\mathrm{ln}f\right)^{}(M_\beta )$$ (11) Because the absolute magnitude $`M_\beta `$, and hence the quantity $`\left(\mathrm{ln}f\right)^{}(M_\beta )`$, is correlated with the random variable $`\zeta _\beta `$, $`C_\beta `$ acts as a correlation coefficient between $`\zeta _\beta `$ and the proposed velocity field model $`u_\beta `$ when $`\beta \beta ^{}`$. On the other hand if $`\beta =\beta ^{}`$ these quantities are statistically independent, since in this case, according to property P2, $`\zeta _\beta \zeta `$ does not depend on the spatial position of galaxies and therefore on any function $`u_\beta (𝐫)`$. It thus turns out that any statistical test of independence between $`\zeta _\beta `$ and $`u_\beta `$ provides us with an unbiased estimate of the value of $`\beta ^{}`$. In particular the coefficient of correlation $`\rho (\zeta _\beta ,u_\beta )`$ has to vanish when $`\beta =\beta ^{}`$, i.e. $$\beta =\beta ^{}\rho (\zeta _\beta ,u_\beta )=0$$ (12) As revealed by Eq. (11), the accuracy of this estimator is related to the amplitude of the correlation between $`\left(\mathrm{ln}f\right)^{}(M_\beta )`$ and $`\zeta _\beta `$. The steeper the function $`\left(\mathrm{ln}f\right)^{}`$, or in other words the smaller the dispersion of the luminosity function $`f(M)`$, the more accurate is the estimate of the velocity parameter $`\beta `$, as expected. In practice, this accuracy can be obtained through numerical simulations by analysing the influence of sampling fluctuations on the coefficient of correlation $`\rho (\zeta _\beta ,u_\beta )`$. The presence of a small-scale velocity dispersion (say of amplitude $`\sigma _v`$), not described by the velocity model $`v_\beta `$, introduces according to Eq. (9) a correlation between the derived quantities $`\mu _\beta `$ and $`M_\beta `$, and consequently between the variables $`\mu _\beta `$ and $`\zeta _\beta `$. Anyway since it is the correlation between the velocity model $`u_\beta `$ and $`\zeta _\beta `$ which is considered herein, and because the random velocity noise is not supposed to be correlated with $`u_\beta `$, the presence of a small-scale velocity dispersion is not expected to drastically bias the estimator proposed Eq. (12), at least as long as the variations of the quantity $`u_\beta (𝐫)`$ are smooth at the scale $`\sigma _v`$. Thanks to the introduction of the random variable $`\zeta `$, an unbiased estimate of the parameter $`\beta `$ has indeed been obtained using a null-correlation technique. Null-correlation approaches are characterised, in general, by their robustness – i.e. some of the functions entering the statistical model are not required to be fully specified (see for example Fliche & Souriau 1979, Bigot et al. 1991, Triay et al. 1994, Rauzy 1997). Unlike the maximum likelihood methods, e.g. the method proposed by Choloniewski (1995) and the VELMOD method of Willick et al. (1997a), no a priori assumptions have been made here concerning the specific shape of the luminosity function and the spatial distribution of the sources. In particular homogeneous as well as inhomogeneous Malmquist biases are automatically accounted for. Note also that selection effects in distance or redshift are allowed since Eq. (10) accepts any extra terms of the form $`\psi (\mu ,u_\beta )`$. #### 2.3.1 Orthonormal decomposition of the velocity field It is worthwhile to mention that, with regard to Eq. (11), the orthonormal decomposition of the velocity field, proposed in the ITF method (Nusser & Davis 1995), may be also applied herein. To see this we proceed as follows. It is assumed hereafter that there exists a N-dimensional vector $`\beta ^{}=(\beta _1^{},\beta _2^{},\mathrm{},\beta _N^{})`$ such that the quantity $`u(𝐫)`$ can be decomposed as $$u(𝐫)u_\beta ^{}(𝐫)=\underset{j=1}{\overset{N}{}}\beta _j^{}u_j(𝐫)$$ (13) where $`u_1(𝐫)`$, $`u_2(𝐫)`$, …, $`u_N(𝐫)`$ is a set of functions verifying the following orthonormality condition, $$\mathrm{Cov}(u_i,u_j)=\delta _{ij}$$ (14) with $`\delta _{ij}`$ the Kronecker symbol and $`\mathrm{Cov}(u_i,u_j)`$ the covariance of $`u_i`$ and $`u_j`$ on the sample. Note that within the approximation $`\beta _j^{}u_j(𝐫)1`$ (or equivalently $`\beta _j^{}v_j(𝐫)cz`$), Eq. (8) implies that $$v(𝐫)\underset{j=1}{\overset{N}{}}\beta _j^{}v_j(𝐫)$$ (15) The coefficient $`C_\beta `$ introduced in Eq. (11) rewrites, as long as $`(\beta _j\beta _j^{})u_j(𝐫)1`$ is satisfied, as $$C_\beta 1+\left(\mathrm{ln}f\right)^{}(M_\beta )\times \underset{j=1}{\overset{N}{}}(\beta _j\beta _j^{})u_j(𝐫)$$ (16) which implies that for each function $`u_i(𝐫)`$, $$\rho (u_i,\zeta _\beta )\underset{j=1}{\overset{N}{}}(\beta _j\beta _j^{})\mathrm{Cov}(u_i,u_j)$$ (17) and thus, because of the orthonormality condition of Eq. (14), that $$\beta _i=\beta _i^{}\rho (\zeta _\beta ,u_i)=0$$ (18) which provides us with the statistically independent estimates of the $`N`$ parameters $`\beta _1^{},\beta _2^{},\mathrm{},\beta _N^{}`$. The procedure for constructing the orthonormal family $`u_1(𝐫)`$, $`u_2(𝐫)`$, …, $`u_N(𝐫)`$ from an arbitrary set of $`N`$ independent functions is described in Nusser & Davis (1995). ## 3 Application to the IRAS sample The method described above is herein applied to the 60 $`\mu `$m IRAS 1.2 Jy sample (Fisher et al. 1995), a magnitude-redshift catalogue containing $`5321`$ galaxies and complete up to a flux $`F_{60}=1.2`$ Jy. Distance modulus and absolute magnitude are computed using $`H_0=100`$ km s<sup>-1</sup> Mpc<sup>-1</sup> for the value of the Hubble constant (the cut-off in apparent magnitude is expressed as $`m_{\mathrm{lim}}=14.3187`$ with this notation). The magnitudes have been k-corrected assuming a spectral slope $`\alpha =2`$, implying that the maximum absolute magnitude introduced in Eq. (3) reads as $`M_{\mathrm{lim}}(\mu )=m_{\mathrm{lim}}\mathrm{k}_{\mathrm{cor}}(\mu )\mu `$. The distribution of the sources in the $`M`$-$`\mu `$ plane is shown figure 1. The peculiar velocity field model tested is the one-parameter predicted IRAS velocity field characterised by the parameter $`\beta =\mathrm{\Omega }_0^{0.6}/b_I`$ (Strauss et al. 1992). ### 3.1 Using apparent magnitude as the distance indicator We have first applied the method using the apparent magnitude of IRAS galaxies as the distance indicator (restricting the analysis to the 4115 galaxies within the redshift range 1000-12000 km s<sup>-1</sup>). The Cumulative Luminosity Function (CLF) $`F(M)`$ is presented in figure 2. Here the CLF has been reconstructed using the $`C^{}`$ method of Lynden-Bell (1971) for two different velocity models, $`\beta =0`$ corresponding to the case where peculiar velocities are neglected and $`\beta =1`$. Note that the presence of peculiar velocities does not affect drastically the CLF reconstruction. The luminosity function of the IRAS galaxies does not exhibit any turnover towards the faint-end tail, at least within the observed range of magnitudes. This means that, because of the large dispersion of such a distance indicator and even if the number of galaxies is large, one cannot expect very strong constraints on the velocity model tested. The correlation between the random variable $`\zeta _\beta `$ and the velocity modulus $`u_\beta `$ for $`\beta =0.6`$ is shown in figure 3. Variations of the coefficient of correlation as a function of the parameter $`\beta `$ are given in figure 4. This curve is a monotonic function, as expected. The preferred value of $`\beta `$ is the one corresponding to $`\rho (\zeta _\beta ,u_\beta )=0`$ (here $`\beta =0.10.15`$). Monte Carlo simulations have been used to calculate the discrepancy between $`\rho (\zeta _\beta ,u_\beta )`$ and zero due to sampling fluctuations. Each simulation is a sample containing $`N_{\mathrm{gal}}=4115`$ galaxies with the same $`\mu _\beta `$ and $`u_\beta `$ as observed and for which the random variable $`\zeta _\beta `$ is computed following Eq. (6) where the rank $`r_i`$ is randomly generated according to a discrete uniform distribution between $`1`$ and $`n_i`$. The cumulative distribution function of $`\rho (\zeta _\beta ,u_\beta )`$ obtained from a large number of simulations, under the null hypothesis that the true value of $`\rho (\zeta _\beta ,u_\beta )=0`$, is shown in figure 5. Note that in practice this distribution does not depend on the amplitude of the quantity $`u_\beta `$ since the coefficient of correlation is a scale-free estimator of the correlation between two random variables. The cumulative distribution function allows us to evaluate the probability that the observed $`\rho `$ is by chance greater than a given value, due to sampling fluctuations. A one-sided rejection test for the $`\beta `$ parameter can thus be constructed. For example, the models with $`\beta 0.7`$ can be rejected with a confidence level of 95% and $`\beta 1.1`$ with a confidence level of 99%. So our method applied to the IRAS sample, using the apparent magnitude of the galaxies as a distance indicator, permits already at this stage to reject high values for the parameter $`\beta `$. ### 3.2 Introduction of a second parameter In a second step, the analysis is refined by taking into account the observed correlation between the absolute magnitude $`M`$ and some “colour index” defined as $`p=2.5\mathrm{log}_{10}(F_{100}/F_{60})`$ (with $`F_{100}`$ the flux at 100 $`\mu `$m). The data have been grouped in $`8`$ classes by interval of $`p`$ (see figure 6). Because of the (weak) correlation between $`p`$ and $`M`$, the spread of the luminosity function for each of these classes taken individually is expected to be slightly smaller than the spread of the global luminosity function, and thus the accuracy of the distance indicator somewhat improved. Figure 7 illustrates such a trend. For each individual class, the random variable $`\zeta _\beta `$ is computed according to Eq. (6). The correlation between $`\zeta _\beta `$ and the velocity modulus $`u_\beta `$ is after that evaluated for the whole sample. The influence of sampling fluctuations is estimated from Monte Carlo simulations, as described in the previous section. The results are presented in figure 8 in terms of the confidence level of rejection for the parameter $`\beta `$. More explicitly, the quantity plotted in ordinate is $`12\mathrm{Prob}(\rho \|\rho _{\mathrm{obs}}(\beta )\|)`$ where the probability that the coefficient of correlation $`\rho `$ is less than $`\|\rho _{\mathrm{obs}}(\beta )\|`$ due to sampling fluctuations is given by the cumulative distribution function of $`\rho `$. The method was first applied to the galaxies in the redshift range 1000-12000 km s<sup>-1</sup>. It is found that $`\beta [0.35,0.25]`$ at $`1\sigma `$, and that models with $`\beta 0.5`$ can be rejected with a confidence level of 95%. This result is in disagreement with most of the analyses based on Tully-Fisher data e.g. VELMOD on MarkIII (Willick & Strauss 1998), ITF method on SFI (Da Costa et al. 1998), ROBUST method on MarkIII MAT sample (see next section), favouring a value of $`\beta 0.5`$. We interpret this discrepancy as follows. When fitting a velocity model to data, the natural weight assigned by the fitting procedure to each galaxy is roughly proportional to the inverse of its redshift, because the accuracy of the distance indicator decreases as $`1/z`$. The mean effective depth of the volume where the velocity model is compared to data has to be estimated using these weights. For our first sample with $`z[1000,12000]`$ km s<sup>-1</sup>, we find a mean effective depth of $`3800`$ km s<sup>-1</sup> (see figure 9). In order to mimic the effective volume sampled by Tully-Fisher data we have applied the method to a truncated version of the IRAS sample containing 1621 galaxies with $`z[500,5000]`$ and galactic latitude $`\|b\|>20`$ (the mean effective depth of this sub-sample is now 2200 km s<sup>-1</sup>, see figure 10). Figure 8 shows that the value of $`\beta `$ estimated from this truncated sample is fully consistent with the values obtained using Tully-Fisher data. An interpretation of these results could be that the predicted IRAS velocity field model, while successful in reproducing locally the cosmic flow, fails to describe the kinematics on larger scales. However, as pointed out by our anonymous referee, the results derived above apply only to the extent that the photometry of the IRAS sample does not suffer from systematic errors. It is worthwhile to stress again that the philosophy of the ROBUST method is to impose very few assumptions – only that the luminosity function is independent of position, that the sample is strictly complete in apparent magnitude and that redshift and apparent magnitude measurements are not affected by systematic biases. Thus, our analysis could be affected if the IRAS photometry were (mildly) non-uniform. To evaluate the amplitude of such effects on our results would essentially require to adopt a realistic model for these systematic photometry variations. Since our primary goal here is to present a new method, we feel that such an error analysis is unwarranted. It is also worth noting that systematic errors in the IRAS photometry – if indeed present – would also affect methods for reconstructing the IRAS predicted peculiar velocity model. In particular, they would generate some systematic variations in the spatial selection function entering the weighting scheme used in density and velocity reconstruction (Strauss et al. 1992, Branchini et al. 1999), leading to systematic discrepancies between the IRAS predicted velocity field and the true cosmic flow. Moreover, it is interesting to note that the ROBUST method could easily be extended to provide a non-parametric test of the uniformity of photometry in redshift surveys; we will consider such an application in future work. ## 4 Application to the MarkIII MAT sample We have shown in the previous section that our method permits to extract valuable information on the cosmic velocity field by using the fluxes of the IRAS galaxies as a distance indicator. The potential of our method is illustrated in this section by treating a more classical case: the Tully-Fisher MarkIII MAT sample (Willick et al. 1997b and references therein). In a first step we have selected from the 1355 galaxies of the MarkIII MAT catalog a subsample for which the selection effects in apparent magnitude are well described by a Heaviside cut-off, $`\theta (mm_{\mathrm{lim}})`$. Assuming that the galaxies are homogeneously distributed in space, the value of $`m_{\mathrm{lim}}`$ can be found by analysing the variations in the logarithm of the cumulative count as a function of the limit in apparent magnitude – see for example figure 4 in Rauzy (1997). In a separate paper (Rauzy, in preparation) we propose anyway a simple tool for assessing the completeness in apparent magnitude of redshift-magnitude catalogues which does not require any assumption concerning the spatial distribution of the sources. This method is closely based on the statistical test presented in Efron & Petrosian (1992). We have applied this test to the MarkIII MAT sample, finding that the completeness in I-band corrected apparent magnitude is satisfied up to $`m_{\mathrm{lim}}=11.25`$ magnitudes. Discarding, in addition, the galaxies with $`z500`$ km s<sup>-1</sup>, we are left with a subsample containing $`318`$ galaxies. Note that the selection effects in redshift which affect the MarkIII MAT sample (Willick et al. 1996) are not a source of problems for our method. The Tully-Fisher relation for the $`318`$ extracted galaxies is shown in figure 11. In a second step, we divided the subsample in $`6`$ classes $`C_i`$ according to the value of $`\eta `$, from the slower rotators in $`C_1`$ to the faster rotators in $`C_6`$ (see figure 11). The C<sup>-</sup> reconstructed cumulative luminosity function is shown for each class in figure 12. If the binning in the log line-width parameter $`\eta `$ were narrow enough, one would expect the reconstructed luminosity function for the class $`C_i`$ to be centered on $`M_i=a_{\mathrm{DTF}}\eta _i+b_{\mathrm{DTF}}`$ and of dispersion $`\sigma _{\mathrm{DTF}}(\eta _i)`$, where $`\eta _i=<\eta >_{C_i}`$ denotes the mean of the $`\eta `$’s in $`C_i`$ and $`a_{\mathrm{DTF}}`$, $`b_{\mathrm{DTF}}`$ and $`\sigma _{\mathrm{DTF}}(\eta _i)`$ are respectively the slope, zero-point and dispersion of the Direct (i.e. Forward) Tully-Fisher relation. We plot in figure 12 the expected culmulative luminosity function when a linear DTF relation with gaussian residuals is assumed and using the values proposed in Willick et al. (1998) for the calibration parameters. The width of the bins has been accounted for by adding in quadrature to the dispersion $`\sigma _{\mathrm{DTF}}(\eta _i)`$ the product of $`\|a_{\mathrm{DTF}}\|`$ and the standard deviation of the $`\eta `$’s in each class $`C_i`$ . For convenience in comparison, Each C<sup>-</sup> reconstructed CLF $`F_{\mathrm{rec}}(M)`$ has been normalised such that $`F_{\mathrm{rec}}(M_i)=0.5`$. It is not clear from figure 12 whether a linear DTF relation with gaussian residuals is plainly successful in reproducing the data. Note in particular that the C<sup>-</sup> reconstructed CLF’s for the slow rotators do not exhibit a turnover towards the faint magnitudes. To question the validity of a linear Tully-Fisher relation is anyway beyond the scope of the present paper. At this stage it is however worthwhile to mention that our method makes a very conservative use of Tully-Fisher information. In fact, we require to assume neither a linear TF relation nor a gaussian distribution for its residuals. The result of the ROBUST method applied to the MarkIII MAT catalog is shown in figure 13. The analysis was performed as described in section 3.2. The mean effective depth of the subsample is 2100 km s<sup>-1</sup> (see figure 14). We find a value of $`\beta =0.6\pm 0.125`$, in complete agreement with the VELMOD and ITF method applied to similar Tully-Fisher data (Willick et al. 1997a and 1998, Davis et al. 1996, Da Costa et al. 1998). We want however to emphasise the robustness of our approach compared to these two fitting methods. Firstly, no assumptions have been made herein concerning the linearity of the Tully-Fisher law, as is required by both the VELMOD and the ITF method. Secondly, we do not need the sample to be free of selection effects in the log line-width parameter $`\eta `$ as is the case for the ITF method. Thirdly, the spatial distribution of the sources, the selection effects in redshift and the shape of the distribution function of the TF residuals need not be specified, as is required by the maximum likelihood VELMOD method. ## 5 Conclusion We presented a method for fitting peculiar velocity models to complete flux limited magnitude-redshift catalogues, using the luminosity function of the sources as a distance indicator, i.e. assuming that the distribution function of the absolute magnitudes of the galaxies does not depend on the spatial position. Our method is based on a null-correlation approach. For a given peculiar velocity field model parametrised by a parameter $`\beta `$, we defined a random variable $`\zeta _\beta `$, computable from the observed redshifts and apparent magnitudes of the sampled galaxies, which has the property of being statistically independent on the position in space (and thus on the modelled radial peculiar velocities themselves) if and only if the parameter $`\beta `$ matches its true value $`\beta ^{}`$. Therefore any test of independence between the random variable $`\zeta _\beta `$ and the modelled velocities or similar quantities provides us with an unbiased estimate of the value of $`\beta ^{}`$. The method can be easily generalised to velocity models parametrised by an $`N`$-dimensional vector $`\beta =(\beta _1,\beta _2,\mathrm{},\beta _N)`$. The method is characterised by its robustness. No assumptions are made concerning the spatial distribution of sources and their luminosity function and selection effects in redshifts are also allowed. The required strict completeness in apparent magnitude can moreover be checked independently (Rauzy, in preparation). Furthermore the inclusion of additional observables correlated with the absolute magnitude is straightforward. The predicted IRAS peculiar velocity model characterised by the density parameter $`\beta `$ has been tested on two samples, the Tully-Fisher MarkIII MAT sample and the 60 $`\mu `$m IRAS 1.2 Jy sample using the fluxes as the distance indicator. The application of our method to the MarkIII MAT sample gives a value of $`\beta =0.6\pm 0.125`$, in excellent agreement with the results obtained previously by the VELMOD and ITF methods on similar datasets. Our method is however more robust than these two fitting methods. In particular, we make a very conservative use of the Tully-Fisher information. We do not require to assume the linearity of the Tully-Fisher relation nor a gaussian distribution of its residuals. We showed that our method allows to extract some valuable informations on the peculiar velocity field from the fluxes of the IRAS 1.2 Jy sample. The poor accuracy of the distance indicator (due to the broad spread of the luminosity function) is balanced in this case thanks to the large number of galaxies contained in the sample. The IRAS sample permits to probe the cosmic flow at larger scales. Indeed, the mean effective depth of the volume in which the velocity model is compared to the data is almost twice the mean effective depth of the MarkIII MAT sample. The application of our method to an IRAS subsample truncated in distance, of an effective depth similar to the MarkIII MAT sample, gives a value of $`\beta `$ in accord with the values obtained using Tully-Fisher data. On the other hand when the application is performed on the whole sample, we found that the predicted IRAS velocity models with $`\beta 0.5`$ can be rejected with a confidence level of $`95\%`$. These results suggest that the predicted IRAS velocity model, while successful in reproducing locally the cosmic flow, fails to describe the kinematics on larger scales. Note that these results do not lead to dismiss the linear “biasing” paradigm. As the errors on the predicted IRAS velocity field increase with distances, it could be that the predictions at the scales considered herein, i.e. beyond $`5000`$ km s<sup>-1</sup>, drastically differ from the true cosmic flow (see for example Davis et al. 1995). ## Acknowledgements We are thankful to Michael Strauss for providing us with the predicted IRAS peculiar velocity model. SR acknowledges the support of the PPARC and both authors acknowledge the use of the STARLINK computer node at Glasgow University.
warning/0004/cs0004006.html	ar5iv	text	# On Redundancy Elimination Tolerant Scheduling RulesThis work has been partially supported by the Italian National Research Council (CNR) research project “Tecniche di taglio dei cicli e loro implementazione in ambiente di Programmazione Logica”, Grant No. 97.02432.CT12. ## 1 Introduction Several different approaches have been considered so far to enrich the SLD resolution in order to improve the performance of top-down interpreters. The usual objective is to reduce the search space without loss of results of the refutation process, possibly obtaining a finite search space. Among the proposed methods, the loop check mechanisms \[Apt, Bol and Klop, 1989\], \[Bol, Apt and Klop, 1991\], \[Smith, Genesereth and Ginsberg , 1986\], \[Van Gelder, 1987\], and the tabulation technique \[Bol & Degersted, 1998\], \[Dietrich, 1987\], \[Ramakrishnan et al., 1999\], \[Tamaki & Sato, 1986\], \[Vieille, 1989\], aim to eliminate redundant computations and to enforce the termination of a query over a logic program. Loop check mechanisms provide the interpreter with the capability of pruning certain nodes of the SLD tree. The pruning is based on excluding some kinds of structural repetitions for the goals in a derivation path. When suitable structure repetitions are found, further rewritings of the current node are ignored, because any solution possibly existing in the cut sub-tree is also present in other parts of the SLD tree. Different forms of loop checks are proposed in the literature. In particular, Bol et al. have defined several simple loop checks, i.e. loop checks whose pruning mechanisms do not depend on the considered logic program, and have analysed them against the basic property of soundness and completeness \[Bol, Apt and Klop, 1991\]. The completeness property concerns with the capability of pruning every infinite derivation. In contrast, soundness concerns with the preservation of the computed answer substitutions. The main idea of tabulation originates from functional programming and consists in building a table during the search of answers in an SLD tree. The table contains entries for atoms with the corresponding answers so far computed. These answers are to be used later, when instances of such atoms should be recomputed. Such instantiated occurrences are named non-admissible atoms (or consumer). In essence, non-admissible atoms are not resolved against clauses but against answers computed in other parts of the SLD tree. The re-using approach exploited by the tabulation technique was already mentioned by Kowalski \[Kowalski, 1979\] and has been proposed several times under different names, such as memo-isation \[Dietrich, 1987\], and AL-technique \[Vieille, 1989\]. The conceptual differences between loop checks and tabulation are reflected in several interesting aspects. In particular, tabulation requires a local selection rule to guarantee the answer preservation, while no missing of solution is possible with (sound) loop checks independently of the used selection rule. On the other hand, the tabulation technique ensures termination for any function-free program and for any program with a finite Herbrand model, while the completeness of loop checks takes place for specific classes of programs possibly with respect to given selection rules \[Bol, 1992\], \[Bol, Apt and Klop, 1991\], \[Pacini & Sessa, 2000\]. Finally, loop checks exploit no auxiliary data structure and the pruning decision usually depends on the current derivation only, while tabulation needs a table to store the answers of atoms solved in the previously traversed portion of the tree. Proposals can be also found in literature for a synergistic use of different techniques aiming to optimise the query evaluation procedure. In particular, in \[Vieille, 1989\] a loop checking mechanism is combined with the tabulation technique in order to eliminate some redundant parts of the search space. In \[Ferrucci, Pacini and Sessa, 1995\] the simple loop check mechanisms proposed in \[Bol, Apt and Klop, 1991\] are combined with another form of redundancy elimination which is named (goal) reduction. Goal reduction is conceptually analogous to the condensing technique proposed by Joyner for the proof of the unsatisfiability of first-order formulas \[Joyner, 1976\]. In both cases redundant atoms are eliminated from resolvents, in order to avoid useless computations and to contain the size of the resolvents at the same time. The main idea of reduction originates from the observation that if there exists a refutation for an atom, then a refutation exists also for any more general version of that atom. In this sense, such more general versions can be seen as potentially redundant and we can imagine to remove them from the resolvent, though suitable cares are to be taken as discussed in \[Ferrucci, Pacini and Sessa, 1995\]. By goal reduction, a generalised form of SLD resolution (named RSLD) can be obtained, where a reduction of the resolvent is performed after each rewriting step. Goal reduction technique has a modus operandi which shows evident affinity with the one of loop checking mechanisms. Indeed, with reduction redundant atoms are definitively ignored, as it is done with loop checks for pruned nodes. This is not the case with tabulation, in the sense that non-admissible atoms, which are indeed solved against previously tabulated answers, are not redundant. Such different philosophy between tabulation and RSLD is highlighted also by the fact that the reduction technique eliminates atoms in their more general version, while non-admissible atoms are instances of previously solved goals. It is evident that RSLD does not need any auxiliary data structure because it considers only the current goal (not even the current derivation path). The soundness of RSLD is shown in \[Ferrucci, Pacini and Sessa, 1995\] independently of the used selection rule. This means that RSLD does not require particular selection rules in order to ensure answer preservation. It is intuitive that redundancy elimination may have positive effects on derivation process. In \[Ferrucci, Pacini and Sessa, 1995\], advantageous combinations are shown with respect to loop checking mechanisms. In particular, it is proven that a well known simple loop check mechanism, namely Equality Variant check of Resultant as Lists ($`EVR_L`$), becomes complete for several classes of programs, provided that RSLD is exploited instead of usual SLD. The specific reason is that the length of resolvents can be maintained within the limit of a finite value through systematic elimination of redundant atoms. In essence, there is clear evidence that the strength of equality loop checks can augment if RSLD resolution is used. However, even though not completely intuitive, redundancy elimination can produce undesirable effects, too. In fact, as exemplified later, problems can arise with program termination, as well as with the completeness of loop checking mechanisms. The rationale behind this is that redundancy elimination can affect the actual sequence of atom rewriting with respect to given selection rules. This can (infinitely) delay the selection of failing atoms, so that termination is missed. On the other hand, the structure of the obtained resolvents can be altered by redundancy elimination, so that loop checks may become unable to detect infinite derivations. As shown in this paper, missing termination and loop detection depends critically on the used selection rule. We say in the sequel that a selection rule is redundancy elimination tolerant if no loss in termination and/or loop detection comes out, passing from SLD to RSLD. In Section 2, we prove that termination and $`EVR_L`$ completeness are preserved if they hold in SLD with respect to all possible selection rules. Then, a more accurate analysis of redundancy elimination tolerance is performed. To this aim, a careful reconsideration of selection rule basic concepts will be required, so that we will be led to a reformulation of selection rule ideas in terms of their operational counterparts, namely scheduling mechanisms, so that we will prefer to talk of tolerant scheduling rules. As a matter of fact, in Section 3 we provide a highly expressive execution model based on priority mechanism for atom selection. A priority is assigned to each atom in a resolvent, and primary importance is given to the event of arrival of new atoms from the body of the applied clause at rewriting time. Indeed, new atoms can be freely positioned with respect to the old ones in the resolvent, through the assignment of priority values according to a given scheduling rule. Then, at any derivation step, the atom with optimum priority is simply selected. This new computational model proves able to address the study of redundancy elimination effects, giving at the same time interesting insights into general properties of selection rules. As a matter of fact, in Section 4 a class of scheduling rules, namely the specialisation independent ones, is defined by using not trivial semantic arguments. Several properties of specialisation independent scheduling rules are also proven. As a quite surprising result, in Section 5 we show that specialisation independent scheduling rules coincide with stack-queue rules, which have an immediate structural characterisation. Indeed, the stack-queue scheduling technique is simply defined so that, in order to obtain the new resolvent at rewriting time, part of new atoms are stacked at the beginning of the old resolvent while the remaining ones are queued. Then in Section 6 we prove that such scheduling rules are tolerant to redundancy elimination, in the sense that neither program termination nor completeness of equality loop check is lost passing from SLD to RSLD. The proof is largely based on properties which we have established for specialisation independent (and stack-queue) scheduling rules. ## 2 Goal reduction, program termination and <br>$`EVR_L`$ completeness Throughout the paper we assume familiarity with the basic concepts of Logic Programming \[Apt, 1990\], \[Apt, 1998\], \[Lloyd, 1987\]. Here, only some notations are given about SLD derivation procedure, which can be described as follows. Let $`G=a_1,a_2,\mathrm{}a_k`$ be a goal, constituted by a conjunction of $`k`$ atoms, and $`c=(htB)`$ a clause, where$`ht`$is an atom and$`B`$is a goal. The goal $`G^{}`$ is a resolvent of$`G`$and$`c`$by a renaming $`\xi `$ and a substitution $`\theta `$, if an atom $`a_i`$ exists, with $`1ik`$, such that $`G^{}=(a_1,\mathrm{}a_{i1},B\xi ,a_{i+1},\mathrm{}a_k)\theta `$, where $`\theta `$ is an idempotent and relevant mgu of $`(ht)\xi `$ and $`a_i`$. In the sequel, given an expression $`E`$, the notation $`var(E)`$ will indicate the set of variables in $`E`$. Moreover, we will denote by $`(G\stackrel{c\xi ,\theta }{}G^{})`$ the fact that $`G^{}`$ is a resolvent of$`G`$and$`c`$by $`\xi `$ and $`\theta `$. Given an initial goal $`G_o`$ and a logic program P, an SLD derivation of $`G_o`$ in P is a possibly infinite sequence of the type: $`G_o\stackrel{c_o\xi _o,\theta _o}{}G_1\mathrm{}G_j\stackrel{c_j\xi _j,\theta _j}{}G_{j+1}\mathrm{}`$ such that, for any $`j0`$, each clause $`c_j`$ belongs to P and each $`c_j\xi _j`$ is standardised apart, i.e. $`var(c_j\xi _j)(var(G_o)var(c_o\xi _o)\mathrm{}var(c_{j1}\xi _{j1}))=\mathrm{}.`$ A selection rule is a function which chooses the atom to be rewritten in the last resolvent of any finite SLD derivation. Given a selection rule$`S,`$ an SLD derivation is via$`S`$if all the selections of atoms are performed in agreement with$`S.`$ An SLD refutation is a finite SLD derivation such that the last resolvent is empty. Now we can introduce the definitions of goal reduction and RSLD derivation. The reduction technique aims to eliminate redundant atoms from the resolvents in order to contain their size. Analogous issue was already been faced for the proof of the unsatisfiability of first-order formulas. Indeed Joyner \[Joyner, 1976\] noted that the increase in size of resolvents is a factor which prevents resolution strategies being decision procedures for solvable classes of first-order formulas (i.e. classes of formulas for which the question of satisfiability or unsatisfiability can be effectively decided). To limit the growth of the number of literals, Joyner introduced a technique for simplifying resolvents, called condensing. The condensation of a clause is defined as the smallest subset of the clauses which is also an instance of it. In other words, the condensation of a clause can be obtained by applying a substitution $`\alpha `$ and eliminating all the atom repetitions. With reference to SLD derivations, the most evident form of redundancy corresponds to multiple occurrences of the same atom in a resolvent. It is obvious that this kind of atom repetition is essentially redundant. However, this is not the only possible case of redundancy. Indeed, the reduction technique, which is introduced in \[Ferrucci, Pacini and Sessa, 1995\] as a variant of Joyner’s condensing technique, is able to perform quite general actions of redundancy elimination from resolvents while preserving the soundness and the completeness of RSLD resolution. By condensation, Joyner obtains a complete and sound resolution procedures, which work as decision procedures for several solvable classes of first order formulas \[Joyner, 1976\]. By reduction, the well known sound $`EVR_L`$ loop check becomes complete for several classes of logic programs \[Ferrucci, Pacini and Sessa, 1995\]. Intuitively, the basic idea of goal reduction technique can be explained as follows. Suppose having to refute a resolvent which contains $`p(x)`$ and $`p(a)`$, where $`x`$ is a variable and $`a`$ is a constant. Obviously, any refutation for $`p(a)`$ implies a refutation for the atom $`p(x)`$, as $`p(x)`$ is more general than $`p(a)`$. In this sense, the atom $`p(x)`$ may appear as a redundant one. Actually, in order to ensure the soundness of the derivation process, the elimination of redundant atoms (such as $`p(x)`$ above) is conditioned in two aspects which can be sketched through the following simple examples: 1. Consider a resolvent like $`p(x),q(x),p(a)`$. In this case, the atom $`p(x)`$ cannot be eliminated, because the connection between the atoms $`p(x)`$ and $`q(x)`$, by the variable $`x`$, is lost. 2. Suppose that $`x`$ is a variable in the initial goal of a derivation, and the actual resolvent is $`p(x),p(a)`$. In this case $`p(x)`$ cannot be dropped, because possible instantiations of $`x`$ in computed answers could be lost. So we would obtain computed answers which are too general with respect the correct answers, thus missing soundness. Now we present a formal definition of goal reduction which takes into account the observations a) and b) and follows the line of Definition 2.1 presented in \[Ferrucci, Pacini and Sessa, 1995\]. We will denote by $`_L`$ the inclusion relation between goals, and $`GN`$ will indicate the goal obtained from $`G`$ by eliminating the atoms which are present in $`N`$. In both cases the goals are regarded as lists. ###### Definition 2.1 (Reduced goal) Let $`X`$ be a set of variables, $`\tau `$ a substitution and$`G`$a goal. A goal$`N`$is a reduced goal of$`G`$by $`\tau `$ up to $`X`$, denoted by $`G>>^\tau N`$, if the following conditions hold: i) $`N_LG,`$ ii) $`b(GN),b\tau N,`$ iii) $`x(var(N)X)`$it is $`x\tau =x.`$ In agreement with the above definition, a part $`(GN)`$ of atoms of$`G`$can be eliminated if a substitution $`\tau `$ exists such that $`b\tau N`$, for any atom $`b(GN)`$, provided that $`\tau `$ does not affect neither the variables in$`N`$nor those in $`X`$. The imposition that $`\tau `$ does not affect the variables in$`N`$prevents the kind of difficulties which are exemplified in a). ###### Example 2.1 Let: $`G=p(z,v),q(w),p(w,v),p(w,x),p(w,y),q(v),q(y),`$ $`X=\{x,w\}.`$ The following goal$`N`$is a reduced goal of$`G`$by $`\tau =\{z/w,y/v\}`$ up to $`X`$: $`N=q(w),p(w,v),p(w,x),q(v).`$ $`\mathrm{}`$ Performing reductions in the resolvents of an SLD derivation corresponds to an actual extension of the SLD resolution process. Then, a generalised version of SLD resolution can be introduced, i.e. the Reduced SLD resolution (RSLD in the sequel), where at any resolution step a reduction of the resolvent is allowed. The following is the formal definition of RSLD derivations. ###### Definition 2.2 (Reduced SLD derivation) Let $`P`$ be a program and $`G_o`$ a goal. A Reduced SLD derivation of $`G_o`$ in $`P`$ (RSLD in the following) is a possibly infinite sequence of the form: $`G_o>>^{\alpha _o}N_o\stackrel{c_o\xi _o,\theta _o}{}G_1\mathrm{}G_h>>^{\alpha _h}N_h\stackrel{c_h\xi _h,\theta _h}{}G_{h+1}>>^{\alpha _{h+1}}N_{h+1}\mathrm{}`$ where, for any $`j0,`$ i) $`c_j`$ is a clause in $`P`$, ii) $`var(c_j\xi _j)(var(G_o)var(c_o\xi _o)\mathrm{}var(c_{j1}\xi _{j1}))=\mathrm{},`$ iii) $`G_j>>^{\alpha _j}N_j`$ up to $`var(G_o\theta _o\mathrm{}\theta _{j1}).`$ It is evident that an SLD derivation is a particular case of RSLD derivation where $`G_j=N_j`$, for any $`j`$. Each $`N_j`$ is called a reduced resolvent. Condition ii) above is the usual standardisation apart requirement. Condition iii) prevents the kind of difficulties which are exemplified in b), guaranteeing the soundness of the mechanism. The soundness and completeness of RSLD resolution are proven in Theorems 2.1 and 2.2 of \[Ferrucci, Pacini and Sessa, 1995\]. ### 2.1 Program termination The completeness of RSLD resolution ensures that missing computed answers is impossible when we pass from SLD to RSLD. This is not the case with termination, as shown by the following Example 2.2. In the example a selection rule$`S`$and a program P are given, such that any SLD derivation of P via$`S`$terminates independently of the initial goal. However, we show that termination is lost, if reduction of resolvents is performed. ###### Example 2.2 Let us consider a selection rule$`S`$such that, given a goal $`G`$, the first atom is chosen for rewriting if the length of$`G`$is odd, and the last atom is chosen otherwise. Let us consider the logic program P consisting of the following clause: $`c=p(x,y)q,p(x,z_1),p(z_1,z_2),p(z_2,y).`$ It can be easily seen that all SLD derivations in P via$`S`$terminate, independently of the initial goal. Indeed, suppose that the initial goal has an odd number of atoms. It is evident that either the derivation via$`S`$fails immediately or the initial goal has the form “$`p(..),Y`$”, so that the first step of the derivation produces a resolvent of an even length as follows: $`p(..),Y\stackrel{𝑐}{}q,p(..),p(..),p(..),Y.`$ Now, either the derivation fails immediately or $`Y=Z,p(..)`$, so that a second derivation step is performed: $`q,p(..),p(..),p(..),Z,p(..)\stackrel{𝑐}{}q,p(..),p(..),p(..),Z,q,p(..),p(..),p(..)`$, and the process fails anyway, since the last resolvent has an odd length. Then, suppose on the contrary that the initial goal has an even number of atoms. Either the derivation fails immediately or the initial goal has the form “$`T,p(..)`$”. In the second case, the first derivation step gives place to a resolvent with an odd length, so that the derivation fails. Now, let us verify that termination can be lost if reduction of resolvents is performed. Indeed, let us consider the RSLD derivation of the goal $`(q,p(x,x))`$ in P via $`S`$ given in Figure 1. It is evident that the number of atoms is even in any reduced resolvent. Thus, the last atom is always selected and the derivation is infinite. $`\mathrm{}`$ As shown by the example in Figure 1, termination with respect to a given selection rule can be missed, if we pass from SLD to RSLD resolution. On the contrary, we show in this section (Theorem 2.1) that termination is preserved, when any SLD derivation of$`G`$in P is finite independently of the used selection rule. Theorem 2.1 will be proven as an immediate consequence of the following Lemma 2.1 ###### Lemma 2.1 Let P be a program and $`G_o`$ a goal. For any possibly infinite RSLD derivation $`D`$ of $`G_o`$ in P, an SLD derivation $`D^{}`$ of $`G_o`$ in P exists, such that every reduced resolvent of $`D`$ is included in the corresponding resolvent of $`D^{}`$ up to renamings. * * Consider a possibly infinite RSLD derivation $`D`$ of $`G_o`$ in P $`D=(G_o>>^{\alpha _o}N_o\stackrel{c_o\xi _o,\theta _o}{}G_1\mathrm{}`$ $`\mathrm{}G_h>>^{\alpha _h}N_h\stackrel{c_h\xi _h,\theta _h}{}G_{h+1}>>^{\alpha _{h+1}}N_{h+1}\mathrm{})`$ (1) Intuitively, the SLD derivation $`D^{}`$ is obtained by choosing, step by step, the same clause and the same atom as in $`D`$. This way, redundant atoms are not eliminated from resolvents of $`D^{}`$, but they have no real influence on the derivation process. More formally, suppose that an SLD derivation of $`G_o`$ in P is already constructed like $`G_o\stackrel{c_o\xi _o^{},\theta _o^{}}{}G_1^{}\mathrm{}G_i^{},`$(2) such that, for any $`0ji`$, a renaming $`\tau _j`$ exists with $`N_j\tau _j_LG_j^{}.`$ It is easy to show that derivation (2) can be extended of one step in agreement with the lemma. Let $`a`$ be the atom which is rewritten in the step $`N_i\stackrel{c_i\xi _i,\theta _i}{}G_{i+1}`$ of derivation (1). It is evident that the clause $`c_i`$ is applicable to the atom $`a\tau _iN_i\tau _i_LG_i^{},`$ so that we have an SLD derivation step of the form: $`G_i^{}\stackrel{c_i\xi _i^{},\theta _i^{}}{}G_{i+1}^{}.`$(3) Now let $`F`$ denote the sublist of atoms in $`G_{i+1}^{}`$ which derives from $`N_i\tau _i`$. It is obvious that the subgoal $`(G_i^{}N_i\tau _i)`$ has no active role in derivation step (3). So, we have that $`F`$ is a variant of $`G_{i+1}`$, i.e. a renaming $`\tau _{i+1}`$ exists with $`F=G_{i+1}\tau _{i+1}`$, which means that $`G_{i+1}\tau _{i+1}_LG_{i+1}^{}`$. But, by definition of goal reduction we have $`N_{i+1}_LG_{i+1}.`$ As a consequence $`N_{i+1}\tau _{i+1}_LG_{i+1}\tau _{i+1}_LG_{i+1}^{}.`$ $`\mathrm{}`$ ###### Theorem 2.1 Let P be a program and$`G`$a goal. If every SLD derivation of$`G`$in P is finite independently of the used selection rule, then every RSLD derivation of$`G`$in P is finite too. ###### Proof 2.2. Suppose that an infinite RSLD derivation of$`G`$ in P exists. By Lemma 2.1, an infinite SLD derivation of$`G`$in P also exists, which contradicts the hypothesis. ### 2.2 $`EVR_L`$ loop check completeness The termination issue of a query to a logic program has attracted much attention over the past few years, both in the logic programming field, and in the deductive database field (see \[De Shreye & Decorte, 1994\] for a survey). A well known approach to the termination problem of a query in a logic program consists in modifying the computation mechanism by adding a capability of pruning, i.e. at some point the interpreter is forced to stop its search through a certain part of the SLD tree \[Apt, Bol and Klop, 1989\], \[Bol, 1992\], \[Bol, Apt and Klop, 1991\], \[Pacini & Sessa, 2000\], \[Smith, Genesereth and Ginsberg , 1986\], \[Van Gelder, 1987\]. These mechanisms are called loop checks, as they are based on discovering some kinds of repetitions in derivation paths. The purpose of a loop check is to reduce the search space for top-down interpreters in order to prune infinite derivations, without loss of results of the refutation process. Thus, two basic properties are considered for loop checks. The completeness property of a loop check concerns the capability of pruning every infinite derivation. In contrast, the soundness property has to do with the preservation of computed answer substitutions. Different forms of loop checking are considered in literature. A systematic analysis of loop checking for SLD resolution is given in \[Bol, Apt and Klop, 1991\]. Simple loop checks have deserved special interest, because the decision of pruning does not depend on the logic program we are confronted with. The more immediate form of simple and sound loop check is the so called Equality Variant of Resultant check, which requires the detection of equal (up to renaming) resultants in the derivation. Such a loop check is formulated with respect to RSLD derivations in the following Definition 2.3 which recalls the essence of the analogous Definition 3.19 in \[Ferrucci, Pacini and Sessa, 1995\]. The notation $`(F=_LG)`$ is used, which means that the goal $`F`$ is equal to $`G`$, where the goals are regarded as lists. ###### Definition 2.3 (Equality Variant Check for Resultants). An RSLD derivation $`G_o>>^{\alpha _o}N_o\stackrel{c_o\xi _o,\theta _o}{}G_1\mathrm{}G_{h1}>>^{\alpha _{h1}}N_{h1}\stackrel{c_{h1}\xi _{h1},\theta _{h1}}{}G_h>>^{\alpha _h}N_h\mathrm{}`$ is pruned by Equality Variant of Resultant as Lists loop check ($`EVR_L`$ in the following), if for some $`i`$ and $`j`$, with $`0i<j`$, a renaming $`\tau `$ exists such that: i) $`G_o\theta _o\mathrm{}\theta _{j1}=G_o\theta _o\mathrm{}\theta _{i1}\tau ,`$ ii) $`N_j=_LN_i\tau .`$ Given an RSLD tree T, the application of $`EVR_L`$ yields a prefix Tp of T which is obtained in this way. The descendants of a node are thrown away iff the derivation associated with the path from the root to the node is pruned. Any couple $`Rs_h=[N_h,G_o\theta _o\mathrm{}\theta _{h1}]`$ is a reduced resultant. Given two resultants $`Rs_j=[N_j,G_o\theta _o\mathrm{}\theta _{j1}]`$ and $`Rs_i=[N_i,G_o\theta _o\mathrm{}\theta _{i1}],`$ for which requirements i) and ii) of Definition 2.3 hold, we will write $`Rs_i_LRs_j.`$ In other words, Definition 2.3 expresses that $`EVR_L`$ check is based on detecting that a resultant is obtained which is related by $`_L`$ to a preceding one in the same derivation. It is worth noting that the relationship $`_L`$ is an equivalence relationship. It is evident that, if reduction of resolvents is always ineffective (i.e. $`G_j=N_j`$, for any $`j`$), the usual $`EVR_L`$ loop check for SLD derivations is found again. It is well known that $`EVR_L`$ is a sound loop check in the case of SLD resolution. The soundness of $`EVR_L`$ is extended to the more general case of RSLD by Theorem 4.1 of \[Ferrucci, Pacini and Sessa, 1995\]. Let us observe that if we do not consider condition i) in Definition 2.3 we obtain the $`EVG_L`$ loop check which is based on detecting that a resolvent is obtained which is a variant of a preceding one in the same derivation. It is worth noting that $`EVG_L`$ is a weakly sound loop check, in sense that it preserves at least a successful, but it does not ensure the preservation of the computed answer substitutions \[Bol, Apt and Klop, 1991\]. The completeness of a loop check is usually referred to given selection rules and classes of programs. A loop check is complete for a program P with respect to a selection rule$`S`$if all infinite derivations of P via$`S`$are pruned. A loop check is complete for a class C of programs, if it is complete for every program in C. Several classes of logic programs are characterised in literature for which complete loop checks can be found. Actually, most of them are classes of function free programs, i.e. programs whose clauses contain no function symbol \[Bol, 1992\], \[Bol, Apt and Klop, 1991\], \[Ferrucci, Pacini and Sessa, 1995\], \[Pacini & Sessa, 2000\]. In the following of this section, and later in Section 6, we consider the problem of preserving the completeness of $`EVR_L`$ check, passing from SLD to RSLD resolution, in the case of function free programs. Let us first show how the completeness of equality loop checks, with respect to a given selection rule, can be lost passing from SLD to RSLD. Indeed, it is sufficient reconsider Example 2.2. In that case $`EVR_L`$ loop check is obviously complete, since no infinite SLD derivation exists. On the other hand, it is obvious that $`EVR_L`$ loop check cannot prune the infinite RSLD derivation developed in the same example, because the length of resolvents increases at each derivation step. Actually, it is immediate to verify that the infinite derivation in Example 2.2 cannot even be pruned by using more complex and powerful checks (like $`SIR_M`$) which are based on subsumption relationships between resultants \[Bol, Apt and Klop, 1991\]. Now we prove that $`EVR_L`$ loop check completeness is preserved for function free programs, in the case that $`EVR_L`$ is complete with respect to all selection rules. Precisely, Theorem 2.6 states that, if $`EVR_L`$ prunes every infinite SLD derivation of a goal$`G`$in a function free program P, then $`EVR_L`$ prunes also every infinite RSLD derivation of$`G`$in P. In order to show this result, let us provide a condition which holds whenever $`EVR_L`$ prunes every infinite derivation of$`G`$in P. Lemma 2.4 states that, if $`EVR_L`$ check prunes all infinite derivations of$`G`$in P, then the length of resolvents in all possible derivations is limited. In the proof of Lemma 2.4 we exploit the notion of S-tree \[Apt & Pedreschi, 1993\]. Given a program P and a goal $`G`$, an S-tree of$`G`$in P is a tree where the descendants of a goal are its resolvents with respect to all selection rules and all input clauses. In other words, an S-tree groups all SLD derivations of$`G`$in P. The notation $`\mathrm{\#}R`$ represents the number of atoms in the goal $`R`$. ###### Lemma 2.4. Let P be a program and$`G`$a goal. Suppose that all infinite SLD derivations of$`G`$in P are pruned by $`EVR_L`$. Then, a finite bound $`l`$ exists such that, for each resolvent$`R`$in any SLD derivation of$`G`$in P, it is $`\mathrm{\#}Rl`$. ###### Proof 2.5. Let T be an S-tree of$`G`$in P. Given a node$`n`$ in T, let $`Dr(n)`$ denote the derivation associated to the path from the root of T to $`n`$, and $`R(n)`$ the final resolvent of $`Dr(n)`$. Then, let Tp be the prefix of T which is obtained by applying the $`EVR_L`$ check to T, i.e. the prefix where the descendants of any node$`n`$ of T are thrown away if and only the derivation $`Dr(n)`$ is not pruned by $`EVR_L`$. By hypothesis, all infinite SLD derivations of$`G`$in P are pruned by $`EVR_L`$, which means that Tp has no infinite path. As a consequence, since T is a finitely branching tree, by Konig’s lemma (see Theorem K, in \[Knuth, 1997\]) the prefix Tp is finite. Now, let $`d`$ be the depth of Tp, and $`l`$ the maximum of the set $`\{\mathrm{\#}R(n)\|n`$ is a node in Tp$`\}`$. We prove that: $`\mathrm{\#}R(n)l,`$ for any node$`n`$ in T. The proof is by induction on the value of $`depth(n)`$. For $`depth(n)d`$ the thesis is trivial. Then consider an integer $`h>d`$, and suppose that $`\mathrm{\#}R(n^{})l`$, for any node $`n^{}`$ with $`depth(n^{})<h.`$ Given a node$`n`$ of T such that $`depth(n)=h`$, we show that also $`\mathrm{\#}R(n)l`$ holds. Since$`n`$Tp, the derivation $`Dr(n)`$ is pruned by $`EVR_L`$, so that two nodes $`n_1`$ and $`n_2`$ exist in the path from the root of T to$`n`$ with: \- $`depth(n_1)<depth(n_2),`$ (1) \- $`R(n_2)`$ is a variant of $`R(n_1)`$. (2) Now, consider the sequence of clauses which has determined the path from $`n_2`$ to$`n`$ in T. Since T contains all SLD derivations of$`G`$in P, the same derivation steps can be repeated in T starting from $`n_1`$. As a consequence, by (1) and (2), a path from $`n_1`$ to a node $`n^{}`$ exists such that: \- $`depth(n^{})=depth(n)(depth(n_2)depth(n_1))<depth(n)=h`$, \- $`R(n^{})`$ is a variant of $`R(n)`$. By inductive hypothesis it is $`\mathrm{\#}R(n^{})l`$. But $`R(n^{})`$ is a variant of $`R(n)`$, so that$`\mathrm{\#}R(n)=\mathrm{\#}R(n^{})l`$. In conclusion, the thesis holds for every node$`n`$ in T. ###### Theorem 2.6. Let P be a function free program and$`G`$a goal. If $`EVR_L`$ prunes every infinite SLD derivation of$`G`$in P independently of the used selection rule, then $`EVR_L`$ prunes every infinite RSLD derivation of $`G`$ in P. ###### Proof 2.7. Let $`D`$ be an infinite RSLD derivation of$`G`$in P. By Lemma 2.1, an SLD derivation $`D^{}`$ of$`G`$in P also exists such that every reduced resolvent of $`D`$ is included in a resolvent of $`D^{}`$ (up to renamings). Since $`EVR_L`$ prunes every infinite SLD derivation of$`G`$in P, by Lemma 2.4 the length of resolvents of $`D^{}`$ is limited. Then, the length of reduced resolvents and resultants of $`D`$ is also limited. Now, since the language of P is function free and has finite many predicate symbols and constants, the relationship denoted by $`_L`$ has only finitely many equivalence classes on resultants of $`D`$. As a consequence, for some $`0i<k`$ we have that the $`k^{th}`$ and the $`i^{th}`$ resultants of $`D`$ are in $`_L`$ relationship. This implies that $`D`$ is pruned by $`EVR_L`$. In this section, redundancy elimination tolerance has been proven on the basis of a rather strong hypothesis, i.e. termination and completeness of loop checking for all possible selection rules. In Section 3 we will introduce a new computational model which will allow us to characterise a class of selection rules which are shown to be redundancy elimination tolerant. As a matter of fact, in Section 6 we will prove that program termination and $`EVR_L`$ loop check completeness are maintained for that class of rules, passing from SLD to RSLD. ## 3 Priority scheduling rules As shown in Section 2, redundancy elimination can determine missing termination and loop check detection. This fact depends critically on the used selection rule, because redundancy elimination can affect the actual sequence of atom rewriting. As a matter of fact, it is widely acknowledged that the analysis of interdependence between derivation processes and the used selection rules is a difficult task. In our study, the necessary insights have been provided by a computation model which is based on a novel mechanism of atom choice, which works in terms of scheduling rules rather than in terms of conventional selection rules. Through this new computational model, a class of scheduling rules is identified in Section 4, which is redundancy elimination tolerant in the sense that no loss in termination and/or loop detection comes out, passing from SLD to RSLD. We start the analysis with an observation about selection rules, as they are normally conceived in literature and used in practice. In SLD derivations, resolvents are usually regarded as lists, nevertheless selection rules are given complete free choice ability of the atom to rewrite. In this sense, two different philosophies are superimposed, because a scheduling (i.e. an ordering) must coexist with an atom choice which can actually overcome the scheduling. Now, in the case that resolvents are viewed as unstructured multisets instead of lists, the obvious solution is that a free choice ability is provided at rewriting time. But, if scheduling policies (i.e. an ordering or a priority assignment) are exploited, it may appear natural that priorities are obeyed at rewriting time, so that the atom with optimum priority is always selected. Indeed, if a scheduling policy is used, the moment of addition of new atoms in the resolvent may be recognised as the really important event, when suitable priority values must be established and assigned. In the following of the paper we consider execution mechanisms for logic programs which are based on priority scheduling policies. In particular we characterise scheduling rules informally as follows: * a priority value is assigned to each atom in the actual resolvent, * assigned priorities are not modified in the following of the derivation, * the atom with optimum priority is always taken for rewriting. In essence a scheduling rule is a rule that defines a priority values for any new atom which enters the actual resolvent. It is crucial that atoms from the body of the applied clause can be freely scheduled with respect to the ones already present in the resolvent, which maintain their own priority values. It is intuitive that this can be easily done if a set of “dense” priority values is adopted. Indeed, as formalised in Section 3.1, we will use rational numbers as priority values. Now, in analogy with Lloyd’s definition of selection rules \[Lloyd, 1987\], we consider the subclass of scheduling rules where the schedule of new atoms is determined only by the last resolvent in the derivation, i.e. by the current state of the computation. Such rules will be named state scheduling rules. A state scheduling rule can be seen as a rule which, for any resolvent $`G`$ and clause $`c`$ (that is applied to the optimum priority atom), determines the schedule positions of the new atoms in the resolvent, through the assignment of appropriate priority values. In other words, a state scheduling rule determines new resolvents, starting from the old ones and from applied clauses. The rewritten atom is necessarily the one with the optimum priority value. It is evident that the transformation from a resolvent to a new one, which is obtained by the addition of new atoms from the applied clause, is nothing more than a step of an SLD derivation. In this sense, we can say that a state scheduling rule characterises a set of derivation steps. Indeed, as formalised in Section 3.5, a state scheduling rule can be straight conceived as a set of derivation steps, that is: the set of derivation steps which are allowed according to the scheduling rule itself. Formal definition of state scheduling rules is provided in Section 3.5. ### 3.1 Atoms, goals and priorities In order to characterise state scheduling rules in a formal way, we introduce the notions of priority goal and priority clause. A priority goal is a goal where each atom has an associated priority value. Thus, a priority goal$`G`$can be thought as a set of couples, where any couple is named priority atom. In the following formal definition, priority atom will be denoted by $`a[p]`$, where $`a`$ is an usual atom and $`p`$ is a rational number which establishes the priority of $`a`$ in $`G`$. The symbol $``$ will be frequently used in the rest of the paper to denote logical implication. ###### Definition 3.1. 1. A priority goal$`G`$(p-goal in the sequel) is defined by a set of priority atoms (or simply p-atoms) of the form: $`G=\{a_1[p_1],\mathrm{}a_k[p_k]\},`$ with $`i,j:ijp_ip_j,`$ where each $`a_m`$ is an usual atom and each $`p_m`$ is a rational number, $`1mk`$. 2. A priority clause (or simply a p-clause) has the form $`c=htB`$, where$`ht`$is an atom (without priority) and$`B`$is a priority goal. In the sequel, priority clauses will be referred as clauses for the sake of simplicity. Capital letters will be used in the following to represent p-goals. In order to denote p-atoms, we will use notations like $`a[p]`$, as well as simple small letters (as $`a,b`$, etc.) when explicit reference to priority values is not important. As a slight abuse of notation, p-goals made of only one p-atom $`a`$ will be often denoted by $`a`$. Given a p-goal $`G`$, the notations $`\mathrm{\#}G`$ will indicate the number of p-atoms in $`G`$. In the sequel, we will exploit very frequently a basic operation which corresponds to the union of two p-goals with no common priority values. This operation is denoted by “+” and is said p-goal merging. During merging operations, atoms retain their priority values. We introduce also the idea of concatenation, which is a particular case of merging. Concatenations will be denoted by the symbol “$`\|`$” (vertical bar). The following are the formal definitions of merging and concatenation. It is worth noting that both these operations are associative. ###### Definition 3.2. 1. A p-goal$`M`$is the merging of$`F`$and$`G`$(denoted by $`M=F+G`$) if$`F`$and$`G`$have no common priority values and $`M=FG`$. 2. Given two p-goals$`F`$and $`G`$, we write$`FG`$to denote that all priorities in$`F`$are less than any priority in $`G.`$ A p-goal$`N`$is the concatenation of$`F`$and$`G`$(denoted by $`N=F\|G`$), if $`N=F+G`$ and$`FG`$. The fact that equal priority values are not admitted in a p-goal has two principal effects. The first one is that a complete ordering (i.e. a scheduling) is imposed on the atoms of a p-goal. In particular we assume that atoms with less priorities precede atoms with greater ones. The second effect is that possible multiple occurrences of atoms are distinguished by different priority values. On the basis of the above observations, the following evident properties of concatenation can be stated. ###### Property 3.1 Given the p-goals $`A_1,A_2,A_3,B_1,B_2`$, and $`B_3`$, the following propositions hold: 1. $`A_1\|A_2=B_1\|B_2,\mathrm{\#}A_1=\mathrm{\#}B_1`$ or $`\mathrm{\#}A_2=\mathrm{\#}B_2A_1=B_1,A_2=B_2.`$ 2. $`A_1\|A_2\|A_3=B_1\|B_2\|B_3,A_2\mathrm{},A_2=B_2A_1=B_1,A_3=B_3.`$ ### 3.2 Shifting and positioning Throughout the paper, we will exploit a basic operator for handling priority values. It will be said (priority) shifting, and corresponds to a modification of priority values which does not alter the scheduling of the atoms in a p-goal. The following is the formal definition of shifting. In the sequel, shiftings will be always denoted by underlined Greek letters. ###### Definition 3.3 (shifting). A shifting $`\underset{¯}{\pi }`$ is an increasing one-to-one application of the type: $`\underset{¯}{\pi }`$ : Rational $``$ Rational. Given a shifting $`\underset{¯}{\pi }`$, and two p-goals$`G`$and$`F`$such that: $`G=\{a_1[p_1],\mathrm{}a_k[p_k]\}`$ and $`F=\{a_1[\underset{¯}{\pi }(p_1)],\mathrm{}a_k[\underset{¯}{\pi }(p_k)]\}`$, we say that$`F`$is a shifting of$`G`$and write$`F=G\underset{¯}{\pi }`$. It is evident that the composition of two shiftings is a shifting, too, as well as the inverse of a shifting. Shifting operations enjoy the following four basic properties. All properties are plain consequence of the definition. The first two properties will be used very often in the sequel without explicit reference. ###### Property 3.2 Ax-i) $`(A_1+A_2+\mathrm{}+A_k)\underset{¯}{\pi }=A_1\underset{¯}{\pi }+A_2\underset{¯}{\pi }+\mathrm{}A_k\underset{¯}{\pi },`$ Ax-ii) $`(A_1\|A_2\|\mathrm{}A_k)\underset{¯}{\pi }=A_1\underset{¯}{\pi }\|A_2\underset{¯}{\pi }\|\mathrm{}A_k\underset{¯}{\pi },`$ Ax-iii) $`G=A_1\underset{¯}{\tau }_1\|A_2\underset{¯}{\tau }_2\|\mathrm{}A_k\underset{¯}{\tau }_k,F=A_1\underset{¯}{\pi }_1\|A_2\underset{¯}{\pi }_2\|\mathrm{}A_k\underset{¯}{\pi }_k`$ $``$ $``$ $`\underset{¯}{\sigma }`$ such that $`F\underset{¯}{\sigma }=G`$, Ax-iv) $`(A_1+A_2+\mathrm{}A_k)\underset{¯}{\pi }=(A_1+A_2+\mathrm{}A_k)\underset{¯}{\tau }`$ $`A_1\underset{¯}{\pi }=A_1\underset{¯}{\tau },A_2\underset{¯}{\pi }=A_2\underset{¯}{\tau },\mathrm{}A_k\underset{¯}{\pi }=A_k\underset{¯}{\tau }.`$ Finally let us consider a combination of shifting and merging which provides the convenient tool to formalise our ideas about scheduling of atoms in resolvents. As outlined in previous section, at any step of derivation, atoms coming from the body of the applied clause are assigned new priority values, while priorities of old atoms are left unchanged. This way, new atoms are positioned (i.e. scheduled) with respect to the old ones. In general, the positioning of atoms from a p-goal $`B`$, with respect to the atoms of another p-goal $`F`$, can be described through a composition of shifting and merging. Indeed, consider an expression like $`F+B\underset{¯}{\pi }`$. The effect of the shifting $`\underset{¯}{\pi }`$ is twofold. First of all, possible conflicts of priority values between$`F`$and$`B`$can be removed, so that the merging $`F+B\underset{¯}{\pi }`$ is correctly performed. At the same time, yet more important, $`\underset{¯}{\pi }`$ allows us to establish the positions which atoms from$`B`$go to occupy. Since priorities are represented by rational values, it is evident that all possible allocations of atoms from $`B`$, with respect to those in $`F`$, can be described through suitable choices of $`\underset{¯}{\pi }`$. ### 3.3 Priority SLD Derivations Now, we are ready to frame well known Logic Programming concepts, as the ones of resolvent and SLD derivation, in terms of priority atoms, goals and scheduling. We start with the following Definition 3.4, which formalises the idea of priority derivation step. Given a p-goal $`a\|F`$, in agreement with our concept of scheduling the atom $`a`$ with minimum priority is always rewritten and atoms coming from the body of the applied clause are positioned with respect to old ones to form the new resolvent. The positioning is obtained through a combination of shifting and merging, as discussed at the end of the previous Section 3.2. With reference to Definition 3.4, the body$`B`$of the applied clause is first shifted by $`\pi `$ and then merged with $`F`$, i.e. with the initial p-goal $`a\|F`$ minus the rewritten atom. ###### Definition 3.4 (priority derivation step). Consider a p-goal$`G=a\|F`$ and a clause $`c=(htB)`$. Let: \- $`\xi `$ be a renaming such that $`var(G)var(c\xi )=\mathrm{}`$, \- $`\theta `$ be an idempotent and relevant mgu of $`a`$ and $`(ht)\xi `$, -$`\underset{¯}{\pi }`$ be a shifting such that$`F`$and $`B\underset{¯}{\pi }`$ have no common priority value. We say that$`R`$is a resolvent of$`G`$and$`c`$by $`\xi `$, $`\theta `$ and $`\underset{¯}{\pi }`$, if: $`R=(F+B\xi \underset{¯}{\pi })\theta `$. The transformation from $`a\|F`$ to $`(F+B\xi \underset{¯}{\pi })\theta `$ will be called a priority derivation step. It is denoted by: $`a\|F\stackrel{𝑐}{}(F+B\xi \underset{¯}{\pi })\theta `$. The notation$`G\stackrel{c\xi ,\theta }{}R`$ will be used to represent a derivation step by $`\theta `$ and $`\xi `$, where the shifting $`\pi `$ is not pointed out. Analogously, we will write$`G\stackrel{c\xi }{}R`$ to represent a derivation step by the renaming $`\xi `$ without specifying the mgu $`\theta `$. By$`G\stackrel{𝑐}{}R`$ we denote a derivation step which generically produces$`R`$as a resolvent of$`G`$and $`c`$. Iterating the process of computing resolvents, we obtain a priority SLD derivation, that is a sequence of priority derivation steps as formalised by the following definition. ###### Definition 3.5 (priority SLD derivation). Let P be a program and $`G_o`$ a p-goal. A priority SLD derivation of $`G_o`$ in P is a possibly infinite sequence of priority derivation steps $`G_o\stackrel{c_o\xi _o,\theta _o}{}G_1\mathrm{}G_k\stackrel{c_k\xi _k,\theta _k}{}G_{k+1}\mathrm{}`$ where, for any $`j0`$, i) $`c_j`$ is a clause in P, ii) $`var(c_j\xi _j)(var(G_o)var(c_o\xi _o)\mathrm{}var(c_{j1}\xi _{j1}))=\mathrm{}.`$ Given a finite priority SLD derivation (p-SLD derivation in the following) of the form: $`G_o\stackrel{c_o\xi _o,\theta _o}{}G_1\mathrm{}G_h\stackrel{c_h\xi _h,\theta _h}{}G,`$ the sequence $`M=c_1,c_2,\mathrm{}c_h`$ of applied clauses will be called template. The whole derivation will be denoted by $`G_o\stackrel{M,\theta }{}G`$, where $`\theta =\theta _1\theta _2\mathrm{}\theta _h`$, or simply $`G_o\stackrel{𝑀}{}G`$, if the substitution $`\theta `$ does not need to be pointed out. We use the notation $`G_o\stackrel{𝑀}{}`$, when there is not interest in specifying the final resolvent. Given a template $`M`$, the notation $`\mathrm{\#}M`$ will indicate the number of clauses in $`M`$. In many cases, we will consider concatenation of templates, which is denoted by a vertical bar “$`\|`$”. It is intuitive that, given a derivation, any subset of atoms in the current resolvent derives from other specific atoms in preceding resolvents. As it will be clear in the sequel, this idea plays an important role in the development of this paper. Thus, it is convenient to give some formal definitions. Precisely, let us consider a p-SLD derivation of the form $`Dr=(F+G\stackrel{𝐻}{}Q)`$. The following two intuitive concepts will be characterised: 1. the sub-resolvent of$`F\text{}`$in $`Dr`$, i.e. the subset of p-atoms in $`Q`$ which derive from the subgoal$`F`$(denoted by $`Q/F`$), 2. the sub-template of$`F\text{}`$in $`Dr`$, i.e. the sequence of clauses which are applied to p-atoms of$`F`$and p-atoms derived from $`F`$, extracted in the order from the template $`H`$ (denoted by $`H/F`$). ###### Definition 3.6 (sub-resolvents and sub-templates). 1. Given a derivation step of the following form, where $`c=(htB)`$: $`a\|(F+G)\stackrel{𝑐}{}(Q=((F+G)+B\xi \underset{¯}{\pi })\alpha )`$, (1) let us define sub-resolvents and sub-templates in (1) as follows: $`Q/a=B\xi \underset{¯}{\pi }\alpha ,Q/F=F\alpha ,Q/(a\|F)=Q/a+Q/F`$ $`c/a=c,c/F=\mathrm{},c/(a\|F)=c.`$ 2. Given a derivation of the form: $`F+G\stackrel{𝑐}{}Q\stackrel{𝐾}{}R,`$ (2) let us recursively define sub-resolvents and sub-templates in (2) as follows: $`R/F=R/(Q/F)`$, $`(c\|K)/F=(c/F)\|(K/(Q/F))`$. It is worth noting that the notation relative to sub-templates and sub-resolvents can be ambiguous. Indeed consider: $`G+F\stackrel{𝐷}{}Q`$(3) $`G+F^{}\stackrel{𝐷}{}Q^{}`$. (4) It is possible that $`D/G`$ with respect to (3) is different from $`D/G`$ with respect to (4). In the following of the paper, when such a kind of ambiguity will possibly arise, we exploit a refined notation of evident meaning, like $`D/^3/G`$ and $`D/^4/G`$. As an example, let us consider $`G=a`$, $`F=b`$, $`F^{}=d`$ and $`D=c`$ such that $`G+F=a\|b\stackrel{𝑐}{}Q`$(3b) $`G+F^{}=d\|a\stackrel{𝑐}{}Q^{}.`$ (4b) Then, $`D/^{3b}/G=c`$ and $`D/^{4b}/G=empty`$. ### 3.4 Congruent lowering of derivation steps This section introduces some important ideas. Precisely, the concepts of specialisation, lowering, and finally congruent lowering are defined and analysed. Congruent lowering is basic for the characterisation of the general concept of scheduling rule, as well as of the class of specialisation independent scheduling rules (see Section 4) to which the results about redundancy elimination tolerance of Section 6 refer. Substitutions and renamings are basic concepts in Logic Programming. In agreement with usual terminology, if a substitution is applied to a goal, an instance is obtained, while, if a renaming is used, a variant of the original goal is produced. Goals which are equal up to renamings are in essence equivalent goals. Practically all the results of Logic Programming are insensible to renamings. An instance may be considered as a specialised version of the original goal, while any goal is more general with respect to its instances. The above concepts are easily adjusted in the frame of priority goals. Intuitively, the application of a renaming/substitution corresponds to the application of a renaming/substitution together with a shifting. Actually, as it will be clear in the following, we are interested in an idea of specialisation of a given p-goal which extends the traditional concept of instantiation. In essence, we will consider couples of p-goals such that the second one is obtained from the first one by performing in the order: \- the application of a generic substitution $`\lambda `$ and a shifting $`\sigma `$, \- the embedding in a generic context $`X`$ of other p-atoms. ###### Definition 3.7 (specialisation). A p-goal$`F`$is a specialisation of a p-goal $`a\|K`$ by $`X`$, if a shifting $`\sigma `$ and a substitution $`\lambda `$ exist such that $`F=a\lambda \underset{¯}{\sigma }\|(K\lambda \underset{¯}{\sigma }+X)`$. It is worth noting that our idea of specialisation is essentially symmetric to the concept of subsumption by an instance (see \[Bol, Apt and Klop, 1991\]). A goal$`G`$subsumes (as list) a goal$`F`$by an instance, if a substitution $`\lambda `$ exists such that $`G\lambda _LF`$. Indeed, considering that any shifting preserves the order of the atoms, it is evident that, if$`F`$is a specialisation of $`a\|K`$ by $`X`$, i.e.$`F=a\lambda \underset{¯}{\sigma }\|(K\lambda \underset{¯}{\sigma }+X)`$, then $`a\|K`$ subsumes (as list)$`F`$by the instance $`(a\|K)\lambda `$. The term “lifting” is used in Logic Programming to express that a derivation step (or a whole derivation) which is possible from a goal $`A\lambda `$ is repeated starting from the more general goal $`A`$. Analogously, we use the term lifting to mean that a derivation step (or a whole derivation) which is possible from a specialisation of $`a\|K`$, i.e. from a p-goal $`a\lambda \underset{¯}{\sigma }\|(K\lambda \underset{¯}{\sigma }+X)`$, is repeated starting from $`a\|K`$. In the sequel of the paper, we will use the dual concept of “lowering”. In other words, the term lowering will mean that a derivation step (or a whole derivation) from a p-goal $`a\|K`$ is repeated, when possible, starting from a specialisation $`a\lambda \underset{¯}{\sigma }\|(K\lambda \underset{¯}{\sigma }+X)`$ of $`a\|K`$. Then, let us give the following definition which refers to single derivation steps. ###### Definition 3.8 (lowering of derivation steps). Let us consider two priority derivation steps of the type$`G\stackrel{𝑐}{}`$ and $`F\stackrel{𝑐}{}`$. We will say that the second step is a lowering of the first one by $`X`$, if the p-goal$`F`$is a specialisation of$`G`$by $`X`$. Let us consider two derivation steps (Ds1) and (Ds2), such that (Ds2) is a lowering of (Ds1) by $`X`$, and let $`c=(htB)`$. By definition of derivation step, they have the following form: $`a\|K\stackrel{𝑐}{}(K+B\xi ^{}\underset{¯}{\theta }^{})\alpha ^{}`$(Ds1) $`a\lambda \underset{¯}{\sigma }\|(K\lambda \underset{¯}{\sigma }+X)\stackrel{𝑐}{}(X+K\lambda \underset{¯}{\sigma }+B\xi ^{\prime \prime }\underset{¯}{\theta }^{\prime \prime })\alpha ^{\prime \prime }.`$ (Ds2) The definition of lowering of derivation steps does not impose any similarity in the way priority values are handled in couples of derivation steps like (Ds1) and (Ds2). In particular, no analogy is required about the positions new atoms go to occupy with respect to old ones in the resolvents produced by (Ds1) and (Ds2). Indeed the shifting $`\underset{¯}{\theta }^{}`$ and $`\underset{¯}{\theta }^{\prime \prime }`$ are completely independent, so that the positions of atoms of $`B\xi ^{\prime \prime }\underset{¯}{\theta }^{\prime \prime }`$, with respect to atoms of $`K\lambda \underset{¯}{\sigma }`$, will be in general different from the positions occupied by atoms of $`B\xi ^{}\underset{¯}{\theta }^{}`$ with respect to atoms of $`K`$. Nevertheless, in the rest of the paper special importance will be given to derivation step lowering such that the positioning of new atoms, with respect to the old ones in $`K`$ and $`K\lambda \underset{¯}{\sigma }`$, is maintained passing from (Ds1) to (Ds2). In such hypothesis, we will say that the lowering is a congruent lowering. As an elementary example, let us consider a clause like $`c=ab_1\|b_2`$ and the following derivation steps, such that (2) is a lowering of (1) by $`x_1\|x_2`$: $`a\|k_1\|k_2\stackrel{𝑐}{}b_1^{}\underset{¯}{\theta }^{}\|k_1\|b_2\underset{¯}{\theta }^{}\|k_2`$(1) $`a\|x_1\|k_1\|x_2\|k_2\stackrel{𝑐}{}x_1\|b_1\underset{¯}{\theta }^{\prime \prime }\|k_1\|x_2\|b_2\underset{¯}{\theta }^{\prime \prime }\|k_2`$(2) In (1) and (2) the relative positions of atoms $`b_1`$ and $`b_2`$ with respect to $`k_1`$ and $`k_2`$ are the same, then (2) is a congruent lowering of (1). Now, let us consider the following other derivation step (3): $`a\|x_1\|k_1\|x_2\|k_2\stackrel{𝑐}{}x_1\|k_1\|b_1\underset{¯}{\tau }\|x_2\|b_2\underset{¯}{\tau }\|k_2`$(3) Also (3) is a lowering of (1) by $`x_1\|x_2`$. However, in this case the positioning of atoms $`b_1`$ and $`b_2`$ with respect to $`k_1`$ and $`k_2`$ is not maintained passing from (1) to (3), so that (3) is not a congruent lowering of (1). Variable substitutions are not considered in the above examples. Indeed, in agreement with the following formal Definition 3.9, they are not really influent for a lowering to be congruent or not. ###### Definition 3.9 (congruent lowering). Let us consider two derivation steps of the form (Ds1) and (Ds2) above, i.e. two derivation steps such that the second one is a lowering of the first one by $`X`$. We will say that step (Ds2) is a congruent lowering of step (Ds1) by $`X`$ if a shifting $`\rho `$ exists with: $`K\underset{¯}{\rho }=K\underset{¯}{\sigma }`$ and $`B\underset{¯}{\theta }^{}\underset{¯}{\rho }=B\underset{¯}{\theta }^{\prime \prime }`$. (c1) It is apparent that the desired analogy, in positioning new atoms in the two derivation steps (Ds1) and (Ds2), is imposed by means of condition (c1) above in Definition 3.9. Indeed, condition (c1) says that the shifting $`\rho `$ creates a correspondence between atoms of $`K+B\underset{¯}{\theta }^{}`$ and atoms of $`K\underset{¯}{\sigma }+B\underset{¯}{\theta }^{\prime \prime }`$, such that old atoms are mapped in old atoms (see $`K\underset{¯}{\rho }=K\underset{¯}{\sigma }`$) and new atoms in new ones (see $`B\underset{¯}{\theta }^{}\underset{¯}{\rho }=B\underset{¯}{\theta }^{\prime \prime }`$). Since any shifting maintains atom precedence, it is intuitive that congruent allocation of new atoms is imposed. More specifically, let us consider the generic atom $`b`$ of$`B`$and assume: $`K+B\underset{¯}{\theta }^{}=M^{}\|b\underset{¯}{\theta }^{}\|N^{}`$, $`K\underset{¯}{\sigma }+B\underset{¯}{\theta }^{\prime \prime }=M^{\prime \prime }\|b\underset{¯}{\theta }^{\prime \prime }\|N^{\prime \prime }`$. It is immediate to verify that<sup>\[</sup><sup>1</sup><sup>1</sup>1The notation $`\mathrm{"}K\underset{¯}{\sigma }+B\underset{¯}{\theta }\mathrm{"}=^{(c1)}K\underset{¯}{\rho }+B\underset{¯}{\theta }^{}\underset{¯}{\rho }\mathrm{"}`$ expresses that the formula (c1) must be used to establish the equality. Similar advising will be used frequently in the sequel.<sup>\]</sup>: $`M^{\prime \prime }\|b\underset{¯}{\theta }^{\prime \prime }\|N^{\prime \prime }=K\underset{¯}{\sigma }+B\underset{¯}{\theta }^{\prime \prime }=^{(c1)}K\underset{¯}{\rho }+B\underset{¯}{\theta }^{}\underset{¯}{\rho }=(K+B\underset{¯}{\theta }^{})\underset{¯}{\rho }=`$ $`=(M^{}\|b\underset{¯}{\theta }^{}\|N^{})\underset{¯}{\rho }=M^{}\underset{¯}{\rho }\|b\underset{¯}{\theta }^{}\underset{¯}{\rho }\|N^{}\underset{¯}{\rho }.`$ Now, by $`B\underset{¯}{\theta }^{}\underset{¯}{\rho }=B\underset{¯}{\theta }^{\prime \prime }`$ in (c1) and Ax-iv in Property 3.2, we have that $`b\underset{¯}{\theta }^{}\underset{¯}{\rho }=b\underset{¯}{\theta }^{\prime \prime }`$. Then, by Property 3.1-ii) it is $`M^{}\underset{¯}{\rho }=M^{\prime \prime }`$, and then also $`\mathrm{\#}M^{}=\mathrm{\#}M^{}\underset{¯}{\rho }=\mathrm{\#}M^{\prime \prime }=n`$, for $`n`$ positive integer. In essence, considered the generic atom $`b`$ of $`B`$, it is found in the $`(n+1)^{th}`$ position in $`K+B\underset{¯}{\theta }^{}`$ as well as in $`K\underset{¯}{\sigma }+B\underset{¯}{\theta }^{\prime \prime }`$. In other words, new atoms from$`B`$are positioned in (Ds1) with respect to old ones (i.e. atoms of $`K`$) exactly as it happens in (Ds2) with respect to $`K\underset{¯}{\sigma }`$. It is evident that the presence of various substitutions in (Ds1) and (Ds2) does not interfere with the above positional considerations. ###### Example 3.10 (lowering and congruent lowering). Let us consider a clause of the form$`c=(aq[1])`$ and the two following derivation steps: $`a[2]\|\{b[\mathrm{𝟑}]\}\stackrel{𝑐}{}\{b[\mathrm{𝟑}],q[\underset{¯}{10}]\}`$, (1) $`a[9]\|\{b[12],b[13],d[15]\}\stackrel{𝑐}{}\{b[12],q[12.5],b[13],d[15]\}`$. (2) In step (1), old atoms are pointed out in bold and new ones are underlined. 1. In agreement with Definition 3.9, step (2) is a lowering of (1) by $`X=\{b[13],d[15]\}`$, with $`K\underset{¯}{\sigma }=\{b[12]\}`$. Pointing out old and new atoms, derivation step (2) can be written as follows: $`a[9]\|\{b[\mathrm{𝟏𝟐}],b[13],d[15]\}\stackrel{𝑐}{}\{b[\mathrm{𝟏𝟐}],q[\underset{¯}{12.5}],b[13],d[15]\}`$. It is evident that (2) is a congruent lowering of (1) by $`X`$, with any shifting $`\rho `$ such that $`\underset{¯}{\rho }\{3/12,10/12.5\}`$. 2. Step (2) is a lowering of (1) also by $`X^{}=\{b[12],d[15]\}`$, with $`K\underset{¯}{\sigma }^{}=\{b[13]\}`$. However (2) is not a congruent lowering of (1) by $`X^{}`$. In fact, in agreement with this second viewpoint, derivation step (2) can be written as follows: $`a[9]\|\{b[12],b[\mathrm{𝟏𝟑}],d[15]\}\stackrel{𝑐}{}\{b[12],q[\underset{¯}{12.5}],b[\mathrm{𝟏𝟑}],d[15]\}`$. As a consequence, for step (2) being a congruent lowering of step (1) by $`X^{}`$, a shifting $`\underset{¯}{\rho }^{}`$ might exist such that $`\underset{¯}{\rho }^{}\{3/13,10/12.5\}`$, which is not an increasing function. $`\mathrm{}`$ We close this section considering a couple of p-goals$`F`$and$`G`$such that they are specialisations of each other, i.e.$`F`$is a specialisation of$`G`$by a subgoal $`X`$ and$`G`$is a specialisation of$`F`$by $`Y`$. In this case it must be$`F=G\lambda \underset{¯}{\sigma }+X`$ and$`G=F\tau \underset{¯}{\rho }+Y`$, which yields: $`G=F\tau \underset{¯}{\rho }+Y=(G\lambda \underset{¯}{\sigma }+X)\tau \underset{¯}{\rho }+Y=G\lambda \tau \underset{¯}{\sigma }\underset{¯}{\rho }+X\tau \underset{¯}{\rho }+Y`$. As a consequence $`\lambda `$ must be a renaming for$`G`$and $`X=Y=\mathrm{}`$ must hold, which means that$`F=G\lambda \underset{¯}{\sigma }`$ where $`\lambda `$ is a renaming. It is evident that the relation “$`F=G\lambda \underset{¯}{\sigma }`$, for a renaming $`\lambda `$ and a shifting $`\underset{¯}{\sigma }`$” can be seen as the translation of the usual notion of “$`F`$ being variant of $`G`$” in the frame of p-SLD resolution. In this sense, we will usually say that$`F`$is a p-variant of a $`G`$, to mean that$`F`$and$`G`$are specialisations of each other. Analogously, two derivation steps may be lowerings of each other, as well as congruent lowerings of each other. Two derivation steps $`Ds_1`$ and $`Ds_2`$ are lowerings of each other if the initial goals are p-variants and the same clause is applied, i.e. it is $`Ds_1=(A\stackrel{𝑐}{})`$ and $`Ds_2=(A\lambda \underset{¯}{\sigma }\stackrel{𝑐}{})`$, where $`\lambda `$ is a renaming. Two derivation steps are congruent lowerings of each other if they have the form: $`Ds_1=a\|K\stackrel{𝑐}{}(K+B\xi ^{}\underset{¯}{\theta }^{})\alpha ^{}`$ and $`Ds_2=(a\|K)\lambda \underset{¯}{\sigma }\stackrel{𝑐}{}(K\lambda \underset{¯}{\sigma }+B\xi ^{\prime \prime }\underset{¯}{\theta }^{\prime \prime })\alpha ^{\prime \prime }`$, where $`c=(htB)`$, $`\lambda `$ is a renaming, and the equalities $`K\underset{¯}{\rho }=K\underset{¯}{\sigma }`$ and $`B\underset{¯}{\theta }^{}\underset{¯}{\rho }=B\underset{¯}{\theta }^{\prime \prime }`$ hold for a shifting $`\underset{¯}{\rho }`$. It is worth noting that by the preceding argument if two derivation steps are lowerings of each other the contexts must be empty. ### 3.5 State priority scheduling rules Now, we use the notion of being congruent lowerings of each other to define the ideas of determinism and completeness of a set of derivation steps. Both concepts are basic for the definition of state priority scheduling rules. ###### Definition 3.11 (determinism). A set$`S`$of priority derivation steps is deterministic if, for each couple of derivation steps $`Ds_1`$ and $`Ds_2`$ in$`S,`$ the following implication holds: $`Ds_1`$ and $`Ds_2`$ are lowerings of each other $``$ $`Ds_1`$ and $`Ds_2`$ are congruent lowerings of each other. In other words, the definition of determinism imposes that two derivation steps, which apply the same clause to p-variant initial goals, give place to congruent allocations of new atoms. Now let us give the definition of completeness of a set of derivation steps. ###### Definition 3.12 (completeness). A set S of priority derivation steps is complete, if the following assertions hold: i) $``$ $`Ds`$ derivation step of the type$`G\stackrel{𝑐}{},`$ $``$ $``$ $`Ds^{}`$ of the type$`G\stackrel{𝑐}{}`$ with $`Ds^{}S`$, ii) $`Ds^{},Ds`$ derivation steps with $`DsS`$, $`Ds^{}`$ and $`Ds`$ are congruent lowerings of each other $``$ $`Ds^{}S`$. Assertion i) of the above definition states that, if a clause$`c`$is applicable to a p-goal $`G`$, i.e. a derivation step exists of the type$`G\stackrel{𝑐}{}`$, the application of the clause$`c`$to$`G`$is indeed possible in any complete set of derivation steps. Assertion ii) assures that$`S`$is closed with respect to being congruent lowerings of each other. In other words, let $`Ds^{}=(G\stackrel{𝑐}{}Q)S`$ be a derivation step, then every other $`Ds^{\prime \prime }=(F\stackrel{𝑐}{}R)`$ must belong to$`S,`$ if$`F`$is a p-variant of$`G`$and new atoms are allocated in$`R`$as it is done in $`Q`$. Now, the formal definition of state priority scheduling rules can be easily given, by combining the properties of determinism and completeness. ###### Definition 3.13 (state priority scheduling rules). A state priority scheduling rule is a complete and deterministic set of priority derivation steps. It can be easily verified that the leftmost selection rule, adopted by the Prolog execution mechanism, is a state priority scheduling rule. The very nature of a state scheduling rule is characterised by the following Definition 3.14. Indeed, the definition simply says that a p-SLD derivation is via a state scheduling rule$`S`$if all derivation steps are admitted in the rule$`S,`$ i.e. they all belong to the set of derivation steps which$`S`$is constituted by. ###### Definition 3.14 (derivations via S). 1. Given a set$`S`$of derivation steps, the notation $`\mathrm{\Delta }(S)`$ represents the whole of p-SLD derivations which are composed of derivation steps in$`S.`$ 2. Given a state scheduling rule$`S,`$ the set $`\mathrm{\Delta }(S)`$ is the set of p-SLD derivations via$`S.`$ In the sequel of the paper we only consider state priority scheduling rules, which therefore will be called just scheduling rules. The following notations will be used frequently. Given a set$`S`$of derivation steps, a clause$`c`$and a template $`M`$, we will denote by $`G\stackrel{S,c}{}R`$and$`G\stackrel{S,M}{}R`$ the fact that the derivation step $`(G\stackrel{𝑐}{}R)S`$and the p-SLD derivation $`(G\stackrel{𝑀}{}R)\mathrm{\Delta }(S)`$, respectively. In the case that the exploited logic program must be pointed out, a notation like $`(G\stackrel{S,M.P}{}R)`$ will be used to specify that the derivation is via$`S`$in the program P, i.e. every clause of the template $`M`$ belongs to P. The notion of p-SLD tree via$`S`$could be characterised in complete analogy with the usual one of SLD tree. Let us close this section with a property, which can be easily shown on the basis of completeness and will be used several times in the sequel. Property 3.3 asserts that if a clause$`c`$can be applied to a p-goal $`a\gamma \|G`$, every complete set of derivation steps allows$`c`$to be applied to any p-goal of the form $`a\|F`$. Since the atom $`a`$ is more general than $`a\gamma `$, the property may also be interpreted as a sort of lifting of derivation steps. However, the subgoals$`G`$and$`F`$are left unrelated at all. The evident explication is that they have no active role in rewriting operations. Moreover, the property recalls that new variables can be always chosen so that conflicts are avoided with arbitrary pre-established sets of variables. The formal proof of this rather intuitive property can be found in Appendix A. In the statement of Property 3.3 and in the sequel of the paper, given a p-SLD derivation $`Dr`$, the notation $`nvar(Dr)`$ will represent the set of standardisation apart variables which are introduced during the derivation $`Dr`$. In the case of a single derivation step $`Ds=(A\stackrel{c\xi }{})`$, it is $`nvar(Ds)=var(c\xi )`$. ###### Property 3.3 Let$`S`$be a complete set of derivation steps. Given two p-goals $`a\gamma \underset{¯}{\tau }\|G`$ and $`a\|F`$, let us fix arbitrarily a finite set $`V`$ of variables. The following implication holds: $`Ds`$ derivation step of the type $`a\gamma \underset{¯}{\tau }\|G\stackrel{𝑐}{}`$ $``$ $`Ds^{}`$ of the type $`a\|F\stackrel{𝑐}{},`$ with $`Ds^{}S`$and $`nvar(Ds^{})V=\mathrm{}`$. ## 4 Specialization independent scheduling rules Now, we will exploit the notion of congruent lowering in order to introduce the concept of specialisation independence. This concept will be used to characterise the class of scheduling rules that are the main object of the paper (specialisation independent scheduling rules). In fact, all our results for termination and loop check completeness preserving will refer to such a class of scheduling rules. In Section 5, a second characterisation of the same class is given which has an operational nature and is surprisingly different in appearance. The definition of specialisation independence enforces the idea of determinism. Indeed, in agreement with the Definition 4.1 below, every lowering is required to be a congruent lowering. In other words, the congruence in the allocation of new atoms must hold any time the initial goals of two derivation steps are related by specialisation and the same clause is used. This can be interpreted saying that the positioning of new atoms with respect to old ones is independent of goal specialisation, which means independent of goal instantiation as well as of the addition of a group $`X`$ of other atoms. ###### Definition 4.1 (specialisation independence). A set$`S`$of priority derivation steps is specialisation independent if, for every couple of steps $`Ds_1`$ and $`Ds_2`$ in$`S,`$ the following implication holds: $`Ds_2`$ is a lowering of $`Ds_1`$ by $`X`$ $`Ds_2`$ is a congruent lowering of $`Ds_1`$ by $`X`$. ###### Definition 4.2 (specialisation independent scheduling rules). A specialisation independent scheduling rule is a complete and specialisation independent set of priority derivation steps. In the next two sections, we provide some results about p-SLD derivations via specialisation independent scheduling rules. The results will be frequently exploited in the sequel. ### 4.1 Derivation lowering In this section we give results which relate resolvents coming from a couple of derivation steps in the congruent lowering relationship. Then, by Lemma 4.4, the analysis is extended to couples of whole derivations, developed via specialisation independent scheduling rules. We start by presenting a preliminary statement (Property 4.1) which holds for every couple of derivation steps that are in the lowering relationship. In reference to derivation steps (1) and (2) below, the preliminary property says that, if we abstract from atom positioning and ignore the additional subgoal $`X`$, the resolvent of (2) is an instance of the resolvent of (1). Property 4.1 can be shown following the line exploited for proving the Variant Lemma (see \[Apt, 1990\]), which is done in Appendix A for the sake of completeness of the paper. ###### Property 4.1 Let $`c=(htB)`$ be a clause. Let us consider two derivation steps like (1) and (2), where (2) is a lowering of (1) by $`X`$. The following implication holds: $`a\|K\stackrel{c\xi ^{}}{}(K+B\xi ^{}\underset{¯}{\theta }^{})\mu ^{}`$, (1) $`a\tau \underset{¯}{\sigma }\|(K\tau \underset{¯}{\sigma }+X)\stackrel{c\xi ^{\prime \prime }}{}(K\tau \underset{¯}{\sigma }+B\xi ^{\prime \prime }\underset{¯}{\theta }^{\prime \prime }+X)\mu ^{\prime \prime }`$ (2) $`\delta `$ such that $`K\tau \mu ^{\prime \prime }=K\mu ^{}\delta `$ and $`B\xi ^{\prime \prime }\mu ^{\prime \prime }=B\xi ^{}\mu ^{}\delta `$, where $`\delta `$ is a renaming, if $`\tau `$ is a renaming. Property 4.2 completes Property 4.1, taking into account the preservation of atom scheduling in the case of congruent lowering. It states that, if we ignore the additional subgoal $`X`$, resolvents are preserved up to a substitution and a shifting. In reference to derivation steps (1) and (2) below, this means that, apart from $`R/X`$, the resolvent$`R`$in (2) is an instance of $`Q`$ such that also atom scheduling is maintained. ###### Property 4.2 Let $`c=(htB)`$ be a clause. Let us consider two derivation steps of the type (1) and (2), such that the second one is a congruent lowering of the first one by $`X`$: $`a\|K\stackrel{𝑐}{}Q`$, (1) $`a\tau \underset{¯}{\pi }\|(K\tau \underset{¯}{\pi }+X)\stackrel{𝑐}{}R`$. (2) The following assertion holds: $`\delta ,\underset{¯}{\rho }`$ such that $`R/((a\|K)\tau \underset{¯}{\pi })=Q\delta \underset{¯}{\rho }`$, where $`\delta `$ is a renaming if $`\tau `$ is a renaming. ###### Proof 4.3. Let $`c=(htB)`$, so that $`Q`$ and$`R`$may be written as follows: $`Q=(K+B\xi ^{}\underset{¯}{\theta }^{})\mu ^{}`$, $`R=(K\tau \underset{¯}{\pi }+X+B\xi ^{\prime \prime }\underset{¯}{\theta }^{\prime \prime })\mu ^{\prime \prime }`$. Since step (2) is a congruent lowering of (1) by $`X`$, a shifting $`\rho `$ exists such that: $`K\underset{¯}{\rho }=K\underset{¯}{\pi },B\underset{¯}{\theta }^{}\underset{¯}{\rho }=B\underset{¯}{\theta }^{\prime \prime }`$. (3) By definition of sub-resolvent and (3), we have: $`R/((a\|K)\tau \underset{¯}{\pi })=B\xi ^{\prime \prime }\underset{¯}{\theta }^{\prime \prime }\mu ^{\prime \prime }+K\tau \underset{¯}{\pi }\mu ^{\prime \prime }=^{(3)}B\xi ^{\prime \prime }\mu ^{\prime \prime }\underset{¯}{\theta }^{}\underset{¯}{\rho }+K\tau \mu ^{\prime \prime }\underset{¯}{\rho }.`$ (4) Now, we apply Property 4.1 to (1) and (2), deriving that a substitution $`\delta `$ exists such that: $`K\tau \mu ^{\prime \prime }=K\mu ^{}\delta `$ and $`B\xi ^{\prime \prime }\mu ^{\prime \prime }=B\xi ^{}\mu ^{}\delta ,`$ (5) where $`\delta `$ is a renaming if $`\tau `$ is a renaming. As a consequence, we have that: $`R/((a\|K)\tau \underset{¯}{\pi })=^{(4)}B\xi ^{\prime \prime }\mu ^{\prime \prime }\underset{¯}{\theta }^{}\underset{¯}{\rho }+K\tau \underset{¯}{\rho }\mu ^{\prime \prime }=^{(5)}B\xi ^{}\mu ^{}\delta \underset{¯}{\theta }^{}\underset{¯}{\rho }+K\mu ^{}\delta \underset{¯}{\rho }=Q\delta \underset{¯}{\rho },`$ where $`\delta `$ is a renaming if $`\tau `$ is a renaming. The following Lemma 4.4 may be seen as the extension of Property 4.2 to whole derivations, provided that the used scheduling rule is specialisation independent. Note that, given a derivation like (1) in the statement below, if a derivation like (2) exists, it can be considered as a lowering of (1). Indeed, the initial p-goal $`X+G\gamma \underset{¯}{\tau }`$ is a specialisation of$`G`$by $`X`$, and the sequence $`E`$ of clauses is applied in the same order to atoms deriving from $`G\gamma \underset{¯}{\tau }`$ in derivation (2). In this sense we will regard Lemma 4.4 as a “specialisation independent lowering lemma”. ###### Lemma 4.4 (specialisation independent lowering lemma). Let$`S`$be a specialisation independent scheduling rule and consider two p-SLD derivations like (1) and (2). The following implication holds: $`G\stackrel{S,E}{}Q`$, (1) $`G\gamma \underset{¯}{\tau }+X\stackrel{S,D}{}R`$, with $`D/(G\gamma \underset{¯}{\tau })=E`$ (2) $`\sigma ,\underset{¯}{\rho }`$ such that $`R/(G\gamma \underset{¯}{\tau })=Q\sigma \underset{¯}{\rho }`$, where $`\sigma `$ is a renaming if $`\gamma `$ is a renaming and $`D/X=\mathrm{}`$. (p1) ###### Proof 4.5. Let us first prove the thesis, apart from the fact (p1). The proof is by induction on the length of $`D`$. If $`\mathrm{\#}D`$ is equal to zero, the thesis is trivially true. Let us suppose that $`\mathrm{\#}D`$ is greater than zero. Two different cases must be considered, i.e. the first clause of $`D`$ (say $`c`$) is applied either to an atom of $`X`$ or to an atom of $`G\gamma \underset{¯}{\tau }`$. First case (The clause$`c`$is applied to an atom of $`X`$). In this case derivation (2) may be rewritten as: $`X+G\gamma \underset{¯}{\tau }\stackrel{S,c\eta ,\alpha }{}G\gamma \underset{¯}{\tau }\alpha +X^{}\stackrel{S,D^{}}{}R`$, (3) with $`D^{}/(G\gamma \underset{¯}{\tau }\alpha )=D/(G\gamma \underset{¯}{\tau })=E`$. By inductive hypothesis, applied to the tail of derivation (3) and derivation (1), we have: $`\sigma ,\underset{¯}{\rho }`$ such that $`R/(G\gamma \underset{¯}{\tau })=R/(G\gamma \underset{¯}{\tau }\alpha )=^{(ind.hyp.)}Q\sigma \underset{¯}{\rho }.`$ Second case (The clause$`c`$is applied to an atom of $`G\gamma \underset{¯}{\tau }`$). In in this case derivations (1) and (2) may be rewritten as (4) and (5), respectively: $`G\stackrel{S,c\xi }{}Y\stackrel{S,E^{}}{}Q`$ (4) $`X+G\gamma \underset{¯}{\tau }\stackrel{S,c\eta ,\alpha }{}X\alpha +Z\stackrel{S,D^{}}{}R`$ (5) with $`D^{}/Z=E^{}`$, $`c/^4/G=c/^5/G\gamma \underset{¯}{\tau }=c`$. (6) Since$`S`$is specialisation independent, the first step of (5) is a congruent lowering of the first one of (4) by $`X`$. Then, by Property 4.2, we have: $`\sigma ^{},\underset{¯}{\rho }^{}`$ such that $`Z=(X\alpha +Z)/(G\gamma \underset{¯}{\tau })=^{(Prop.\text{4.2})}Y\sigma ^{}\underset{¯}{\rho }^{}`$. (7) As a consequence, recalling the first fact in (6), the inductive hypothesis can be applied to the tails of derivations (4) and (5). Then, we have: $`\sigma ,\underset{¯}{\rho }`$ such that $`R/Z=Q\sigma \underset{¯}{\rho }`$. (8) In conclusion, we have that: $`R/(G\gamma \underset{¯}{\tau })=R/Z=Q\sigma \underset{¯}{\rho }`$. In order to show the fact (p1), i.e. $`\sigma `$ is a renaming if $`\gamma `$ is a renaming and $`D/X=\mathrm{}`$, it is sufficient to note that: \- the “first case” does not occur at all, \- the substitutions $`\sigma ^{}`$ and $`\sigma `$, mentioned in (7) and (8), are renamings. The following example shows that the hypothesis of specialisation independence is crucial for the validity of Lemma 4.4. ###### Example 4.6. Let us consider a scheduling rule$`S`$such that new atoms are positioned in the centre of the old resolvent. New atoms are positioned immediately before the centre if the length of the resolvent (the rewritten atom excluded) is odd. It is easy to recognise that Lowering Lemma 4.4 does not hold for such a rule. Indeed, let P be the following program: $`c1=p(x)q(x)[1]`$ $`c2=sp(b)[1]`$. Now, in reference to the statement of Lemma 4.4, let: $`G=s[1],p(a)[2]`$ $`G\gamma \underset{¯}{\tau }=s[1],p(a)[1.5]`$ and $`X=r[2].`$ The following are two derivations of$`G`$in P and $`(G\gamma \underset{¯}{\tau }+X)`$ in P, respectively: $`\{s[1],p(a)[2]\}\stackrel{S,c2}{}\{p(b)[1],p(a)[2]\}\stackrel{S,c1}{}(\{q(b)[1],p(a)[2]\}=Q)`$ $`\{s[1],p(a)[1.5],r[2]\}\stackrel{S,c2}{}\{p(a)[1.5],p(b)[1.7],r[2]\}\stackrel{S,c1}{}`$ $`(\{p(b)[1.7],q(a)[1.8],r[2]\}=R).`$ Thus, no $`\sigma `$ and $`\underset{¯}{\rho }`$ can exist such that: $`R/(G\gamma \underset{¯}{\tau })=\{q(a)[1.8],p(b)[1.7]\}=\{q(b)[1],p(a)[2]\}\sigma \underset{¯}{\rho }=Q\sigma \underset{¯}{\rho }.`$ Note that $`R/(G\gamma \underset{¯}{\tau })`$ and $`Q`$ are essentially different, even if they are considered as multisets abstracting from priority values. It is easy to check that the used scheduling rule is not specialisation independent, in agreement with Definition 4.2.$`\mathrm{}`$ ### 4.2 Derivation lifting and combining The following Lemma 4.7 is a result about p-SLD derivation lifting which is valid for specialisation independent scheduling rules. In reference to derivation (1) below, the lemma asserts that the sub-template of clauses, applied to the part $`G\gamma \underset{¯}{\tau }`$ of the initial p-goal $`(X+G\gamma \underset{¯}{\tau })`$ in (1), can be applied again in the order starting from the more general goal $`G`$, via the same scheduling rule. The lemma also recalls that standardisation apart variables can be chosen in order to avoid conflicts with any fixed finite set of variables. The lemma does not relate resolvents. Indeed, Lemma 4.4 can be exploited to this purpose. ###### Lemma 4.7 (specialisation independent lifting lemma). Let$`S`$be a specialisation independent scheduling rule. Given any finite set $`V`$ of variables, the following implication holds: $`X+G\gamma \underset{¯}{\tau }\stackrel{S,D}{}`$ (1) $`Dr=(G\stackrel{S,D/G\gamma \underset{¯}{\tau }}{})`$, with $`nvar(Dr)V=\mathrm{}`$. ###### Proof 4.8. The proof is by induction on the length of the template $`D`$. If $`\mathrm{\#}D`$ is zero, the assert is evident. Let us suppose that $`\mathrm{\#}D>0`$. Two cases must be considered, i.e. either the first clause in $`D`$ (say $`c`$) is applied to an atom of $`X`$ or the clause$`c`$is applied to an atom of $`G\gamma \underset{¯}{\tau }`$. First case (The clause$`c`$is applied to an atom of $`X`$). Derivation (1) may be rewritten as: $`X+G\gamma \underset{¯}{\tau }\stackrel{S,c\eta ,\beta }{}X^{}+G\gamma \underset{¯}{\tau }\beta \stackrel{S,D^{}}{}.`$ (2) By inductive hypothesis applied to the tail of (2), for any finite set $`V`$ of variables, a derivation $`Dr`$ exists such that: $`Dr=(G\stackrel{S,D^{}/G\gamma \underset{¯}{\tau }\beta }{})`$, with $`nvar(Dr)V=\mathrm{}`$. But, by construction of (2), it is $`D^{}/G\gamma \underset{¯}{\tau }\beta =D/G\gamma \underset{¯}{\tau }`$, so that the thesis is verified. Second case (The clause$`c`$is applied to an atom of $`G\gamma \underset{¯}{\tau }`$). Derivation (1) may be rewritten as follows: $`X+G\gamma \underset{¯}{\tau }\stackrel{S,c\eta ,\beta }{}X\beta +G^{}\stackrel{S,D^{}}{}`$, (3) where $`c\|(D^{}/G^{})=D/G\gamma \underset{¯}{\tau }`$. (4) Let$`G=a\|Z`$, that is $`X+G\gamma \underset{¯}{\tau }=a\gamma \underset{¯}{\tau }\|(X+Z\gamma \underset{¯}{\tau })`$. By (3) and Property 3.3, we can assert that a derivation step exists like: $`Ds^{}=((G=a\|Z)\stackrel{S,c}{}R^{}),`$ (6a) with $`nvar(Ds^{})V=\mathrm{}`$. (6b) Since, by hypothesis$`S`$is specialisation independent, the first step of derivation (3) is a congruent lowering of step (6a) by $`X`$. As a consequence, by Property 4.2, a substitution $`\pi `$’ and a shifting $`\underset{¯}{\rho }`$’ exist with: $`G^{}=(X\beta +G^{})/(G\gamma \underset{¯}{\tau })=R^{}\pi ^{}\underset{¯}{\rho }^{}.`$ (7) Then, by inductive hypothesis applied to the tail of (3), we may assert that, a derivation $`Dr^{\prime \prime }`$ exists: $`Dr^{\prime \prime }=(R^{}\stackrel{S,D^{}/G^{}}{})`$ (8a) with $`nvar(Dr^{\prime \prime })(nvar(Ds^{})var(G)V)=\mathrm{}`$. (8b) So, derivation (8a) is standardised apart with respect to (6a). Since$`S`$is a state scheduling rule, (6a) and (8a) can be combined in order to give place to an unique derivation $`Dr`$ such that: $`Dr=(G\stackrel{𝑐}{}R^{}\stackrel{D^{}/G^{}}{})\mathrm{\Delta }(S)`$, where, by (6b) and (8b), we have also that: $`nvar(Dr)V=(nvar(Ds^{})nvar(Dr^{\prime \prime }))V=\mathrm{}.`$ By (4), the thesis is proven. It is worth noting that Lowering Lemma 4.4 and Lifting Lemma 4.7 consider couples of p-goals in a specialisation relationship, i.e. p-goals of the form$`G`$and $`(G\gamma \underset{¯}{\tau }+X)`$. The distinctive point is that a group $`X`$ of additional atoms may be present in the second p-goal, besides the instantiation of$`G`$by $`\gamma `$. The correspondence is obvious with the fact that Definition 4.1 requires that positioning of new atoms is independent of goal specialisation. As it will be clear in the following, this kind of independence is basic in order to assure tolerance to redundancy elimination. In \[Gabrielli, Levi and Meo, 1996\] a class of selection rules is introduced for which independence of atom choices from goal instantiation is assured. These rules are named skeleton selection rules. Indeed, they are sensible only to a specific structural extract (the skeleton) of the applied clauses and the initial goal in the story of a derivation. As shown in \[Gabrielli, Levi and Meo, 1996\], instantiation independence is sufficient to proof a Strong Lifting Lemma which asserts that, for any skeleton rule$`S,`$ an SLD derivation of a goal $`G\gamma `$ via$`S`$can be lifted to a derivation of$`G`$via the same rule$`S,`$ relating in a quite strong sense the mgu’s and the resolvents. On the other hand, instantiation independence seems not sufficient to assure redundancy elimination tolerance. For example, in agreement to the definition in \[Gabrielli, Levi and Meo, 1996\], the selection rule of Example 2.2 is a skeleton rule, because choices only depend on the length of the initial goal and the ones of applied clauses. Really, choices are performed on the unique basis of the length of the actual resolvent, so that the rule of Example 2.2 can be seen as a case of $`state`$ skeleton selection rule. Anyhow, the rule is not tolerant to redundancy elimination. In order to point out the role of the hypothesis of specialisation independence with respect to derivation lifting, let us give the following example where Lifting Lemma 4.7 does not hold. Note that the used scheduling rule is instantiation independent, but it is not specialisation independent. ###### Example 4.9. Let us consider again the scheduling rule of Example 4.6. It is easy to recognise that Lifting Lemma 4.7 does not hold for such a rule. Indeed, let P be the following program: $`c1=pp[1],r[2],r[3]`$ $`c2=r`$. Now, in reference to the statement of Lemma 4.7, let: $`G=\{p[1],s[2],s[3]\}`$ $`G\gamma \underset{¯}{\tau }=\{p[1.1],s[2],s[2.5]\}`$ and $`X=\{r[1.5],r[1.6]\}`$. In Figure 2 an infinite p-SLD derivation of $`(G\gamma \underset{¯}{\tau }+X)`$ in P is shown . On the contrary, the only p-SLD derivation of $`G`$ in P is the following one {$`p[1],s[2],s[3]\}\stackrel{S,c1}{}\{s[2],p[2.5],r[2.6],r[2.7],s[3]\}.`$ which fails at the second resolvent. $`\mathrm{}`$ From the proofs of Lemmata 4.4 and 4.7, the proof of two corresponding assertions can be easily drawn. They are given in Lemma 4.10 below, and are valid for all scheduling rules in the case of two p-SLD derivations which are lowerings of each other. Part a) of the lemma may be viewed as a form of Variant Lemma. ###### Lemma 4.10 (determinism lemma). Let$`S`$be any scheduling rule and $`V`$ any arbitrary finite set of variables. Then let$`G`$and $`G^{}`$ be two p-goals such that $`G^{}`$ is a p-variant of $`G`$. The following implications hold: a)$`G\stackrel{S,D}{}Q`$ and $`G^{}\stackrel{S,D}{}R`$ $`R`$is a p-variant of $`Q`$, b)$`G^{}\stackrel{S,D}{}`$ $``$ $`Dr=(G\stackrel{S,D}{})`$, with $`nvar(Dr)V=\mathrm{}.`$ ###### Proof 4.11. Let us consider part a) of the lemma. By definition of p-variant it is $`G^{}=G\gamma \underset{¯}{\tau }`$, for a renaming $`\gamma `$ and a shifting $`\underset{¯}{\tau }`$. By fact (p1) in Lemma 4.4, i.e. “where $`\sigma `$ is a renaming if $`\gamma `$ is a renaming and $`D/X=\mathrm{}`$”, the result appears as an immediate consequence of the proof of Lemma 4.4 itself. It is sufficient to note that, if $`X`$ is empty and $`\gamma `$ is a renaming, the fact that the used scheduling rule is specialisation independent becomes useless. Indeed, in reference to the proofs of Lemma 4.4, though the hypothesis of specialisation independence is dropped, the first steps of (5) and (4) are congruent lowerings of each other, because every scheduling rule is deterministic. Similar considerations are possible for part b) of the lemma, in reference to the proof of Lemma 4.7. Now, let us give a property that is valid for all scheduling rules and derives easily from Lemma 4.10. It asserts that two p-SLD derivations $`Dr_1`$ and $`Dr_2`$, via the same scheduling rule$`S,`$ can be composed giving place to a longer derivation via$`S,`$ if the last resolvent of $`Dr_1`$ coincides with the first of $`Dr_2`$. ###### Property 4.3 (combination) Let$`S`$be any (state) scheduling rule. The following implication holds: $`Dr_1,Dr_2`$ with $`Dr_1=(G\stackrel{S,E}{}F),Dr_2=(F\stackrel{S,H}{}Q)`$ $`Dr=(G\stackrel{𝐸}{}F\stackrel{𝐻}{}R)`$, with $`Dr\mathrm{\Delta }(S)`$, where$`R`$is a p-variant of $`Q`$. ###### Proof 4.12. By Lemma 4.10-b) applied to $`Dr_2`$, a p-SLD derivation $`Dr^{}=(F\stackrel{S,H}{}R)`$ exists with $`nvar(Dr^{})(nvar(Dr_1)var(G))=\mathrm{}`$. Thus, $`Dr^{}`$ is standardised apart with respect to $`Dr_1`$. Since$`S`$is a state scheduling rule, $`Dr`$ is obtained as the composition of $`Dr_1`$ and $`Dr^{}`$. The fact that$`R`$is a p-variant of $`Q`$ follows from Lemma 4.10-a), applied to $`Dr_2`$ and $`Dr^{}`$. ## 5 Stack-queue selection rules Prolog interpreters adopt a leftmost scheduling policy such that the first atom in the goal is always selected for rewriting and is replaced in the resolvent by the body of the applied clause. In other words, the actual resolvent is maintained as a stack, the atom on the top of the stack is always selected for rewriting, while new atoms from the applied clause are pushed on the top of the stack. In analogy, a queue scheduling policy may be considered, which corresponds to a very simple case of fair selection rule (see \[Lloyd, 1987\]). As for the stack scheduling policy the first atom in the resolvent is always selected, but new atoms are positioned at the end of the old resolvent. Thus, the resolvent is treated as a queue of atoms and any queued atom is eventually selected in the case of infinite derivations In this section the class of stack-queue scheduling rules is defined, which is a generalisation of both stack and queue scheduling policies. According to stack-queue rules, for any clause$`c=(htB)`$, two p-goals $`M_s`$ and $`M_q`$ can be identified, with$`B=M_s\|M_q`$, such that the atoms in $`M_s`$ are always scheduled in stack mode while the atoms in $`M_q`$ are scheduled in queue mode. More formally, we have the following definition. As shown in the sequel of this Section 5, the stack-queue class turns out to be an operational characterisation of the class of specialisation independent scheduling rules ###### Definition 5.1 (stack-queue derivation steps). A set $`SQ`$ of derivation steps is said to be of stack-queue type, if it verifies the following condition. Given any clause$`c=(htB)`$, two p-goals $`M_s`$ and $`M_q`$ exist with $`M_s\|M_q=B`$, such that for any p-goal $`(a\|K)`$: $`a\|K\stackrel{SQ,c\xi ,\mu }{}RR=(M_s\xi \underset{¯}{\gamma }\|K\|M_q\xi \underset{¯}{\gamma })\mu .`$ The following property states that any set of stack-queue derivation steps is specialisation independent. Then, as stated in Theorem 5.3, any set of stack-queue derivation steps which satisfies the completeness property is a specialisation independent scheduling rule. ###### Property 5.1 (stack-queue implies specialisation independence) Let $`SQ`$ be a stack-queue set of derivation steps. Then $`SQ`$ is specialisation independent. ###### Proof 5.2. Let us consider two derivation steps in $`SQ`$ and suppose that derivation step (2) is a lowering of (1) by $`F`$. This means that (1) and (2) have the following form, where$`c=(htM_s\|M_q)`$: $`a\|K\stackrel{SQ,c}{}(M_s\xi ^{}\underset{¯}{\gamma }^{}\|K\|M_q\xi ^{}\underset{¯}{\gamma }^{})\alpha ^{}`$ (1) $`a\lambda \underset{¯}{\sigma }\|(K\lambda \underset{¯}{\sigma }+F)\stackrel{SQ,c}{}(M_s\xi ^{\prime \prime }\underset{¯}{\gamma }^{\prime \prime }\|(K\lambda \underset{¯}{\sigma }+F)\|M_q\xi ^{\prime \prime }\underset{¯}{\gamma }^{\prime \prime })\alpha ^{\prime \prime }.`$ (2) In order to show that $`SQ`$ is specialisation independent, we have to verify that derivation step (2) is a congruent lowering of (1) by $`F`$, i.e. a shifting $`\underset{¯}{\rho }`$ exists, such that: $`M_s\underset{¯}{\gamma }^{\prime \prime }\|M_q\underset{¯}{\gamma }^{\prime \prime }=(M_s\underset{¯}{\gamma }^{}\|M_q\underset{¯}{\gamma }^{})\underset{¯}{\rho },K`$$`\sigma `$$`=K\underset{¯}{\rho }`$. (3) By Property 3.2, a shifting $`\underset{¯}{\rho }`$ exists such that: $`M_s\underset{¯}{\gamma }^{\prime \prime }\|K\underset{¯}{\sigma }\|M_q\underset{¯}{\gamma }^{\prime \prime }=^{(Axiii)}(M_s\underset{¯}{\gamma }^{}\|K\|M_q\underset{¯}{\gamma }^{})\underset{¯}{\rho }=M_s\underset{¯}{\gamma }^{}\underset{¯}{\rho }\|K\underset{¯}{\rho }\|M_q\underset{¯}{\gamma }^{}\underset{¯}{\rho }.`$ Since it is evident that $`\mathrm{\#}M_s\underset{¯}{\gamma }^{}\underset{¯}{\rho }=\mathrm{\#}M_s\underset{¯}{\gamma }^{\prime \prime }`$ and $`\mathrm{\#}K\underset{¯}{\rho }=\mathrm{\#}K\underset{¯}{\sigma }`$, by Property 3.1-i) we have: $`M_s\underset{¯}{\gamma }^{\prime \prime }=M_s\underset{¯}{\gamma }^{}\underset{¯}{\rho },M_q\underset{¯}{\gamma }^{\prime \prime }=M_q\underset{¯}{\gamma }^{}\underset{¯}{\rho },K\underset{¯}{\sigma }=K\underset{¯}{\rho }`$, which immediately implies assertion (3). ###### Theorem 5.3 (stack-queue scheduling rules). Let $`SQ`$ be a complete set of stack-queue derivation steps. Then $`SQ`$ is a specialisation independent scheduling rule. ### 5.1 Specialisation independence implies stack-queue Now, we prove (Theorem 5.6) that any specialisation independent scheduling rule is actually a stack-queue rule. Thus, combining this fact with Theorem 5.3, we have that Definition 4.2 and the operational characterisation of Definition 5.1 identify the same family of scheduling rules. To this aim, let us show the following lemma. ###### Lemma 5.4 (not internal positioning). Let$`S`$be a specialisation independent scheduling rule. Given any clause$`c=(htB)`$, for every derivation step of the form: $`a\|K\stackrel{S,c\xi ,\eta }{}R,`$ (1) two subgoals $`M_s`$ and $`M_q`$ exist, with$`B=M_s\|M_q`$, such that: $`R=(M_s\xi \underset{¯}{\gamma }\|K\|M_q\xi \underset{¯}{\gamma })\eta .`$ ###### Proof 5.5. Let us consider a p-goal like: $`a\|K\underset{¯}{\omega }_1\|K\underset{¯}{\omega }_2\|\mathrm{}\|K\underset{¯}{\omega }_n`$, with $`n>\mathrm{\#}B`$. On the basis of (1), by Property 3.3 a derivation step also exists of the following form: $`a\|K\underset{¯}{\omega }_1\|\mathrm{}\|K\underset{¯}{\omega }_n\stackrel{S,c\gamma ,\mu }{}(Q=((K\underset{¯}{\omega }_1\|\mathrm{}\|K\underset{¯}{\omega }_n)+B\gamma \underset{¯}{\tau })\mu )`$. (2) Since $`n>\mathrm{\#}B`$, an index $`j`$ must exist such that no atom of$`B`$has been positioned inside $`K\underset{¯}{\omega }_j`$. A priori several $`j`$’s might exist. Without loss of generality, we take any one of them. Thus, two p-goals $`M_s`$ and $`M_q`$ must exist, with $`M_s\|M_q=B`$, such that: $`Q=(M_s\underset{¯}{\tau }\gamma +(K\underset{¯}{\omega }_1\|\mathrm{}\|K\underset{¯}{\omega }_{j1}))\|K\underset{¯}{\omega }_j\|(M_q\underset{¯}{\tau }\gamma +(K\underset{¯}{\omega }_{j+1}\|\mathrm{}\|K\underset{¯}{\omega }_n))\mu .`$ (3) Now, by definition, derivation step (1) has the form: $`a\|K\stackrel{S,c}{}(R=(K+B\xi \underset{¯}{\sigma })\eta )`$. (1a) Since$`S`$is a specialisation independent rule, step (2) is a congruent lowering of step (1a) by the subgoal $`(K\underset{¯}{\omega }_1\|\mathrm{}\|K\underset{¯}{\omega }_{j1}\|K\underset{¯}{\omega }_{j+1}\|\mathrm{}\|K\underset{¯}{\omega }_n)`$, so that a shifting $`\underset{¯}{\rho }`$ exists with $`K\underset{¯}{\rho }=K\underset{¯}{\omega }_j`$ and $`B\underset{¯}{\sigma }\underset{¯}{\rho }=B\underset{¯}{\tau }=M_s\underset{¯}{\tau }\|M_q\underset{¯}{\tau }`$. Then, recalling that (3) implies $`M_s\underset{¯}{\tau }K\underset{¯}{\omega }_jM_q\underset{¯}{\tau }`$, we obtain: $`(K+B\xi \underset{¯}{\sigma })\underset{¯}{\rho }=K\underset{¯}{\omega }_j+B\xi \underset{¯}{\tau }=(M_s\xi \underset{¯}{\tau }\|K\underset{¯}{\omega }_j\|M_q\xi \underset{¯}{\tau })=(M_s\xi \underset{¯}{\tau }\|K\underset{¯}{\rho }\|M_q\xi \underset{¯}{\tau })`$. Finally: $`R=(K+B\xi \underset{¯}{\sigma })\eta =(K+B\xi \underset{¯}{\sigma })\eta \underset{¯}{\rho }\underset{¯}{\rho }^1=(M_s\xi \underset{¯}{\tau }\underset{¯}{\rho }^1\|K\|M_q\xi \underset{¯}{\tau }\underset{¯}{\rho }^1)\eta `$. The following Theorem 5.6 shows that, for any scheduling rule, specialisation independence implies that the rule is stack-queue. Together with Theorem 5.3, this result proofs that stack-queue is an operational characterisation of the set of specialisation independent scheduling rules. ###### Theorem 5.6 (specialisation independence implies stack-queue). Let$`S`$be a specialisation independent scheduling rule. Given any clause$`c=(htB),`$ two p-goals $`M_s`$ and $`M_q`$ exist, with $`M_s\|M_q=B`$, such that for every derivation step of the form: $`a\|K\stackrel{S,c\xi ,\eta }{}R`$(1) it is: $`R=(M_s\xi \underset{¯}{\pi }\|K\|M_q\xi \underset{¯}{\pi })\eta `$. ###### Proof 5.7. Let $`p`$ be the predicate symbol of atom $`ht`$. Consider a p-atom $`b`$ of the form $`b=p(x_1,\mathrm{},x_k)[s]`$, where $`x_1,\mathrm{},x_k`$ are distinct variables. Then, consider a ground p-atom $`r`$ such that $`br`$. By construction of $`b`$ and completeness of$`S,`$ a derivation step of the type $`(b\|r\stackrel{S,c}{})`$ exists, which necessarily has the following form because $`r`$ is a single atom: $`b\|r\stackrel{S,c}{}(M_s\lambda \underset{¯}{\epsilon }\|r\|M_q\lambda \underset{¯}{\epsilon })\mu ,`$ with$`B=M_s\|M_q`$. (2) Now, let us prove that $`M_s\|M_q`$ is the partition of$`B`$which is required by the thesis. Consider derivation step (1). Two cases are possible, either $`K=\mathrm{}`$ or $`K\mathrm{}`$. Case 1 $`(K=\mathrm{}).`$ In this case we have: $`a\stackrel{S,c}{}(R=B\xi \underset{¯}{\pi }\eta =(M_s\xi \underset{¯}{\pi }\|M_q\xi \underset{¯}{\pi })\eta ).`$ Case 2 $`(K\mathrm{}).`$ On the basis of (1), we have that also p-atom $`a`$ has $`p`$ as a predicate symbol, so that a substitution $`\tau `$ and a shifting $`\underset{¯}{\sigma }`$ exist with $`a=b\tau \underset{¯}{\sigma }`$. By (1) and Property 3.3, a derivation step exists like: $`(b\tau \underset{¯}{\sigma }\|(r\tau \underset{¯}{\sigma }+K)=a\|(K+r\underset{¯}{\sigma }))\stackrel{S,c\xi ^{},\eta ^{}}{}Q`$, (4) where by Lemma 5.4 we have that: $`Q=(N_s\xi ^{}\underset{¯}{\gamma }\|(r\underset{¯}{\sigma }+K)\|N_q\xi ^{}\underset{¯}{\gamma })\eta ^{}`$, with$`B=N_s\|N_q`$. (5) The proof can be now completed by exploiting derivation step (4) as a sort of “bridge” between (1) and (2). In fact, since$`S`$is specialisation independent rule, derivation step (4) is a congruent lowering of step (2) by $`K`$, so that a shifting $`\underset{¯}{\rho }^{}`$ exists with $`r\underset{¯}{\rho }^{}=r\underset{¯}{\sigma }`$ and $`(M_s\underset{¯}{\epsilon }\|M_q\underset{¯}{\epsilon })\underset{¯}{\rho }^{}=N_s\underset{¯}{\gamma }\|N_q\underset{¯}{\gamma }`$. As a consequence (see (5) and (2)), we can write: $`N_s\underset{¯}{\gamma }\|r\underset{¯}{\sigma }\|N_q\underset{¯}{\gamma }=N_s\underset{¯}{\gamma }\|N_q\underset{¯}{\gamma }+r\underset{¯}{\sigma }=(M_s\underset{¯}{\epsilon }\|M_q\underset{¯}{\epsilon })\underset{¯}{\rho }^{}+r\underset{¯}{\rho }^{}=(M_s\underset{¯}{\epsilon }\|r\|M_q\underset{¯}{\epsilon })\underset{¯}{\rho }^{}`$, with $`r\underset{¯}{\sigma }=r\underset{¯}{\rho }^{}`$. Then, by Property 3.1-ii we have that $`N_s\underset{¯}{\gamma }=M_s\underset{¯}{\epsilon }\underset{¯}{\rho }^{}`$, which obviously implies: $`\mathrm{\#}N_s=\mathrm{\#}M_s`$. (6) Now, let us note that by Lemma 5.4 it must be: $`R=(A_s\xi \underset{¯}{\pi }\|K\|A_q\xi \underset{¯}{\pi })\eta ,`$ with$`B=A_s\|A_q`$. (7) Since$`S`$is a specialisation independent rule, derivation step (4) is a congruent lowering of step (1) by $`r\underset{¯}{\sigma }`$, so that a shifting $`\underset{¯}{\rho }^{\prime \prime }`$ exists with $`K\underset{¯}{\rho }^{\prime \prime }=K`$ and $`(A_s\underset{¯}{\pi }\|A_q\underset{¯}{\pi })\underset{¯}{\rho }^{\prime \prime }=N_s\underset{¯}{\gamma }\|N_q\underset{¯}{\gamma }`$. As a consequence (see (5) and (7)) we can write: $`N_s\underset{¯}{\gamma }\|K\|N_q\underset{¯}{\gamma }=N_s\underset{¯}{\gamma }\|N_q\underset{¯}{\gamma }+K=(A_s\underset{¯}{\pi }\|A_q\underset{¯}{\pi })\underset{¯}{\rho }^{\prime \prime }+K\underset{¯}{\rho }^{\prime \prime }=(A_s\underset{¯}{\pi }\|K\|A_q\underset{¯}{\pi })\underset{¯}{\rho }^{\prime \prime }`$, with $`K=K\underset{¯}{\rho }^{\prime \prime }`$. Then, by Property 3.1-ii, we have that $`N_s\underset{¯}{\gamma }=A_s\underset{¯}{\pi }\underset{¯}{\rho }^{\prime \prime }`$, which obviously implies: $`\mathrm{\#}N_s=\mathrm{\#}A_s`$. (8) By (2) and (7), it is $`M_s\|M_q=A_s\|A_q=B`$. By (6), (8) and Property 3.1-i), we have that: $`A_s=M_s`$ and $`A_q=M_q`$. Substituting in (7) , the thesis is obtained. ### 5.2 Notes on the structure of stack-queue derivations Let us consider a stack-queue derivation like: $`A\|B\stackrel{SQ,M,\sigma }{}`$, where $`M=c_1,c_2,\mathrm{}c_h`$ and $`M/B=\mathrm{}.`$(1) By definition of stack-queue scheduling rules, only atoms in $`A`$ together with atoms deriving from $`A`$ and allocated in stack mode can be rewritten in derivation (1). Thus, derivation (1) has the form: $`A\|B\stackrel{SQ,c_1\xi _1,\sigma _1}{}X_1\|A_1\|B\sigma _1\|Y_1\stackrel{SQ,c_2\xi _2,\sigma _2}{}\mathrm{}`$ $`X_i\|A_i\|B\sigma _1\mathrm{}\sigma _i\|Y_i\stackrel{SQ,c_{i+1}\xi _{i+1},\sigma _{i+1}}{}\mathrm{}\stackrel{SQ,c_h\xi _h,\sigma _h}{}X_h\|A_h\|B\sigma _1\mathrm{}\sigma _h\|Y_h`$, (1a) where: * each $`X_i`$ is formed by new atoms deriving from $`A`$ which are allocated in stack mode, * each $`A_i`$ is formed by atoms of $`A`$ which are not yet rewritten, * each $`Y_i`$ is formed by new atoms deriving from $`A`$ which are allocated in queue mode. The above structural considerations suggest the following formal definition. ###### Definition 5.8 (A-preq type derivations). A p-SLD derivation, of the form $`A\|B\stackrel{SQ,M}{}`$, is of pre-queued type w.r.t. the subgoal $`A`$ (simply written $`A`$-preq type in the following) if the only rewritten atoms are: \- atoms from the subgoal $`A`$, \- atoms deriving from $`A`$ and allocated in stack mode. Note that Definition 5.8 is significant even if $`B=\mathrm{}`$. It is evident that any $`A`$-preq derivation has the form (1a). In the sequel we use the following shortened notation to represent $`A`$-preq type derivations: $`A\|B\stackrel{SQ,M,\sigma }{}A^s\|B\sigma \|A^q,`$ (Ap) where, with reference to (1a), $`A^s=X_h\|A_h`$ stands for “stacked subgoal derived from $`A`$”, and $`A^q=Y_h`$ means “queued subgoal derived from $`A`$”. It is evident that in any preq type derivation we have $`M/B=\mathrm{}`$. The following definition characterises an $`A`$-queued derivation as an $`A`$-preq derivation where all atoms of $`A`$ are rewritten together with all atoms deriving from $`A`$ and allocated in stack mode, i.e. $`A^s`$= $`\mathrm{}`$. Intuitively, an $`A`$-queued derivation is an $`A`$-preq derivation which cannot be extended without loosing its $`A`$-preq nature. Indeed, the acronym “$`A`$-preq” stands for “$`A`$-pre-queued” derivation. ###### Definition 5.9 (A-queued derivations). Let $`SQ`$ be a stack-queue scheduling rule. A derivation which is of $`A`$-preq type and has the form: $`A\|B\stackrel{SQ,K,\sigma }{}B\sigma \|A^q`$ (Aq) is said to be queued w.r.t. $`A`$ (simply written $`A`$-queued in the following). In the following Section 5.3, we will exploit the notations introduced in (Ap) and (Aq) to represent $`A`$-preq type and $`A`$-queued derivations, respectively. It is worth noting that starting from a p-goal of the form $`A\|B`$, when the $`A`$-queued derivation is reached, the last resolvent presents a situation where the roles of $`A`$ and$`B`$are exchanged. In practice, restarting from $`B\sigma \|A^q`$, the derivation can attempt to proceed towards a $`(B\sigma )`$-queued derivation. The proof of an important result in Section 5.3 (Duplication Theorem 5.10) is based on this cyclic behaviour of stack-queue derivations. ### 5.3 Duplication tolerance In this section an important property is shown for stack-queue scheduling rules. Let us give an intuitive presentation of this result, which is stated in the full duplication theorem (Theorem 5.12). Suppose that a p-SLD derivation $`Dr`$ of$`G`$in P can be developed via a stack-queue scheduling rule $`SQ`$. Then consider a p-goal $`G^{}`$ which is equal to$`G`$apart from the duplication of some atoms. Furthermore, suppose that each copy is scheduled after the corresponding original atom. In this hypothesis, the full duplication theorem asserts that a p-SLD derivation of $`G^{}`$ in P exists via the same scheduling rule $`SQ`$, where all derivation steps of $`Dr`$ are redone in the order. The full duplication theorem is basic for the proof of the final results of the paper, i.e. results about redundancy elimination tolerance which are given in Section 6. Indeed, let us consider the problem of preserving program termination. Intuitively, program termination is preserved if the introduction of redundancy elimination does not provoke any really different new derivations. Reversing the viewpoint, termination is retained if any derivation, developed in presence of redundancy elimination, can be traced again when redundancy is left in place. The full duplication theorem asserts this kind of fact in the simplest case, i.e. when redundancy has the form of a replica of atoms already present in the initial p-goal, provided that the scheduling rule is of stack-queue type. First we show a duplication theorem (Theorem 5.10) which is valid when only one atom or group of adjacent atoms is duplicated. Then the result is easily extended to obtain the full theorem. Though intuitive in appearance, Theorem 5.10 has a relatively complex proof. In this section we give only a sketch of the argument. In the sketch, we will make reference to the particular case of completely ground derivations, i.e. derivations such that all resolvents are ground. This simplification will allow us to highlight the essence of the argument, without having to do with technical problems deriving from variable instantiations. Formal presentation of the proof of Theorem 5.10 is given in Appendix B. Note that the hypothesis of ground resolvents is verified in the case that no new variable is present in clause bodies and initial goals are ground. ###### Theorem 5.10 (duplication theorem). Let P be a logic program and $`SQ`$ a stack-queue scheduling rule. Given two p-goals of the form $`A\|B\|C\|D`$ and $`A\|B\|C\|B\underset{¯}{\pi }\|D`$, the following implication holds: $`A\|B\|C\|D\stackrel{SQ,X.P}{}Q`$ (1) $`Y`$ such that $`A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,Y.P}{}R`$ with $`X_LY`$ and $`\mathrm{\#}Q\mathrm{\#}R.`$ ###### Proof 5.11 (Proof (sketch)). Let $`\mathrm{\Delta }(SQ,n)`$ denote the subset of $`\mathrm{\Delta }(SQ)`$ such that, for any derivation $`Dr`$ in $`\mathrm{\Delta }(SQ,n)`$ , it is $`\mathrm{\#}Drn`$, where $`\mathrm{\#}Dr`$ denotes the length of $`Dr`$. We show the thesis by induction on $`n`$. In other words, we show that the thesis holds when derivation (1) belongs to $`\mathrm{\Delta }(SQ,n)`$, for any $`n0`$. The fact is obvious for $`\mathrm{\Delta }(SQ,0)`$. In order to justify the inductive step from $`\mathrm{\Delta }(SQ,n1)`$ to $`\mathrm{\Delta }(SQ,n)`$, for $`n>0`$, let us consider a derivation like: $`(A\|B\|C\|D\stackrel{X.P}{}Q)\mathrm{\Delta }(SQ,n)`$(1a) and show that $`(A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,Y.P}{}R)`$ exists with $`X_LY`$ and $`\mathrm{\#}Q\mathrm{\#}R`$. The following three possible situations must be taken into account. Then, we start with case 3, which is the most significant one. 1. derivation (1a) is of $`(A\|B\|C)`$-preq type, 2. derivation (1a) is of $`(A\|B\|C\|D)`$-preq type, and not of $`(A\|B\|C)`$ -preq type, 3. derivation (1a) is not of $`(A\|B\|C\|D)`$-preq type. Case 3. As already said, the simplified argument, which we use in this sketch, works in the hypothesis that all resolvents are ground, so that derivation (1a) has the following form: $`A\|B\|C\|D\stackrel{𝐻}{}B\|C\|D\|A^q\stackrel{𝐾}{}C\|D\|A^q\|B^q\stackrel{𝑀}{}`$ $`D\|A^q\|B^q\|C^q\stackrel{𝑁}{}A^q\|B^q\|C^q\|D^q\stackrel{𝑇}{}Q`$, where $`H\|K\|M\|N\|T=X`$. (2) Then, it is intuitive that a derivation can be constructed like the following, where $`\underset{¯}{\varphi }`$ is a suitable shifting: $`A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,H}{}B\|C\|B\underset{¯}{\pi }\|D\|A^q\stackrel{SQ,K}{}`$ $`C\|B\underset{¯}{\pi }\|D\|A^q\|B^q\stackrel{SQ,M}{}B\underset{¯}{\pi }\|D\|A^q\|B^q\|C^q\stackrel{SQ,K}{}`$ (3) $`D\|A^q\|B^q\|C^q\|B^q\underset{¯}{\varphi }\stackrel{SQ,N}{}A^q\|B^q\|C^q\|B^q\underset{¯}{\varphi }\|D^q.`$ By construction of (2), $`A^q\|B^q\|C^q\|D^q\stackrel{𝑇}{}Q`$ is a derivation belonging to $`\mathrm{\Delta }(SQ,m)`$, with $`m<n`$. By inductive hypothesis, a derivation exists such that: $`A^q\|B^q\|C^q\|B^q\underset{¯}{\varphi }\|D^q\stackrel{SQ,Y^{}.P}{}R^{}`$, (5) with $`T_LY^{}`$ and $`\mathrm{\#}Q\mathrm{\#}R^{}`$. (5a) By Property 4.3, derivations (3) and (5) can be combined to yield a derivation of the form: $`A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,(H\|K\|M\|K\|N).P}{}A^q\|B^q\|C^q\|B^q\underset{¯}{\varphi }\|D^q\stackrel{SQ,Y^{}.P}{}R`$, (6) where$`R`$is a p-variant of $`R^{}`$, which implies $`\mathrm{\#}R=\mathrm{\#}R^{}`$. Finally: $`X=H\|K\|M\|N\|T_L^{(5a)}(H\|K\|M\|K\|N\|Y^{}),\mathrm{\#}Q^{(5a)}\mathrm{\#}R^{}=\mathrm{\#}R`$. Case 2. Derivation (1a) has the form $`A\|B\|C\|D\stackrel{H\|K\|M\|N}{}D^s\|A^q\|B^q\|C^q\|D^q`$, where $`H\|K\|M\|N=X`$. Analogously to case 3), a derivation can be constructed like: $`A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,(H\|K\|M\|K\|N).P}{}D^s\|A^q\|B^q\|C^q\|B^q\underset{¯}{\pi }\|D^q`$. Case 1. Derivation (1a) has the form $`A\|B\|C\|D\stackrel{𝑋}{}(A\|B\|C)^s\|D\|(A\|B\|C)^q`$. A derivation exists like: $`(A\|B\|C)\|B\underset{¯}{\pi }\|D\stackrel{SQ,X}{}(A\|B\|C)^s\|B\underset{¯}{\pi }\|D\|(A\|B\|C)^q.`$ Now we can state and prove the full duplication theorem, which extends the previous Theorem 5.10 to the duplication of two or more not adjacent atoms in the initial goal of a p-SLD derivation. ###### Theorem 5.12 (full duplication theorem). Let P be a logic program and $`SQ`$ a stack-queue scheduling rule. Given a p-goal $`N+F`$ such that: $`b[s]F,b[s^{}]N`$ with $`s^{}<s`$, the following implication holds: $`N\stackrel{SQ,M.P}{}Q`$ $``$ $`Y`$ such that $`N+F\stackrel{SQ,Y.P}{}R`$ with $`M_LY`$ and $`\mathrm{\#}Q\mathrm{\#}R`$. ###### Proof 5.13. By hypothesis, the subgoal$`F`$is made of duplicated atoms. Then, the proof is by induction on the length of $`F`$. Indeed, if$`F`$is empty the thesis is true. Now, suppose that the thesis is already proven for any$`F`$with $`\mathrm{\#}F=n0`$. Then let us consider any p-goal$`G=F\|b[s]`$ with $`\mathrm{\#}F=n`$. By inductive hypothesis a derivation exists such that: $`N+F\stackrel{SQ,Z.P}{}S,`$ with $`M_LZ`$ and $`\mathrm{\#}Q\mathrm{\#}S`$. By hypothesis, three p-goals $`A,C`$ and $`D`$ exist together with a p-atom $`b[s^{}]`$, such that: $`N+(F\|b[s])=A\|b[s^{}]\|C\|b[s]\|D`$ and $`N+F=A\|b[s^{}]\|C\|D`$. As a consequence, Theorem 5.10 can be applied to $`N+F`$ and $`N+(F\|b[s])`$ yielding: $`N+(F\|b[s])\stackrel{SQ,Y.P}{}R`$, with $`Z_LY`$ and $`\mathrm{\#}S\mathrm{\#}R`$. Now the induction step is completed, because: $`M_LZ_LY`$ and $`\mathrm{\#}Q\mathrm{\#}S\mathrm{\#}R`$. ## 6 Redundancy elimination tolerance In this section, the tolerance of stack-queue scheduling rules to redundancy elimination is considered. The preservation of program termination in shown in Section 6.1. The preservation of the completeness of $`EVR_L`$ loop check is shown in Section 6.2 for function free programs. First, the idea of goal reduction, which is originally given in \[Ferrucci, Pacini and Sessa, 1995\] and is recalled in Definition 2.1 of this paper, is restated. Indeed, in Section 2 little attention is paid to the positions of atoms which are removed from a resolvent. However, if the execution is based on atom priority values, it is intuitive that removing an atom without any convenient expedient may overthrow the essence of previous atom scheduling. Thus, a refined definition of goal reduction is given below (Definition 6.1) which fits the frame of priority SLD derivation mechanisms. The inspiring idea of priority reduction is quite simple. According to Definition 2.1, for any removed atom $`b`$, an eliminating atom $`a=b\tau `$ exists which remains in the reduced resolvent. Several removed atoms may share the same eliminating one. In reference to Definition 6.1 below, for any eliminating atom $`a_j[p_j]`$, the corresponding subset $`A_j`$ of eliminated atoms is pointed out. Then, except for the case $`a_j[p_j]A_j`$, any $`a_j[p_j]`$ is advanced to the least priority value in $`A_j`$. In other words, each eliminating atom is advanced to replace the first scheduled atom among its eliminated ones. Intuitively, the aim is to restore the essence of the previous atom priorities. The notation $`\{+A_j,1jh\}`$ will represent the merging $`A_1+A_2+\mathrm{}+A_h`$, and the notation $`prs(A_j)`$ the set of priority values in $`A_j`$. ###### Definition 6.1 (priority reduced goals). Let $`X`$ be a set of variables, $`\tau `$ a substitution and$`G`$a p-goal. A p-goal$`N`$is a reduced p-goal of$`G`$by $`\tau `$ up to X, denoted by$`G>>^\tau N`$, if the following conditions hold: i)$`G=F+\{+a_j[p_j],\mathrm{\hspace{0.33em}1}jh\}+\{+A_j,\mathrm{\hspace{0.33em}1}jh\}`$, where $`b[s]A_j,b\tau =a_j,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}jh`$, ii)$`N=F+\{+a_j[r_j],\mathrm{\hspace{0.33em}1}jh\},`$ where $`r_j=min(\{p_j\}prs(A_j)),\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}jh,`$ iii) $`x(Xvar(N))`$ it is $`x\tau =x`$. ###### Example 6.2. Given the p-goal $`G=p(z)[1],q(w)[2],p(a)[3],p(y)[4],q(v)[5]`$, the following$`N`$is a reduced p-goal of$`G`$by the substitution $`\tau =\{z/a,y/a,v/w\}`$: $`N=p(a)[1],q(w)[2]`$. Note that $`p(a)[3]`$ has been advanced to replace the first of the atoms it eliminates, that is $`p(z)[1]`$$`\mathrm{}`$ Now, the idea of priority reduced SLD derivation can be defined as a generalisation of Definition 3.5. In essence a priority reduced SLD derivation is a p-SLD derivation where, at any step, a priority reduction of the resolvent according to Definition 6.1 is allowed. ###### Definition 6.3 (priority Reduced SLD derivation). Let P be a program and $`G_o`$ a p-goal. A priority reduced SLD derivation of $`G_o`$ in P (p-RSLD derivation for short) is a possibly infinite sequence of priority reductions and derivation steps $`G_o>>^{\alpha _o}N_o\stackrel{c_o\xi _o,\theta _o}{}G_1\mathrm{}G_k>>^{\alpha _k}N_k\stackrel{c_k\xi _k,\theta _k}{}G_{k+1}>>^{\alpha _{k+1}}N_{k+1}\mathrm{}`$ where, for any $`j0`$, i) $`c_j`$ is a clause in P, ii) $`var(c_j\xi _j)(var(G_o)var(c_o\xi _o)\mathrm{}var(c_{j1}\xi _{j1}))=\mathrm{}`$, iii) $`G_j>>^{\alpha _j}N_j`$ up to $`var(G_o\theta _o\mathrm{}\theta _{j1})`$. The notation $`G\stackrel{S,D}{}>>N`$ will be used to represent a p-RSLD derivation which is developed in agreement with the scheduling rule$`S`$using the template $`D`$. The last resolvent$`N`$is intended to be a reduced resolvent. ### 6.1 Termination preserving In this section, the redundancy elimination tolerance of stack-queue scheduling rules is shown, with reference to program termination (Theorem 6.6). The following lemma is fundamental for proving the preservation of termination, as well as the preservation of $`EVR_L`$ loop check completeness. ###### Lemma 6.4. Let P be a program and $`SQ`$ a stack-queue scheduling rule. The following implication holds: $`G\stackrel{SQ,X.P}{}>>Q`$(1) $`Z`$ such that$`G\stackrel{SQ,Z.P}{}R`$, with $`X_LZ,\mathrm{\#}Q\mathrm{\#}R`$. ###### Proof 6.5. The proof is by induction on the length of $`X`$. If $`\mathrm{\#}X=0`$, the thesis is trivially verified with $`Z=\mathrm{}`$. Then let us consider $`X=c\|H`$. Derivation (1) may be rewritten as: $`(G>>^\tau N)\stackrel{SQ,c}{}F\stackrel{SQ,H}{}>>Q.`$ (2) Since $`\mathrm{\#}H<\mathrm{\#}X`$, by inductive hypothesis, a p-SLD derivation exists of the form: $`F\stackrel{SQ,K.P}{}T`$, with $`H_LK,\mathrm{\#}Q\mathrm{\#}T`$. (3) By Property 4.3, the first derivation step of (2) and derivation (3) can be combined to yield a derivation of the following form: $`N\stackrel{SQ,c}{}F\stackrel{SQ,K}{}S,`$ where$`S`$is a p-variant of $`T`$. (4) Now, let us consider the p-goal $`G\tau `$. With reference to Definition 6.1, we have that: $`G\tau =(F+\{+a_j[p_j],\mathrm{\hspace{0.33em}1}jh\})\tau +\{+A_j\tau ,\mathrm{\hspace{0.33em}1}jh\}=^{(Def.\text{6.1}iii)}`$ $`F+\{+a_j[p_j],\mathrm{\hspace{0.33em}1}jh\}+\{+A_j\tau ,\mathrm{\hspace{0.33em}1}jh\}=^{(Def.\text{6.1}iii)}`$ $`F+\{+a_j[r_j],\mathrm{\hspace{0.33em}1}jh\}+\{+A_j\tau \{r_j/p_j\},\mathrm{\hspace{0.33em}1}jh\}=`$ $`=N+\{+A_j\tau \{r_j/p_j\},\mathrm{\hspace{0.33em}1}jh\}^[`$<sup>2</sup><sup>2</sup>2The notation $`A_j\tau \{r_j/p_j\}`$ means that the priority value $`r_j`$ is replaced by $`p_j`$ in the p-goal $`A_j\tau `$.<sup>\]</sup>, where $`a_j[r_j]A_j\tau \{r_j/p_j\}`$ and any atom in $`A_j\tau \{r_j/p_j\}`$ is a duplicate of $`a_j,`$ $`1jh`$. Then,$`N`$and $`G\tau `$ verify the hypothesis of Theorem 5.12. As a consequence, by (4) a derivation also exists such that: $`G\tau \stackrel{SQ,Z.P}{}V,`$ with $`c\|K_LZ`$ and $`\mathrm{\#}S\mathrm{\#}V`$ (5) Now, let us apply Lifting Lemma 4.7 to (5). We obtain that a p-SLD derivation exists like: $`G\stackrel{SQ,Z.P}{}R`$. (6) where, applying Lowering Lemma 4.4 to (5) and (6), we have that $`\mathrm{\#}V=\mathrm{\#}R`$. Finally, we conclude: $`X=c\|H_L^{(3)}c\|K_L^{(5)}Z,`$ $`\mathrm{\#}Q^{(3)}\mathrm{\#}T=^{(4)}\mathrm{\#}S^{(5)}\mathrm{\#}V=\mathrm{\#}R.`$ ###### Theorem 6.6 (termination preserving). Let P be a program,$`G`$a p-goal and $`SQ`$ a stack-queue scheduling rule. If every p-SLD derivation of$`G`$in P via $`SQ`$ is finite, then any p-RSLD derivation via $`SQ`$ is finite too. ###### Proof 6.7. Let $`T`$ be the p-SLD tree of$`G`$in P via $`SQ`$. By hypothesis, every p-SLD derivation of$`G`$in P via $`SQ`$ is finite. As a consequence, since $`T`$ is a finitely branching tree, by Konig’s lemma (see Theorem K, in \[Knuth, 1997\]) $`T`$ is a finite tree. Let $`f`$ be the depth of $`T`$. Given any p-RSLD derivation of the form$`G\stackrel{SQ,X.P}{}>>`$, by Lemma 6.4 a p-SLD derivation of the form$`G\stackrel{SQ,Z.P}{}`$ exists in $`T`$, with $`X_LZ`$. But $`\mathrm{\#}Zf`$, so that we obtain $`\mathrm{\#}X\mathrm{\#}Zf.`$ In conclusion, the length of all p-RSLD derivations of$`G`$in P via $`SQ`$ is limited by $`f`$. Let us close this section with two examples which show that both stack-queue scheduling and eliminating atom advancement are essential for redundancy elimination tolerance. The first example shows the necessity of advancement of eliminating atoms. The second one is an example of state scheduling rule which is not tolerant to redundancy elimination, though goal reduction is performed in agreement with Definition 6.1. Of course, the scheduling rule is not of stack-queue type. In the following sketches of p-SLD and p-RSLD derivations, explicit indication of priority values is omitted, for the sake of brevity. ###### Example 6.8. Let us consider the stack scheduling rule (i.e. the usual leftmost rule) and the following single clause program P: $`c=pq(x)\|p`$. It is evident that all p-SLD derivations fail. However, if advancement of eliminating atoms is not performed, an infinite p-RSLD derivation of P exists, as shown in Figure 3. $`\mathrm{}`$ ###### Example 6.9. Let$`S`$be a scheduling rule which behaves as a stack rule, with an exception when atoms having $`s`$ as a predicate symbol are rewritten. In this case new atoms are positioned immediately after the first old atom, if one exists. Then, let us consider the logic program P consisting of the following clauses: $`c1=r`$ $`c2=s(x,y)t(x,y)`$ $`c3=q(x,y)r\|s(z,y)\|r\|q(x,z)`$. It is easy to verify that all p-SLD derivations of P terminate independently of the initial p-goal. In fact, given a p-SLD derivation of$`G`$in P, where$`G`$is any p-goal, two cases are possible: either an atom with predicate symbol $`q`$ is rewritten or not. If no atom with predicate symbol $`q`$ is rewritten, the derivation terminates evidently. Otherwise the derivation fails, as described below: $`G\stackrel{𝑆}{}q(..)\|K\stackrel{S,c3}{}r\|s(..)\|r\|q(..)\|K\stackrel{S,c1}{}`$ $`s(..)\|r\|q(..)\|K\stackrel{S,c2}{}r\|t(..)\|q(..)\|K\stackrel{S,c1}{}t(..)\|q(..)\|K.`$ Now let us show that, if reduction of resolvents is allowed, an infinite p-RSLD derivation of P exists. It is easy to verify that the infinite RSLD derivation in Figure 4 cannot be pruned neither by $`EVR_L`$ loop check nor by more powerful checks (like $`SIR_M`$) which are based on subsumption relationships between resultants \[Bol, Apt and Klop, 1991\]$`\mathrm{}`$ ### 6.2 Preserving the completeness of $`EVR_L`$ loop check In this section we prove the preservation of $`EVR_L`$ loop check completeness, passing from p-SLD to p-RSLD. The result holds for function free programs, provided that stack-queue scheduling rules are used in combination with priority reduction of resolvents, as introduced in Definition 6.1. The section starts with a characterisation of $`EVR_L`$ loop check which exploits the concept of priority shifting and is equivalent to the one stated in Definition 2.3. In essence, passing from Definition 2.3 to Definition 6.10 below, only assertion ii) is modified. On the other hand, the requirement $`N_j=N_i\tau \underset{¯}{\tau }`$ is plainly equivalent to $`N_i\tau =_LN_j`$, since any shifting $`\underset{¯}{\tau }`$ implies that the order of atoms is preserved. ###### Definition 6.10 (priority Equality Variant Check for Resultants). A p-RSLD derivation $`G_o>>^{\alpha _o}N_o\stackrel{c_o\xi _o,\theta _o}{}G_1\mathrm{}G_{h1}>>^{\alpha _{h1}}N_{h1}\stackrel{c_{h1}\xi _{h1},\theta _{h1}}{}G_h>>^{\alpha _h}N_h\mathrm{}`$ is pruned by priority Equality Variant of Resultant check (called p-$`EVR_L`$ check, in the following), if for some $`i`$ and $`j`$, with $`0i<j`$, a renaming $`\tau `$ and a shifting $`\underset{¯}{\tau }`$ exist such that: i) $`G_o\theta _o\mathrm{}\theta _{j1}=G_o\theta _o\mathrm{}\theta _{i1}\tau `$, ii) $`N_j=N_i\tau \underset{¯}{\tau }`$. With reference to the above definition, any couple $`Rs_h=[N_h,G_o\theta _o\mathrm{}\theta _{h1}]`$ is a reduced resultant. Given two reduced resultants $`Rs_j=[N_j,G_o\theta _o\mathrm{}\theta _{j1}]`$ and $`Rs_i=[N_i,G_o\theta _o\mathrm{}\theta _{i1}]`$, for which requirements i) and ii) of Definition 6.10 hold, we will write $`Rs_iRs_j`$. In other words, Definition 6.10 expresses that p-$`EVR_L`$ loop check is based on detecting that a reduced resultant is obtained which is connected by the relationship $``$ to a preceding one in the same derivation. It is worth noting that $``$ is an equivalence relationship. Now let us prove Theorem 6.13, which states that the completeness of p-$`EVR_L`$ loop check is preserved passing from p-SLD to p-RSLD, if stack-queue scheduling rules are used. To this aim we provide a necessary condition which holds whenever p-$`EVR_L`$ prunes every infinite p-SLD derivation of a goal$`G`$in a program P via a scheduling rule$`S.`$ Indeed, as shown in Lemma 6.11, in this hypothesis the length of resolvents of all possible derivations of$`G`$in P via$`S`$is limited. The structure and the proof of Lemma 6.11 are strictly analogous to the ones of Lemma 2.4. Note also that Lemma 6.11 holds for any scheduling rule. On the contrary, the stack-queue hypothesis is necessary in Theorem 6.13, which concludes the section. ###### Lemma 6.11. Let P be a program and$`G`$a p-goal. Suppose that all infinite p-SLD derivations of$`G`$in P via a scheduling rule$`S`$are pruned by p-$`EVR_L`$. Then, a finite bound $`l`$ exists such that, for each resolvent$`R`$in any p-SLD derivation of$`G`$in P via$`S,`$ it is $`\mathrm{\#}Rl`$. ###### Proof 6.12. The proof of this lemma can be obtained from the one of Lemma 2.4, by means of the following replacements: “Let $`T`$ be the p-SLD tree of$`G`$in P via $`S`$” for “Let $`T`$ be an S-tree of$`G`$in P”, “By Determinism Lemma 4.10” for “Since $`T`$ contains all SLD derivations of$`G`$in P”, “p-$`EVR_L`$” and “p-variant” for “$`EVR_L`$” and “variant”, respectively. ###### Theorem 6.13 (p-$`EVR_L`$ loop check completeness preservation). Let P be a function free program, $`G_o`$ a p-goal and $`SQ`$ a stack-queue scheduling rule. Suppose that all infinite p-SLD derivations of $`G_o`$ in P via $`SQ`$ are pruned by p-$`EVR_L`$, then all infinite p-RSLD derivations of $`G_o`$ in P via $`SQ`$ are pruned by p-$`EVR_L`$. ###### Proof 6.14. Let $`D`$ be an infinite p-RSLD derivation of $`G_o`$ in P via $`SQ`$. Let $`(G_o\stackrel{SQ,X}{}>>Q)`$ be any finite prefix of $`D`$. By Lemma 6.4, a p-SLD derivation $`D^{}=(G_o\stackrel{SQ,Z}{}R)`$ exists with $`\mathrm{\#}Q\mathrm{\#}R`$. On the other hand, by Lemma 6.11 a bound $`l`$ exists such that $`\mathrm{\#}Q\mathrm{\#}Rl`$. But $`Q`$ is the generic reduced resolvent in $`D`$, so that the number of atoms in all reduced resolvents of $`D`$ is bounded by $`l`$. As a consequence, the number of atoms in all reduced resultants of $`D`$ is also limited. Since the program P has finite many predicate symbols and constants and no function symbol is allowed, the relationship $``$ between reduced resultants of $`D`$ has only finitely many equivalence classes. Then, for some $`0i<k`$ in $`D`$, we have that the $`k^{th}`$ reduced resultant is related by $``$ to the $`i^{th}`$ one. This implies that $`D`$ is pruned by p-$`EVR_L`$. ## 7 Conclusions In the paper the problem of possible undesirable effects of redundancy elimination from resolvents is addressed. In particular we have shown that program termination and loop check completeness can be lost. Conditions are characterised which ensure the redundancy elimination tolerance, in the sense that program termination and completeness of equality loop check are preserved when redundancy is eliminated. However, difficulties in analysing interdependence of redundancy elimination effects from the used selection rule have arisen, and the necessity of a framework to formalise suitable features of selection rules has been highlighted. To this aim, a highly expressive execution model based on priority mechanism for atom selection is developed in the paper. The distinctive aspect is that primary importance is given to the event of arrival of new atoms from the body of the applied clause at rewriting time, when new atoms can be freely positioned with respect to old ones in the resolvent. Then, at any derivation step, the atom with optimum priority is simply selected. The results presented in the paper show that the new computational model is able to give remarkable insights into general properties of selection rules. As a matter of fact, the priority model allows us to formalise the delicate concepts on which the axiomatic definition of specialisation independent scheduling rules is based. As a quite unexpected result, the specialisation independence turns out to be equivalent to stack-queue scheduling technique, which has a very simple operational characterisation. In other words, the priority mechanism is necessary to formalise the real semantic features of specialisation independent scheduling rules. On the contrary, the full generality of the same mechanism can be abandoned if only operational aspects of specialisation independent rules are of interest, in the sense that all we need is a “watershed” between the stacked and the queued atoms. It is widely acknowledged that the study of selection rules is a difficult subject which deserves attention. We are confident that the computational model proposed in the paper can be usefully exploited in future work to get further insights into topics which are related to selection rule theory and application, such as loop check, termination and optimisation of derivation processes. ## Appendix A Appendix This Appendix contains the formal proofs of Properties 3.3 and 4.1. The very simple Property A1 is considered before proving Property 3.3. ###### Property A.1 Let$`S`$be a complete set of derivation steps. Given a p-goal$`G`$and a clause $`c`$, the following implication holds: $`Ds`$ derivation step of the type $`(G\stackrel{c\xi }{}),`$(1) $`Ds^{}`$ of the type $`(G\stackrel{c\xi }{})`$, with $`Ds^{}S`$. ###### Proof A.1. Let$`G=a\|K`$ and $`c=(htB)`$. By (1) and the completeness of$`S`$(part i), a derivation step exists of the form: $`(a\|K\stackrel{𝑐}{}(K+B\xi ^{}\underset{¯}{\theta })\mu ^{})S`$(2) By definition, the derivation step in (1) has the form: $`a\|K\stackrel{𝑐}{}(K+B\xi \underset{¯}{\tau })\mu .`$ Then, it is evident that a derivation also exists like: $`Ds^{}=(a\|K\stackrel{𝑐}{}(K+B\xi \underset{¯}{\theta })\mu ).`$ By construction, derivation steps (2) and $`Ds^{}`$ are congruent lowerings of each other. Then, by completeness of$`S`$(part ii), derivation step $`Ds^{}`$ belongs to$`S`$. ###### Property A.2 (Property 3.3) Let$`S`$be a complete set of derivation steps. Given two p-goals $`a\gamma \underset{¯}{\tau }\|G`$ and $`a\|F`$, let us fix arbitrarily a finite set $`V`$of variables. The following implication holds: $`Ds`$ derivation step of the form $`a\gamma \underset{¯}{\tau }\|G\stackrel{𝑐}{}`$(1) $`Ds^{}`$ of the form $`a\|F\stackrel{𝑐}{}`$, with $`Ds^{}S`$and $`nvar(Ds^{})V=\mathrm{}.`$ ###### Proof A.2. Let$`c=(htB)`$. On the basis of (1), by definition of derivation step, a standardisation apart renaming $`\xi ^{}`$ for$`c`$and an mgu $`\beta `$ exist, with $`a\gamma \beta =(ht)\xi ^{}\beta `$. Then, let us consider a renaming $`\xi `$ of $`c\xi ^{}`$, such that the following assertions hold for the range of $`\xi `$: $`var(a\|F)var(c\xi ^{}\xi )=\mathrm{},`$ (2a) $`domain(\gamma )var((ht)\xi ^{}\xi )=\mathrm{},`$ (2b) $`domain(\xi ^1)var(a\gamma )=\mathrm{}`$ (2c) $`var(c\xi ^{}\xi )V=\mathrm{}.`$ (2d) By facts (2b) and (2c), we have that: $`a\gamma \xi ^1\beta =^{(2c)}a\gamma \beta =(ht)\xi ^{}\beta =(ht)\xi ^{}\xi \xi ^1\beta =^{(2b)}(ht)\xi ^{}\xi \gamma \xi ^1\beta .`$ In other words, $`a`$ and $`(ht)\xi ^{}\xi `$ unify through the unifier $`\gamma \xi ^1\beta `$. On the other hand, the fact (2a) says that $`\xi ^{}\xi `$ is a standardisation apart renaming for$`c`$with respect to $`a\|F`$. Then, a derivation step exists of the form $`a\|F\stackrel{c\xi ^{}\xi }{}`$. By hypothesis the set$`S`$is complete, so that by Property A1 we have also a derivation step such that: $`Ds^{}=(a\|F\stackrel{c\xi ^{}\xi }{})S.`$ Since it is $`nvar(Ds^{})=var(c\xi ^{}\xi )`$, by (2d) we have that $`nvar(Ds^{})V=\mathrm{}`$. ###### Property A.3 (Property 4.1) Let$`c=(htB)`$ be a clause. Let us consider two derivation steps $`Ds_1`$ and $`Ds_2`$ such that the $`Ds_2`$ is a lowering of $`Ds_1`$ by $`X`$. The following implication holds: $`Ds_1=(a\|K\stackrel{𝑐}{}(K+B\xi ^{}\underset{¯}{\theta }^{})\mu ^{}),`$ (1) $`Ds_2=(a\tau \underset{¯}{\sigma }\|(K\tau \underset{¯}{\sigma }+X)\stackrel{𝑐}{}(K\tau \underset{¯}{\sigma }+B\xi ^{\prime \prime }\underset{¯}{\theta }^{\prime \prime }+X)\mu ^{\prime \prime })`$ (2) $`\delta `$such that $`K\tau \mu ^{\prime \prime }=K\mu ^{}\delta ,B\xi ^{\prime \prime }\mu ^{\prime \prime }=B\xi ^{}\mu ^{}\delta ,`$ where $`\delta `$ is a renaming, if $`\tau `$ is a renaming. ###### Proof A.3. By definition of derivation step, we have: $`var(a\|K)var((htB)\xi ^{})=\mathrm{},`$(3) $`var((a\|K)\tau )var((htB)\xi ^{\prime \prime })=\mathrm{}`$, (4) $`\mu ^{}=mgu(a,(ht)\xi ^{}),\mu ^{\prime \prime }=mgu(a\tau ,(ht)\xi ^{\prime \prime })`$. (5) Let $`\pi =\tau /var(a\|K)^[`$<sup>3</sup><sup>3</sup>3The notation $`\tau /var(a\|K)`$ represents $`\tau `$ restricted to the variables of $`a\|K`$.<sup>\]</sup> and $`\varphi =((\xi ^{})^1\xi ^{\prime \prime })/var((htB)\xi ^{})`$. By (3) it is: $`domain(\pi )domain(\varphi )=\mathrm{}`$, (6a) $`(htB)\xi ^{}\pi =(htB)\xi ^{}`$ and $`(a\|K)\varphi =(a\|K)`$. (6b) As a consequence of (6a), the union $`(\pi \varphi )`$ is a well defined substitution. Then, we may write that: $`a(\pi \varphi )\mu ^{\prime \prime }=^{(6b)}a\pi \mu ^{\prime \prime }=a\tau \mu ^{\prime \prime }=^{(5)}(ht)\xi ^{\prime \prime }\mu ^{\prime \prime }=(ht)\xi ^{}(\xi ^{})^1\xi ^{\prime \prime }\mu ^{\prime \prime }=`$ $`=(ht)\xi ^{}\varphi \mu ^{\prime \prime }=^{(6b)}(ht)\xi ^{}(\pi \varphi )\mu ^{\prime \prime },`$ so that $`(\pi \varphi )\mu ^{\prime \prime }`$ is an unifier of $`a`$ and $`(ht)\xi ^{}`$. Since $`\mu ^{}`$ is an mgu of $`a`$ and $`(ht)\xi ^{}`$, a substitution $`\delta `$ exists with: $`(\pi \varphi )\mu ^{\prime \prime }=\mu ^{}\delta `$. (7) Then, we have: $`K\tau \mu ^{\prime \prime }=K\pi \mu ^{\prime \prime }=^{(6b)}K(\pi \varphi )\mu ^{\prime \prime }=^{(7)}K\mu ^{}\delta ,`$ (8a) $`B\xi ^{\prime \prime }\mu ^{\prime \prime }=B\xi ^{}(\xi ^{})^1\xi ^{\prime \prime }\mu ^{\prime \prime }=B\xi ^{}\varphi \mu ^{\prime \prime }=^{(6b)}B\xi ^{}(\pi \varphi )\mu ^{\prime \prime }=^{(7)}B\xi ^{}\mu ^{}\delta .`$ (8b) Now let us suppose that $`\tau `$ is a renaming. In this case, facts (3) and (4) become symmetric at all. As a consequence, by symmetry with respect to (8a) and (8b), a substitution $`\gamma `$ exists such that $`K\mu ^{}=K\tau \mu ^{\prime \prime }\gamma `$ and $`B\xi ^{}\mu ^{}=B\xi ^{\prime \prime }\mu ^{\prime \prime }\gamma `$. Then we have: $`(K\mu ^{}+B\xi ^{}\mu ^{})\delta \gamma =(K\tau \mu ^{\prime \prime }+B\xi ^{\prime \prime }\mu ^{\prime \prime })\gamma =K\mu ^{}+B\xi ^{}\mu ^{}.`$ It is evident that $`\delta `$ is a renaming for $`K\mu ^{}+B\xi ^{}\mu ^{}`$, then the thesis is verified. ## Appendix B Appendix In this Appendix we provide a formal proof of the duplication theorem (Theorem 5.10). Such a proof exploits two lemmata which are given below. Lemma B1 establishes a condition which allows us to repeat derivations via a specialisation independent scheduling rule, when we pass from a goal$`G`$to a suitable kind of instantiations of $`G`$. Lemma B1 is a correspondent, for p-SLD derivations, of part (ii) of Strong Lifting Lemma \[Gabrielli, Levi and Meo, 1996\]. Indeed, both part (ii) of Strong Lifting Lemma and Lemma B1 can be seen as results about sufficient conditions for derivation lowering from a goal$`G`$to instantiations of$`G`$itself. Here a direct proof of Lemma B1 is given which takes into account technical aspects concerning our priority value mechanism. Lemma B1 does not relate resolvents, because it is not important for the purposes of this Appendix. ###### Lemma B.1. Let$`S`$be a specialisation independent scheduling rule,$`G`$a p-goal and $`\varphi `$ a substitution. The following implication holds: $`G\stackrel{S,X,\theta }{}G\theta \varphi \stackrel{S,X}{}.`$(1) ###### Proof B.2. The proof is by induction on the length of $`X`$. If $`\mathrm{\#}X=0`$, the thesis is trivially true. For $`\mathrm{\#}X>0`$, let$`G=a\|F`$, $`X=c\|H`$ with$`c=(htB)`$, and rewrite derivation (1) as follows: $`a\|F\stackrel{S,c\xi ,\gamma }{}(Q=(F+B\xi \underset{¯}{\pi })\gamma )\stackrel{S,H,\mu }{},`$ where $`\gamma \mu =\theta `$. (2) Then, let us consider the substitution $`\varphi _g=\varphi \sigma _g,`$ where $`\sigma _g`$ is such that $`(a\|F)\theta \varphi _g=G\theta \varphi _g`$ is ground. Since $`\gamma `$ is an mgu of $`a`$ and $`(ht)\xi `$, we have $`a\gamma =(ht)\xi \gamma `$, which means $`a\theta \varphi _g=a\gamma \mu \varphi _g=(ht)\xi \gamma \mu \varphi _g=(ht)\xi \theta \varphi _g.`$ But $`a\theta \varphi _g`$ is ground, so that we obtain the equality $`(a\theta \varphi _g)\theta \varphi _g=a\theta \varphi _g=((ht)\xi )\theta \varphi _g`$. In other words, $`a\theta \varphi _g`$ and $`(ht)\xi `$ unify through the unifier $`\theta \varphi _g`$. Moreover, the renamed clause $`c\xi `$ is obviously standardised apart with respect to the ground p-goal $`(a\|F)\theta \varphi _g`$, so that a derivation step like $`(a\|F)\theta \varphi _g\stackrel{c\xi }{}`$ exists. Thus, by completeness of$`S`$and Property A1, a derivation step also exists of the form: $`(G\theta \varphi _g=(a\|F)\theta \varphi _g)\stackrel{S,c\xi ,\eta }{}(R=(F\theta \varphi _g+B\xi \underset{¯}{\pi }^{})\eta ).`$ (3) Now, the substitution $`\eta `$ is an mgu of $`a\theta \varphi _g`$ and $`(ht)\xi `$, so that a substitution $`\pi `$ exists with: $`\theta \varphi _g=\eta \pi .`$ (4) On the other hand, since$`S`$is specialisation independent, step (3) is a congruent lowering of the first step of (2) by $`\mathrm{}`$, i.e. $`\underset{¯}{\rho }`$ such that $`F\underset{¯}{\rho }=F,B\underset{¯}{\pi }\underset{¯}{\rho }=B\underset{¯}{\pi }^{}`$, (5) which implies: $`Q\mu \varphi _g=(F+B\xi \underset{¯}{\pi })\gamma \mu \varphi _g\underset{¯}{\rho }\underset{¯}{\rho }^1=(F\underset{¯}{\rho }+B\xi \underset{¯}{\pi }\underset{¯}{\rho })\theta \varphi _g\underset{¯}{\rho }^1=^{(5)}(F+B\xi \underset{¯}{\pi }^{})\theta \varphi _g\underset{¯}{\rho }^1`$. But $`F\theta \varphi _g`$ is ground, so that $`(F\theta \varphi _g)\theta \varphi _g=F\theta \varphi _g`$. As a consequence: $`Q\mu \varphi _g=(F\theta \varphi _g+B\xi \underset{¯}{\pi }^{})\theta \varphi _g\underset{¯}{\rho }^1=^{(4)}(F\theta \varphi _g+B\xi \underset{¯}{\pi }^{})\eta \pi \underset{¯}{\rho }^1=R\pi \underset{¯}{\rho }^1.`$ By inductive hypothesis applied to the tail of (2), we have that $`(Q\mu \varphi _g=R\pi \underset{¯}{\rho }^1)\stackrel{S,H}{},`$ which by Lifting Lemma 4.7 implies that$`R\stackrel{S,H}{}`$. Now, by Property 4.3, the last obtained derivation can be combined with (3) yielding: $`(G\theta \varphi _g=G\theta \varphi \sigma _g)\stackrel{S,X}{}.`$ By Lifting Lemma 4.7, we conclude $`(G\theta \varphi \stackrel{S,X}{})`$, so that the inductive step is completed. The following Lemma B2 is a special form of determinism lemma which holds for preq type stack-queue derivations. Roughly speaking, the lemma states that an $`A`$-preq type derivation, starting from a p-goal of the form $`A\|X`$, can be replicated from a p-goal like $`A\lambda \underset{¯}{\lambda }\|Y`$, where $`\lambda `$ is a renaming. Note that no hypothesis is made on $`X`$ and $`Y`$ which can be completely unrelated. The intuitive explication is that only atoms deriving from $`A`$ are rewritten so that neither $`X`$ nor $`Y`$ have any active role in the derivations. The formal statement and the proof of Lemma B2 are preceded by the quite simple Property B1. ###### Property B.1 Let $`SQ`$ be a stack-queue scheduling rule. The following implication holds: $`A\|X\stackrel{SQ,D}{}Q`$, of $`A`$-preq type(1) $`A\gamma \underset{¯}{\lambda }\|Y\stackrel{SQ,D}{}R`$, of $`(A\gamma \underset{¯}{\lambda })`$-preq type, where $`\gamma `$ is a renaming(2) $``$ $`\delta ,\underset{¯}{\delta }`$ such that $`R/(A\gamma \underset{¯}{\lambda })=(Q/A)\delta \underset{¯}{\delta },`$ where $`\delta `$ is a renaming. ###### Proof B.3. By hypothesis, derivations (1) and (2) are of $`A`$-preq and $`(A\gamma \underset{¯}{\lambda })`$-preq type respectively, so that $`D/A=D/(A\gamma \underset{¯}{\lambda })=D`$. Then, by Lifting Lemma 4.7, a derivation exists like: $`A\stackrel{S,D}{}T`$. (3) By Lowering Lemma 4.4, applied to (3) and (1), a renaming $`\alpha `$ and a shifting $`\underset{¯}{\alpha }`$ exist with $`Q/A=T\alpha \underset{¯}{\alpha }`$. By Lowering Lemma 4.4 applied to (3) and (2), a renaming $`\beta `$ and a shifting $`\underset{¯}{\beta }`$ exist with $`R/(A\gamma \underset{¯}{\lambda })=T\beta \underset{¯}{\beta }`$. Finally, we derive that: $`R/(A\gamma \underset{¯}{\lambda })=T\alpha \underset{¯}{\alpha }\alpha ^1\underset{¯}{\alpha }^1\beta \underset{¯}{\beta }=(Q/A)\alpha ^1\underset{¯}{\alpha }^1\beta \underset{¯}{\beta }`$. ###### Lemma B.4 (preq type determinism). Let $`SQ`$ be a stack-queue scheduling rule and $`V`$ any finite set of variables. Let $`A\|X`$ and $`A\lambda \underset{¯}{\lambda }\|Y`$ two p-goals, where $`\lambda `$ is a renaming. The following implication holds: $`A\|X\stackrel{SQ,K,\psi }{}A^s\|X\psi \|A^q`$, of $`A`$-preq type, (1) $``$ $`\delta ,\underset{¯}{\delta }`$ and $`D=(A\lambda \underset{¯}{\lambda }\|Y\stackrel{SQ,K,\theta }{}A^s\delta \underset{¯}{\delta }\|Y\theta \|A^q\delta \underset{¯}{\delta })`$, of $`(A\lambda \underset{¯}{\lambda })`$-preq type, where $`\delta `$ is a renaming and $`nvar(D)V=\mathrm{}.`$ ###### Proof B.5. Let $`A=a\|F`$. We show Lemma B2 by induction on the length of the template $`K`$. If $`\mathrm{\#}K=0`$, the assert is evident. If $`\mathrm{\#}K>0`$, let $`K=c\|H`$ with$`c=(htM_s\|M_q)`$. Derivation (1) can be rewritten as follows: $`a\|F\|X\stackrel{SQ,c\xi ,\mu }{}(Q=a^s\|F\mu \|X\mu \|a^q)\stackrel{SQ,H,\sigma }{}A^s\|X\psi \|A^q`$ (1a) with $`a^s=M_s\xi \underset{¯}{\alpha }\mu `$ and $`A^s=(a^s\|F\mu )^s`$. (1b) Then, by Property 3.3 with reference to the first step of (1a), a derivation step $`Ds`$ exists such that: $`Ds=((a\|F)\lambda \underset{¯}{\lambda }\|Y\stackrel{SQ,c\tau ,\eta }{}R)`$, with $`nvar(Ds)V=\mathrm{}`$. (2) Since the selection rule $`SQ`$ is stack queue, it must be: $`R=M_s\tau \underset{¯}{\gamma }\eta \|F\lambda \underset{¯}{\lambda }\eta \|Y\eta \|M_q\tau \underset{¯}{\gamma }\eta `$. By Property B1 applied to derivation step (2) and the first step of (1a), a renaming $`\beta `$ and a shifting $`\underset{¯}{\beta }`$ exist with: $`M_s\tau \underset{¯}{\gamma }\eta \|F\lambda \underset{¯}{\lambda }\eta \|M_q\tau \underset{¯}{\gamma }\eta =R/^2/(A\lambda \underset{¯}{\lambda })=(Q/^{1a}/A)\beta \underset{¯}{\beta }=(a^s\|F\mu \|a^q)\beta \underset{¯}{\beta }.`$ Now, by (1b) it is $`\mathrm{\#}a^s=\mathrm{\#}M_s\tau \underset{¯}{\gamma }\eta ,`$ so that the equality $`M_s\tau \underset{¯}{\gamma }\eta \|F\lambda \underset{¯}{\lambda }\eta =(a^s\|F\mu )\beta \underset{¯}{\beta }`$ holds, by Property 3.1-i). Thus, the inductive hypothesis can be applied to the tail of (1a). As a consequence, a p-SLD derivation $`D^{}`$ exists, which is of $`(M_s\tau \underset{¯}{\gamma }\eta \|F\lambda \underset{¯}{\lambda }\eta )`$-preq type and has the following form: $`R\stackrel{SQ,H,\pi }{}((a^s\|F\mu )^s\delta ^{}\underset{¯}{\delta }^{}\|Y\eta \pi \|Z=^{(1b)}A^s\delta ^{}\underset{¯}{\delta }^{}\|Y\eta \pi \|Z)`$, (3) with $`nvar(D^{})(var((a\|F)\lambda \underset{¯}{\lambda }\|Y)nvar(Ds)V)=\mathrm{}`$. (3a) On the basis of (3a) above, derivation step (2) and derivation (3) can be combined to yield the derivation $`D`$: $`D=(A\lambda \underset{¯}{\lambda }\|Y\stackrel{SQ,K}{}A^s\delta ^{}\underset{¯}{\delta }^{}\|Y\eta \pi \|Z)`$, (4) where $`D`$ is of $`(A\lambda \underset{¯}{\lambda })`$-preq type. The thesis is now proven. Indeed, by Property B1 applied to derivations (1) and (4), a renaming $`\delta `$ and a shifting $`\underset{¯}{\delta }`$ exist with $`A^s\delta ^{}\underset{¯}{\delta }^{}\|Z=(A^s\|A^q)\delta \underset{¯}{\delta }`$, so that by Property 3.1-i) we have $`A^s\delta ^{}\underset{¯}{\delta }^{}=A^s\delta \underset{¯}{\delta }`$ and $`Z=A^q\delta \underset{¯}{\delta }`$. The fact that $`nvar(D)V=\mathrm{}`$ follows from (2) and (3a). ###### Theorem B.6 (Theorem 5.10 – duplication theorem). Let $`SQ`$ be a stack-queue scheduling rule. Given two p-goals of the form $`A\|B\|C\|D`$ and $`A\|B\|C\|B\underset{¯}{\pi }\|D`$, the following implication holds: $`(A\|B\|C\|D\stackrel{SQ,X.P}{}Q)`$(1) $``$ $`Y`$ such that $`(A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,Y.P}{}R)`$ with $`X_LY`$ and $`\mathrm{\#}Q\mathrm{\#}R.`$ ###### Proof B.7. Let $`\mathrm{\Delta }(SQ,n)`$ denote the subset of $`\mathrm{\Delta }(SQ)`$, such that for any derivation $`Dr`$ in $`\mathrm{\Delta }(SQ,n)`$ it is $`\mathrm{\#}Drn`$, where $`\mathrm{\#}Dr`$ denotes the length of $`Dr`$. We show the thesis by induction on $`n`$, i.e. we show that the thesis holds when derivation (1) belongs to $`\mathrm{\Delta }(SQ,n)`$, for any $`n0`$. The fact is obvious for $`\mathrm{\Delta }(SQ,0)`$. In order to prove the inductive step from $`\mathrm{\Delta }(SQ,n1)`$ to $`\mathrm{\Delta }(SQ,n)`$, for $`n>0`$, let us consider a derivation like: $`(A\|B\|C\|D\stackrel{X.P,\theta }{}Q)\mathrm{\Delta }(SQ,n)`$, (1a) and show that $`(A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,Y.P}{}R)`$ exists with $`X_LY`$ and $`\mathrm{\#}Q\mathrm{\#}R`$. Actually, the proof of the inductive step will be organised in two phases: * first, the inductive step is shown in the case that the initial p-goal $`A\|B\|C\|D`$ is ground, * then, the validity of the inductive step is extended to generic initial p-goals. Let us recall that the sketch of Section 5.3 was given in the simplifying hypothesis that: every clause body introduces no new variable and initial p-goals are ground. In this sense, we may say that the first phase removes the first restriction, while the second one is retained. In the second phase, also the restriction on the groundness of initial goals is overcome. First phase (the initial p-goal $`A\|B\|C\|D`$ is ground). With reference to (1a), the following three possible situations must be taken into account. Then, we start with Case 3, which is the most significant one. 1. derivation (1a) is of $`(A\|B\|C)`$-preq type, 2. derivation (1a) is of $`(A\|B\|C\|D)`$-preq type, and not of $`(A\|B\|C)`$-preq type, 3. derivation (1a) is not of $`(A\|B\|C\|D)`$-preq type. Case3. Derivation (1a) has the following form: $`A\|B\|C\|D\stackrel{𝐻}{}B\|C\|D\|A^q\stackrel{𝐾}{}C\|D\|A^q\|B^q\stackrel{𝑀}{}`$ $`D\|A^q\|B^q\|C^q\stackrel{𝑁}{}A^q\|B^q\|C^q\|D^q\stackrel{𝑇}{}Q,`$ (2) where $`H\|K\|M\|N\|T=X`$. In fact, since $`A\|B\|C\|D`$ is ground, in each of the four initial segments of (2) only standardisation apart variables, introduced in the same segment, can be instantiated. In particular, $`var(A^q)`$ are not instantiated in the second segment, $`var(A^q\|B^q)`$ are not in the third segment, and $`var(A^q\|B^q\|C^q)`$ are not in the fourth one. It is evident, as a consequence, that the p-goal $`A^q\|B^q\|C^q\|D^q`$ consists of four subgoals without common variables. Now, since also $`B\underset{¯}{\pi }`$ is ground, a derivation can be constructed through five successive applications of Lemma B2, as depicted below: $`A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,H}{}B\|C\|B\underset{¯}{\pi }\|D\|A^q\alpha \underset{¯}{\alpha }\stackrel{SQ,K}{}`$ $`C\|B\underset{¯}{\pi }\|D\|A^q\alpha \underset{¯}{\alpha }\|B^q\beta \underset{¯}{\beta }\stackrel{SQ,M}{}B\underset{¯}{\pi }\|D\|A^q\alpha \underset{¯}{\alpha }\|B^q\beta \underset{¯}{\beta }\|C^q\gamma \underset{¯}{\gamma }\stackrel{SQ,K}{}`$ (3) $`D\|A^q\alpha \underset{¯}{\alpha }\|B^q\beta \underset{¯}{\beta }\|C^q\gamma \underset{¯}{\gamma }\|B^q\varphi \underset{¯}{\varphi }\stackrel{SQ,N}{}(Z=A^q\alpha \underset{¯}{\alpha }\|B^q\beta \underset{¯}{\beta }\|C^q\gamma \underset{¯}{\gamma }\|B^q\varphi \underset{¯}{\varphi }\|D^q\delta \underset{¯}{\delta })`$ where $`\alpha `$, $`\beta `$, $`\gamma `$, $`\varphi `$ and $`\delta `$ are renamings. At each application of Lemma B2, a segment of derivation (3) is obtained on the basis of a corresponding segment of derivation (2). Moreover, Lemma B2 assures that each new segment can be freely standardised apart, so that each segment can be readily added to the sequence of its predecessors in (3). Note that the second segment of (2) is considered twice, in order to generate both the second and the fourth segment of (3). In analogy with derivation (2), the final p-goal $`Z`$ of derivation (3) consists of five subgoals without common variables. As a consequence, the five renamings $`\alpha ^1`$, $`\beta ^1`$, $`\gamma ^1`$, $`\varphi ^1`$ and $`\delta ^1`$ have disjoint domains, so that they can be joined in order to form a unique substitution $`\xi =(\alpha ^1\beta ^1\gamma ^1\varphi ^1\delta ^1)`$. (4a) Then, let us consider the p-goal $`A^q\|B^q\|C^q\|B^q\underset{¯}{\pi }^{}\|D^q`$, where $`\underset{¯}{\pi }^{}`$ is a suitable shifting. By (4a) and Property 3.2, we have that: $`A^q\|B^q\|C^q\|B^q\underset{¯}{\pi }^{}\|D^q=^{(Axiii)}(A^q\underset{¯}{\alpha }\|B^q\underset{¯}{\beta }\|C^q\underset{¯}{\gamma }\|B\underset{¯}{\varphi }\|D^q\underset{¯}{\delta })\underset{¯}{\sigma }=^{(4a)}`$ $`(A^q\alpha \underset{¯}{\alpha }\|B^q\beta \underset{¯}{\beta }\|C^q\gamma \underset{¯}{\gamma }\|B^q\varphi \underset{¯}{\varphi }\|D^q\delta \underset{¯}{\delta })\xi \underset{¯}{\sigma }=Z\xi \underset{¯}{\sigma }`$. (4b) By construction of (2) the derivation $`(A^q\|B^q\|C^q\|D^q\stackrel{𝑇}{}Q)`$ belongs to $`\mathrm{\Delta }(SQ,m)`$, with $`m<n`$. By inductive hypothesis, a derivation exists of the form: $`(Z\xi \underset{¯}{\sigma }=A^q\|B^q\|C^q\|B^q\underset{¯}{\pi }^{}\|D^q)\stackrel{SQ,Y^{}.P}{}R^{}`$, (5) with $`T_LY^{}`$ and $`\mathrm{\#}Q\mathrm{\#}R^{}`$. (5a) By (5) and Lifting Lemma 4.7 a derivation exists like: $`Z\stackrel{SQ,Y^{}}{}`$, (6) By Property 4.3, derivations (3) and (6) can be combined to yield a derivation of the form: $`A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,(H\|K\|M\|K\|N).P}{}Z\stackrel{SQ,Y^{}.P}{}R`$, (7) where, by Lowering Lemma 4.4 applied to (5) and the tail of (7), it is $`\mathrm{\#}R=\mathrm{\#}R^{}`$. Finally: $`X=H\|K\|M\|N\|T_L^{(5a)}(H\|K\|M\|K\|N\|Y^{})`$ and $`\mathrm{\#}Q^{(5a)}\mathrm{\#}R^{}=\mathrm{\#}R`$. Case 2. Derivation (1a) has the following form: $`A\|B\|C\|D\stackrel{H\|K\|M\|N}{}D^s\|A^q\|B^q\|C^q\|D^q`$, with $`H\|K\|M\|N=X`$. Analogously to preceding case 3), through Lemma B2 a derivation can be constructed like: $`A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,(H\|K\|M\|K\|N).P}{}D^s\delta \underset{¯}{\delta }\|A^q\alpha \underset{¯}{\alpha }\|B^q\beta \underset{¯}{\beta }\|C^q\gamma \underset{¯}{\gamma }\|B^q\varphi \underset{¯}{\varphi }\|D^q\delta \underset{¯}{\delta }`$. Case 1. Derivation (1a) has the form: $`A\|B\|C\|D\stackrel{𝑋}{}(A\|B\|C)^s\|D\|(A\|B\|C)^q.`$ Through Lemma B2, a derivation can be constructed like: $`(A\|B\|C)\|B\underset{¯}{\pi }\|D\stackrel{SQ,X}{}(A\|B\|C)^s\gamma \underset{¯}{\gamma }\|B\underset{¯}{\pi }\|D\|(A\|B\|C)^q\gamma \underset{¯}{\gamma }.`$ Second phase (the initial p-goal $`A\|B\|C\|D`$ is generic). In the preceding first phase of this proof, the inductive step is verified in the hypothesis that the initial p-goal $`A\|B\|C\|D`$ is ground. Now consider a generic p-goal of the form $`A\|B\|C\|D`$. With reference to (1a), let $`\varphi _g`$ be a grounding substitution for $`(A\|B\|C\|D)\theta `$. By Lemma B1, a derivation exists such that: $`((A\|B\|C\|D)\theta \varphi _g\stackrel{𝑋}{}Q^{})\mathrm{\Delta }(SQ,n)`$, (8) where, by Lowering Lemma 4.4 applied to (1a) and (8), we have: $`\mathrm{\#}Q^{}=\mathrm{\#}Q`$. (8a) Since the inductive hypothesis is already proven for ground initial goals, by (8) a derivation exists: $`(A\|B\|C\|B\underset{¯}{\pi }\|D)\theta \varphi _g\stackrel{SQ,Y.P}{}R^{}`$, (9) with $`X_LY`$ and $`\mathrm{\#}Q^{}\mathrm{\#}R^{}`$. (9a) Then, by Lifting Lemma 4.7 a derivation exists: $`A\|B\|C\|B\underset{¯}{\pi }\|D\stackrel{SQ,Y.P}{}R`$, (10) where, by Lowering Lemma 4.4 applied to (9) and (10), we have $`X_L^{(9a)}Y`$ and $`\mathrm{\#}Q=^{(8a)}\mathrm{\#}Q^{}^{(9a)}\mathrm{\#}R^{}=^{(Lem.\text{4.4})}\mathrm{\#}R`$. As a consequence, the induction step is completely verified.
warning/0004/cond-mat0004486.html	ar5iv	text	# Slow Logarithmic Decay of Magnetization in the Zero Temperature Dynamics of an Ising Spin Chain: Analogy to Granular Compaction ## Abstract We study the zero temperature coarsening dynamics in an Ising chain in presence of a dynamically induced field that favors locally the ‘$``$’ phase compared to the ‘$`+`$’ phase. At late times, while the ‘$`+`$’ domains still coarsen as $`t^{1/2}`$, the ‘$``$’ domains coarsen slightly faster as $`t^{1/2}\mathrm{log}(t)`$. As a result, at late times, the magnetization decays slowly as, $`m(t)=1+\mathrm{const}./\mathrm{log}(t)`$. We establish this behavior both analytically within an independent interval approximation (IIA) and numerically. In the zero volume fraction limit of the ‘$`+`$’ phase, we argue that the IIA becomes asymptotically exact. Our model can be alternately viewed as a simple Ising model for granular compaction. At late times in our model, the system decays into a fully compact state (where all spins are ‘$``$’) in a slow logarithmic manner $`1/\mathrm{log}(t)`$, a fact that has been observed in recent experiments on granular systems. PACS numbers: 05.40.-a, 82.20.Mj The effect of quenched disorder on the relaxation dynamics of many body systems has been studied quite extensively. In systems such as structural glasses, where quenched disorder is absent, an alternative approach has been put forward that considers the slow relaxation due to kinetic disorder, induced by the dynamics itself. Another important system where kinetic disorders give rise to slow relaxation is granular material. The density relaxation of loosely packed glass beads has been studied in recent experiments and it was found that the density $`\rho (t)`$ compactified slowly as, $`\rho (\mathrm{})\rho (t)1/\mathrm{log}(t)`$, under mechanical tapping. It is natural to expect that such kinetic disorders may play important roles in the dynamics of other systems as well. In this Letter we study, for the first time, the effect of a dynamically self-induced field in an important class of out of equilibrium problems, namely the domain growth problems, and show that such systems also exhibit logarithmic relaxation, suggesting that inverse logarithmic relaxation is quite robust. Domain growth following a rapid quench in temperature in ferromagnetic spin systems is one of the better understood out of equilibrium phenomena. For example, if an Ising system is quenched rapidly from infinite temperature to zero temperature without breaking the symmetry between the two ground states, domains of up and down spins form and grow with time. The average linear size of a domain grows with time as $`l(t)t^{\frac{1}{2}}`$ for zero temperature nonconserved dissipative dynamics. However if one puts on a small uniform external field (say in the down direction), then even at zero temperature the system quickly reaches the pure state of magnetization $`1`$ in a finite time proportional to the initial size of the up domains. A natural question is: what happens when, instead of applying a global external bias, the symmetry between the pure states is broken locally by the dynamics itself ? In this Letter we address this question in the context of a simple Ising spin chain with spins $`S_i=\pm 1`$. Starting from a given initial configuration, the system evolves by single spin flip continuous time dynamics. Let $`W(S_i;S_{i1},S_{i+1})`$ denote the rate at which the flip $`S_iS_i`$ occurs when the two neighboring spins are $`S_{i1}`$ and $`S_{i+1}`$. In our model the rates are specified as follows: $`W(+;++)=W(;)=0`$ (1) $`W(+;+)=W(+;+)=W(;+)=W(;+)={\displaystyle \frac{1}{2}}`$ (2) $`W(+;)=1`$ (3) $`W(;++)=\alpha `$ (4) Here we restrict ourselves to the case $`\alpha =0`$. We note that the case $`\alpha =1`$ corresponds to the usual zero temperature Glauber dynamics. The only difference is that for $`\alpha =0`$, the move $`(+,,+)(+,+,+)`$ is not allowed and thereby the symmetry between ‘$`+`$’ and ‘$``$’ spins is locally dynamically broken. Thus isolated ‘$``$’ spins (surrounded on both sides by a ‘$`+`$’) block the coalescence of ‘$`+`$’ domains and locally favor the ‘$``$’ spins. As a result, we show below, the system eventually decays into the state where all spins are ‘$``$’ but does so in a very slow manner, $`m(t)+11/\mathrm{log}(t)`$. Our main results can be summarized as follows. In contrast to the case $`\alpha =1`$ (where the average size of both ‘$`+`$’ and ‘$``$’ domains grow as $`ł_\pm (t)t^{1/2}`$ at late times and the average magnetization $`m(t)=(l_+l_{})/(l_++l_{})`$ is a constant of motion), for $`\alpha =0`$ we show that at late times, while $`l_+(t)t^{1/2}`$, $`l_{}(t)t^{1/2}\mathrm{log}(bt)`$ where $`b`$ is a constant that depends on the initial condition. Thus due to the dynamically generated local bias, the ‘$``$’ domains grow slightly faster than the ‘$`+`$’ domains and as a result the magnetization decays as, $`m(t)=1+\mathrm{const}/\mathrm{log}(bt)`$ for large $`t`$. Notice that the average domain size grows faster for $`\alpha =0`$ than for $`\alpha =1`$, i.e. paradoxically coarsening is enhanced by putting one of the rates to zero. Our model can alternately be viewed as a toy model of granular compaction if one identifies the ‘$``$’ spins as particles,‘$`+`$’ spins as holes and the $`1`$-d lattice as a section of the bottom layer of a granular system. The final state where all spins are negative (magnetization, $`m=1`$) then corresponds to the fully compact state with particle density $`1`$. The basic mechanism for granular compaction can be summarized as follows. There exist local kinetic ‘defects’ and the system can gain in compaction only by relaxing such local defects. Such relaxation happens via the tapping process. However, these defects become rarer with time and it becomes harder and harder for the system to find such a local defect, relax it and thereby gain in compaction. This is the origin of the slow logarithmic relaxation. In our model, the triplets ‘$`++`$’ play the role of such local defects which decay only via the diffusion of kinks. Thus the diffusion effectively plays the role of tapping. The density of these triplets decays with time and the system finds it progressively harder to relax. Such logarithmic relaxation has been found previously in various models of granular systems. However our model differs from these models in an important way. In previous models, the density, while increasing as $`\rho (\mathrm{})\rho (t)1/\mathrm{log}(t)`$ at intermediate times, finally saturates to its steady state value exponentially fast for finite values of the rates of microscopic processes. In contrast, in our model, the approach to the final state is logarithmic asymptotically even at very late times. In terms of the motion of the domain walls between ‘$`+`$’ and ‘$``$’ phases, our model can also be viewed as a reaction diffusion process. We need to distinguish between the two types of domain walls $`+A`$ and $`+B`$. Note that by definition (originating from a spin configuration) the $`A`$’s and $`B`$’s always occur alternately. The $`A`$’s and $`B`$’s diffuse and when an $`A`$ and a $`B`$ meet, they annihilate only if $`A`$ is to the left of $`B`$, otherwise there is hard core repulsion between them. To start with, we set up our notations. We define $`P_n(t)`$ and $`R_n(t)`$ to be the number of domains of size $`n`$ per unit length of ‘$`+`$’ and ‘$``$’ types respectively. Then, $`N(t)=_nP_n=_nR_n`$ is the number of domains of either ‘$`+`$’ or ‘$``$’ spins per unit length. The density of kinks is therefore $`2N(t)`$. We also define the normalized variables, $`p_n=P_n/N`$ and $`r_n=R_n/N`$. $`p_n`$ (or $`r_n`$) denotes the conditional probability that given a domain of ‘$`+`$’ (or ‘$``$’) has occurred, it is of length $`n`$. Let $`L_+(t)=nP_n`$ and $`L_{}(t)=nR_n`$ denote the densities of ‘$`+`$’ and ‘$``$’ spins. Clearly $`L_+(t)+L_{}(t)=1`$ and the magnetization per unit length is $`m(t)=L_+(t)L_{}(t)`$. The average size of a ‘$`+`$’ and a ‘$``$’ domain are denoted respectively by $`l_+(t)=np_n=L_+(t)/N`$ and $`l_{}(t)=nr_n=L_{}(t)/N`$. Following Glauber’s calculation for the $`\alpha =1`$ case, it is easy to show that for $`\alpha =0`$ case, the domain density $`N(t)=(1S_iS_{i+1})/4`$ of either phase, and the fraction of ‘$`+`$’ spins, $`L_+(t)=(1+S_i)/2`$ evolve according to the exact equations, $$\frac{dN}{dt}=P_1,$$ (5) and $$\frac{dL_+}{dt}=R_1.$$ (6) where $`P_1(t)=(1S_{i1})(1+S_i)(1S_{i+1})/8`$ and $`R_1(t)=(1+S_{i1})(1S_i)(1+S_{i+1})/8`$ are respectively the density of ‘$`+`$’ and ‘$``$’ domains of unit length, i.e., the density of triplets ‘$`+`$’ and ‘$`++`$. It is easy to see physically the origin of these two exact equations. Eq.(5) arises from the fact that the domain density can decrease only via the annihilation of the triplets ‘$`+`$’. Also, on average, the fraction of ‘$`+`$’ spins can decrease only due to the blockage by ‘$`++`$’ triplets giving rise to Eq.(6). Using $`m(t)=S_i=2L_+(t)1`$, we find from Eq.(6) that $`dm/dt=2R_1`$. We note that for the case $`\alpha =1`$, $`dm/dt=0`$ , indicating that the magnetization does not evolve with time. In our case, due to the triplet defects ‘$`++`$’, the average magnetization decays with time. We also note that, unlike the $`\alpha =1`$ case, the evolution equation (6) for the single point correlation function involves two and three point correlations (via $`R_1(t)`$). Writing down the analogous equations for $`R_n(t)`$ gives an infinite hierarchy which makes an exact solution difficult for $`\alpha =0`$. Using $`R_1=r_1N`$ and $`L_+=l_+N`$ in Eq.(6), one can formally solve for $`N(t)`$ in terms of $`r_1`$ and $`l_+`$ as, $$\frac{N(t)}{N(t_0)}=\frac{l_+(t_0)}{l_+(t)}\mathrm{exp}\left(_{t_0}^t\frac{r_1(t^{})}{l_+(t^{})}𝑑t^{}\right).$$ (7) Furthermore if the density of the ‘$`+`$’ phase is $`L_+(t_0)=ϵ`$, then, using the relation $`N(t)=1/[l_{}(t)+l_+(t)]`$ in Eq.(7), we find $$\frac{l_{}(t)}{l_+(t)}=\frac{1}{ϵ}\mathrm{exp}\left(_{t_0}^t\frac{r_1(t^{})}{l_+(t^{})}𝑑t^{}\right)1,$$ (8) clearly showing that the ratio $`l_{}(t)/l_+(t)`$ is growing due to the presence of the triplets ‘$`++`$’. Note that the asymmetry between the growth of ‘$``$’ and ‘$`+`$’ domains is evident due to the triplet defects ‘$`++`$’, present with density $`R_1=r_1N`$. In order to make further analytic progress, we first use the IIA where correlations between neighboring domains are neglected. The IIA was used previously for the pure Glauber-Ising model, i.e., the $`\alpha =1`$ case. It yielded results in agreement, qualitatively as well as quantitatively to a fair degree of accuracy, with the exact results available . Following the derivation of the IIA equations in the $`\alpha =1`$ case, it is straightforward to derive the corresponding equations for the $`\alpha =0`$ case. Under this approximation, the domain densities $`P_n(t)`$ and $`R_n(t)`$ evolve as $$\frac{dP_n}{dt}=P_{n+1}+P_{n1}2P_n+\frac{R_1}{N}(P_nP_{n1})$$ (9) for all $`n1`$ with $`P_0=0`$ (absorbing boundary condition) and $`{\displaystyle \frac{dR_n}{dt}}`$ $`=`$ $`R_{n+1}+R_{n1}2R_n{\displaystyle \frac{P_1}{N}}(R_n+R_{n1})`$ (10) $`+{\displaystyle \frac{P_1}{N^2}}{\displaystyle \underset{i=1}{\overset{n2}{}}}R_iR_{ni1};n2`$ (11) $`{\displaystyle \frac{dR_1}{dt}}`$ $`=`$ $`R_2R_1{\displaystyle \frac{P_1}{N}}R_1,`$ (12) where $`N(t)=P_n=R_n`$. It can be easily checked that these two IIA equations satisfy Eqs.(5) and (6) exactly, and consequently also Eqs.(7) and (8). To calculate $`N(t)`$ using Eq.(7), we need to evaluate two quantities from the IIA equations: (i) $`r_1(t)=R_1/N`$ and (ii) $`l_+(t)=np_n`$ where $`p_n=P_n/N`$. In order to calculate these two quantities, it is useful to write the IIA equations in terms of the normalized variables, $`p_n=P_n/N`$ and $`r_n=R_n/N`$. From Eqs.(9) and (12), we then get, $$\frac{dp_n}{dt}=p_{n+1}+p_{n1}2p_n+r_1(p_np_{n1})+p_1p_n$$ (13) for all $`n1`$ with $`p_0=0`$ (absorbing boundary condition) and $`{\displaystyle \frac{dr_n}{dt}}`$ $`=`$ $`r_{n+1}+r_{n1}2r_np_1r_{n1}`$ (14) $`+p_1{\displaystyle \underset{i=1}{\overset{n2}{}}}r_ir_{ni1};n2`$ (15) $`{\displaystyle \frac{dr_1}{dt}}`$ $`=`$ $`r_2r_1.`$ (16) It is easy to check that the normalization condition $`p_n=r_n=1`$ is satisfied by these two equations. The two IIA equations above are coupled nonlinear equations with infinite number of variables and hence are difficult to solve exactly. Our approach is a combination of a scaling assumption and then rechecking this assumption for self-consistency. Consider first the $`p_n`$ equation, i.e. Eq.(13). We substitute $`p_n(t)=t^{1/2}f(nt^{1/2},t)`$ in Eq.(13) and ask if the resulting equation allows for a steady state scaling solution as $`t\mathrm{}`$, i.e., if the scaling function becomes explicitly independent of $`t`$ as $`t\mathrm{}`$. It is easy to verify that if $`r_1(t)`$ decays faster than $`t^{1/2}`$, such a time-independent scaling solution is possible with $`f(x)=\frac{x}{2}\mathrm{exp}(x^2/4)`$. In this case, $`l_+(t)=np_nt^{1/2}_0^{\mathrm{}}xf(x)𝑑x=\sqrt{\pi t}`$ at late times. Next we consider the $`r_n`$ equation, i.e., Eq.(16). Since $`p_1=d\mathrm{log}N/dt`$, a natural choice would be to write $`r_n(t)=N(t)g(nN(t),t)`$. Substituting this in Eq.(16), we find that in the $`t\mathrm{}`$ limit, the equation allows for a time independent scaling function, $`g(x)=c\mathrm{exp}(cx)`$ ($`c`$ is a constant), provided $`N(t)`$ decays faster than $`t^{1/2}`$. In this case, $`r_1=N(t)g(0)=cN(t)`$. Thus if scaling starts holding beyond some time $`t_0`$, then $`c=r_1(t_0)/N(t_0)`$. Using the results (i) $`r_1(t)=r_1(t_0)N(t)/N(t_0)`$ and (ii) $`l_+(t)=\sqrt{\pi t}`$ in the exact equation Eq.(7), we find $$\frac{N(t)}{N(t_0)}=\sqrt{\frac{t_0}{t}}\frac{\mathrm{log}(b)}{\mathrm{log}(bt/t_0)}$$ (17) where $`\mathrm{log}(b)=\frac{\sqrt{\pi }}{r_1(t_0)\sqrt{t_0}}`$. Substituting this result in the expression for $`r_1(t)`$, we find $$r_1(t)=\frac{\sqrt{\pi }}{\sqrt{t}\mathrm{log}(bt/t_0)}.$$ (18) We now use the late time result Eq.(17) in the exact relation Eq.(5) and find, $$p_1=\frac{1}{2t}+\frac{1}{t\mathrm{log}(bt/t_0)}.$$ (19) From the above expressions, it is clear that both $`r_1(t)`$ and $`N(t)`$ decay faster than $`t^{1/2}`$ and hence our scaling solutions are completely self-consistent. It is easy to see that these IIA results become exact in the zero density limit of the ‘$`+`$’ phase ($`ϵ0`$). In this limit, the average size of a ‘$``$’ domain is $`1/ϵ`$ times larger than the average size of a ‘$`+`$’ domain. As time increases, the ‘$`+`$’ domains will certainly grow in size. But a typical ‘$`+`$’ domain will disappear (via the absorbing boundary condition) much before encountering other ‘$`+`$’ domains, i.e., before feeling the presence of the constraint due to triplets ‘$`++`$’. The probability of such an event is of order $`O(ϵ)`$. Thus, effectively, the dynamics of the system will proceed via eating up of the ‘$`+`$’ domains. Hence, if there is no correlation between domains in the initial condition, the dynamics is not going to generate correlations between them and hence IIA becomes exact. The picture in this limit is similar to the zero temperature dynamics of the $`q`$ state Potts model in the limit $`q1^+`$. For other volume fractions, it is likely that IIA predicts the correct fixed point picture at late times. This is confirmed by Monte Carlo simulations of the model. To improve efficiency, these simulations were made for a version of the model with simultaneous updating. This should not change any of the above conclusions. For convenience, we chose initial conditions such that all domains of minority spins had length 1, while the lengths of majority spin domains were distributed exponentially. At each time $`t`$, all kink positions were written into an array `K`, and only this array is used to generate the array `K’` for the next time step. For each value of $`m(0)`$ we simulated between 20 and 200 lattices of $`2^{26}`$ sites for $`3\times 10^7`$ time steps. Data for $`l_{}(t)/l_+(t)`$, plotted in Fig.1, show the predicted monotonic increase with $`t`$. This increase is logarithmic for $`t>t_0`$, while it is faster for $`t<t_0`$. For a detailed comparison with the above theory we need to know how $`t_0`$ (and thus also $`b`$) depends on $`m(0)`$. We expect it to be exponential for $`m(0)<<1`$, but this is not sufficient for a detailed analysis. Thus we determine for each $`m(0)`$ a $`\tau `$ such that the data for $`l_{}(t)/l_+(t)`$, for $`\sqrt{t}N(t)`$, for $`tp_1(t)`$ and for $`\sqrt{t}r_1(t)`$ collapse when plotted against $`t/\tau `$. The fact that a single $`\tau [m(0)]`$ exists which gives a good data collapse in all four plots is highly nontrivial. We just show such plots for $`N(t),p_1(t)`$ and $`r_1(t)`$ in Figs.2 and 3. We see good agreement with the theoretical predictions. In particular, data collapse is excellent (showing that the only memory left from the initial conditions is the current value of the magnetization). But the detailed predictions for the scaling function show substantial corrections which seem however to disappear for $`t\mathrm{}`$. We thank M. Barma and C. Sire for useful discussions. A related reaction diffusion model has recently been studied numerically by Odor and Menyhard \[cond-mat/ 0002199\]. We thank G. Odor for communicating his results to us.
warning/0004/cond-mat0004010.html	ar5iv	text	# Energy Storage in a Hamiltonian System in Partial Contact with a Heat Bath ## Abstract To understand the mechanism allowing for the long-term storage of excess energy in proteins, we study a Hamiltonian system consisting of several coupled pendula in partial contact with a heat bath. It is found that energy storage is possible when the motion of each pendulum switches between oscillatory (vibrational) and rotational (phase-slip) modes. The storage time increases almost exponentially to the square root of the injected energy. The relevance of our mechanism to protein motors is discussed. Proteins are among the most important biopolymers for living systems. They transform chemical energy to mechanical energy, and vice versa, and contribute to biological functions. However, the question of how proteins work dynamically remains unanswered. Recently, a noteworthy experiment concerning protein motors was performed . In this experiment, the working process of a single molecule was directly investigated. The results suggest that proteins often store energy obtained from a reaction with ATP (adenosine triphosphate) and use it later (e.g., for enzymatic reactions with other proteins). The interval for energy storage was found to sometimes be very long, up to the order of seconds, while typical timescales for normal vibrations are several picoseconds. How can proteins store excess energy for such a long time, somehow overcoming the relaxation process toward thermal equilibrium? In order for a protein to store energy for a sufficiently long time, energy must be absorbed into a certain part of the protein, in accordance with its own dynamics. Furthermore, some characteristic type of dynamics is required to store the excess energy without losing it to the surrounding aqueous solution. As a first approach to understand the working mechanism of proteins, we construct a Hamiltonian system (in partial contact with a heat bath), which stores energy for a given time span, in spite of its eventual relaxation to thermal equilibrium on a much longer timescale. In this Letter, we adopt a system consisting of several coupled pendula , each of which possesses two modes of motion, oscillation and rotation. We clarify the characteristic dynamics necessary for energy storage in connection with the coexistence of these two modes, and also show that partial contact with the heat bath is necessary for this storage. The relationship between the storage time and the injected energy is obtained, and is experimentally verifiable. The relevance of our results to protein motors is also discussed. A protein consists of a folded chain of amino acids that assumes a globular shape . The main chain is accompanied by side chains arranged around it, with each side chain hanging on the main chain similar to a pendulum. Some side chains are gathered in a globular shape, and a certain assembly of them can play an important role in the function of the protein. As an abstract model for the angular motion of side chains in such a functional assembly, we choose a system of coupled pendula. In particular, we study the idealized case of $`N`$ identical pendula equally coupled to each other. Here, the oscillation of the pendula corresponds to the vibration of the side chains. With this simple model as an example, we demonstrate that long-term storage is generally possible in a class of Hamiltonian systems . Our study is restricted to a classical mechanical description, since quantum mechanical effects are believed to be irrelevant to protein dynamics (except for the choice of the potential). Our Hamiltonian is given by $$H=K+V=\underset{i=1}{\overset{N}{}}\frac{p_i^2}{2}+\underset{i,j=1}{\overset{N}{}}V(\theta _i,\theta _j),$$ (1) where $`p_i`$ is the momentum of the $`i`$-th pendulum. The potential $`V`$ is constructed so that each pair of pendula interacts through their phase difference with an attractive force to align the phases : $$V(\theta _i,\theta _j)=\frac{1}{2(2\pi )^2N}\{1\mathrm{cos}(2\pi (\theta _i\theta _j))\}.$$ (2) Hence, the evolution equations for the momentum $`p_i`$ and the phase $`\theta _i`$ are given by $`\dot{p}_j`$ $`=`$ $`{\displaystyle \frac{1}{2\pi N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{sin}(2\pi (\theta _i\theta _j)),`$ (3) $`\dot{\theta }_j`$ $`=`$ $`p_j.`$ (4) In this model, all pairs of pendula interact identically. If the protein were a straight chain without folding, it could be modeled by a one-dimensional chain of pendula. In reality, due to its globular shape, a global interaction occurs between pendula. Although the assumed “mean-field coupling” with uniform strength represents an extreme simplification, the resultant model can capture some general characteristics of protein dynamics. For the (dynamical) function of proteins, it is essential that they exist in aqueous solution, which functions as a heat bath. Since the hydrophobic part of a protein molecule is segregated from the aqueous solution, contact with the heat bath is restricted. Accordingly, we define the equation of motion so that only some pendula in the system contact the heat bath: $$\dot{p}_i=\frac{H}{\theta _i}\gamma p_i+\sqrt{2\gamma T}\xi _i(t)(iN_h).$$ (5) Here, the heat bath is described by the Langevin equation, in which $`T`$ represents its temperature, $`\gamma `$ is a relaxation coefficient, and $`\xi _i(t)`$ is a Gaussian random form satisfying $`\xi _i(t)=0`$ and $`\xi _i(t_1)\xi _j(t_2)=\delta _{ij}\delta (t_2t_1)`$, with $``$ as the temporal average. An important point here is that contact with the heat bath is restricted to the few pendula satisfying $`iN_h`$. This restriction may be interpreted as only allowing for interaction of the hydrophilic part with the aqueous solution. Note that this restricted heat bath (i.e., partial contact with the heat bath) is sufficient for realizing thermal equilibrium. Here, we briefly review the typical behavior of eqs. (3) and (4) in the conservative case , without coupling to the heat bath. When the total energy is small, all pendula vibrate almost synchronously (i.e., $`\|\theta _i\theta _j\|`$ and $`\|p_ip_j\|`$ are small for all $`i`$, $`j`$). The typical timescale of vibration is $`1`$ \- $`10`$, while there is collective motion giving rise to a periodic change in the degree of synchronization, with a timescale of $`10^2`$ \- $`10^3`$. As the total energy increases, rotational (phase-slip) behavior of pendula begins to appear. For a sufficiently large total energy, almost all pendula exhibit rotational motion. Here, each $`p_i`$ changes slowly over time, with relatively large values of $`\|p_ip_j\|`$. In this case, the pendula rotate almost freely, since the correlation of their phase with those of other pendula cannot be maintained, and the interaction term cancels out even with a short time average, except in the rare situation that two momenta take very close values. The timescale of rotation is $`O(1)`$. In the medium energy regime, the rotational motion of a single pendulum appears intermittently from an assembly of vibrating pendula (as displayed in the upper part of Fig. 1). Once a pendulum starts to rotate, it typically continues to rotate over many cycles, which is longer than the typical timescale of each pendulum’s vibration (and rotation). We emphasize here that the effective interaction for a rotating pendulum is much weaker than those for vibrating ones, as mentioned above. In the thermodynamic limit with $`N\mathrm{}`$ , there appears a phase transition at the energy $`E_c/N=0.0190`$ ($`T_c=0.0127`$), where the vibrational mode is dominant in the ‘solid’-like phase and the free rotation is dominant in the ‘gas’-like phase. If $`N`$ is not sufficiently large (as is the case studied here), the transition point is blurred by the finite size effect, and the fluctuations are large around the transition point. Hereafter, we refer to this temperature range ($`0.01\stackrel{<}{}T\stackrel{<}{}0.03`$) as the medium temperature range. Now, we consider the Hamiltonian system in contact with the heat bath described by eq. (5), where the total energy is now time-dependent. It is found that the fluctuations of energy are sometimes correlated with the dynamics of the pendula. For instance, in the medium temperature regime in partial contact with the heat bath, the intermittent appearance of rotating pendula, similar to the Hamiltonian dynamics, becomes accompanied with a large fluctuation of the total energy (see Fig. 1). Here, deviation toward a larger total energy is supported by the concentration of energy in one (or few) pendulum . While local fluctuations are thus dependent on the dynamics, the system reaches thermal equilibrium for a sufficiently long timescale, irrespective of the values of $`\gamma `$, $`N_h`$ ($`1`$), or the temperature. Note that the equilibrium property as a canonical ensemble is identical for any number of $`N_h`$, even for $`N_hN`$. When $`N_h=N`$ (full contact), the timescale to reach thermal equilibrium is determined completely by the value of $`\gamma `$, independently of the dynamical properties of the system. On the other hand, in the case of partial contact, the relaxation process is affected by inherent Hamiltonian dynamics. It progresses beyond the timescale $`\mathrm{\Gamma }^1`$, where $`\mathrm{\Gamma }\gamma N_h/N`$ gives the dissipation rate of the system. Even in this case, the dynamics of the pendula corresponding to $`iN_h`$ are governed by the timescale $`\gamma ^1`$ ($`<\mathrm{\Gamma }^1`$). As the next step, we discuss the energetic behavior of the system far from equilibrium. Consider a special enzymatic event such as ATP attachment or its reaction. With such an event, proteins are moved far from equilibrium. We study how such an event is related to the storage of energy. Due to such a ‘reaction’ event, some portion of the protein is forced far from the previous equilibrium state. This situation is modeled by the addition of an instantaneous kick to a certain pendulum at the reaction event. The kicked pendulum comes to possess a larger amount of momentum and kinetic energy than other pendula . An example of temporal evolution after the kick at $`t=t_0`$ is shown in Fig. 2, where the system with $`N=10`$ and $`N_h=1`$ is adopted; that is only the first pendulum is in contact with the heat bath. Kinetic energy with the amount of $`E_0`$ is added to the kicked pendulum. In Fig. 2, the high-energy state continues up to $`2\times 10^5`$ while the typical relaxation time $`\mathrm{\Gamma }^1`$ of the heat bath is equal to $`10^3`$. There, the kicked pendulum continues to rotate in isolation, maintaining the large energy, while it is affected only slightly by the other pendula over a long time interval. This allows for long-term energy storage. In order to study the lifetime of energy storage, we define it as the interval from the kick until the relaxation of the total energy (see Fig. 2). This is determined by the time $`t_R`$ at which the total kinetic energy $`K`$ decreases to $`NT/2`$, that is the value at thermal equilibrium. The results, however, do not depend on the specific choice of the relaxation time. The distribution of the lifetime is shown in Fig. 3, where evolutions for 200 stochastic processes $`\xi _i(t)`$ are sampled. For comparison, we show the distribution for two cases with the same value of $`\mathrm{\Gamma }`$: $`N_h=1`$ with $`\gamma =10^2`$ and $`N_h=N(=10)`$ with $`\gamma =10^3`$. It is noted that the distribution of the lifetime is quite different for the two cases. For the case of partial contact, $`N_h=1`$, the typical lifetime of the energy storage is very long and reaches $`10^6`$, in contrast with $`10^3`$ for the $`N_h=N`$ case. Moreover, the long-term energy storage is obtained as far as the kicked pendulum is not in contact with the heat bath, even if $`N_hN`$ as long as $`N_hN`$. This result suggests that partial contact with the heat bath is necessary for long-term energy storage. Although the thermal equilibrium properties are similar in the above two cases ($`N_h<N`$ and $`N_h=N`$), a large difference appears when the system is placed far from equilibrium. The pendula $`iN_h`$, interacting directly with the heat bath, cannot rotate freely over the $`1/\gamma `$ timescale (in contrast with the Hamiltonian dynamics) and the dynamics is replaced by Brownian motion due to the heat bath. In the case of partial contact, a long duration of rotation is possible for pendula corresponding to $`i>N_h`$, if the pendula possess sufficiently large energy to remain far from equilibrium. The pendula there become free from the thermal effect, due to the distinctively weaker interaction with the other pendula, and follow nearly pure Hamiltonian dynamics. Although the prototype of this mechanism is observed as a large fluctuation around equilibrium in the medium temperature regime (see Fig. 1), it works well far from equilibrium. The lifetime of energy storage increases with the increase of the kicked energy $`E_0`$, injected to a single pendulum. Figure 4 shows the relationship between the average lifetime and $`E_0`$. The cases in partial contact $`N_h=1`$ with two different $`\gamma `$ values are compared with the Hamiltonian (i.e., microcanonical) case without the heat bath . In the Hamiltonian case, we find the following relation $$t_Rt_0\mathrm{exp}(\alpha \sqrt{E_0}),$$ (6) between the average lifetime and $`E_0`$ with $``$ as the ensemble average, and $`\alpha `$ as a temperature-dependent constant. In a dissipative case with the heat bath, a similar increase of the relaxation time is maintained, although there exists slight suppression of the lifetime at a high energy. In Hamiltonian systems, the exponential law of relaxation time is known as the ‘Boltzmann-Jeans conjecture’ . The energy relaxation between slow and fast modes generally requires a considerable amount of time. The conjecture is confirmed in numerical experiments for a classical gas of diatomic molecules , and some analytic estimates for a generalized version of the conjecture are given within a classical perturbation theory of Hamiltonian systems. In our case, there is only slight energy exchange between the fast rotational and slow vibrational modes. In the above diatomic molecules with translation, rotation, and vibration, when it is sufficiently high, the rotational energy almost freezes, and the transfer of the energy to the translational mode requires a time exponential to the angular momentum, similar to eq. (6). We may expect that this exponential form is generally valid, if the excited mode has a weaker interaction with other modes as the excited energy becomes higher. Thus far, the exponential law has been proven for a class of Hamiltonian systems. Although our model has partial contact with the heat bath, a similar increase of the relaxation time is confirmed. Moreover, the same increase is obtained for a system with the same $`\mathrm{\Gamma }`$ value, for any value of $`N_h`$ (up to $`N1`$), as long as the rotating pendulum does not come into direct contact with the heat bath. On the other hand, if all pendula are in contact with the heat bath, the energy of the rotating pendulum decays much faster than the case with partial contact depicted in Fig. 4. The typical lifetime for this full contact case is of the order $`\mathrm{\Gamma }^1`$ over a wide range of $`E_0`$. The system in partial contact with the heat bath retains a mechanism for slow relaxation, similar to the Hamiltonian system. What type of Hamiltonian is required for long-term energy storage? Consider a system with two possible phases (e.g., gas-like and solid-like phases). Here, some elements possessing higher energy “melt” and come to have a distinctively weaker interaction with elements in the ‘solid’-like phase. Such differentiation is indispensable to the storage of a large amount of energy. Partial contact with a heat bath is essential to maintain the differentiation of states required to store the absorbed energy. Our coupled pendulum model provides a simple example for this behavior. The dynamic mechanism for energy storage presented here is simple enough to be realized in real protein motors. It is only necessary that the protein dynamics exists near a phase transition region with ‘solid’-like (vibrational) modes and ‘gas’-like (rotational) modes and that the folded structure prevents some parts of the protein from experiencing the random effect of a heat bath. These conditions are expected to be satisfied for many kinds of proteins with a sufficiently large size, in addition to the protein motors. It is important to note that one can confirm experimentally if the present mechanism of the energy storage is valid or not by examining the relationship between the storage time and the injected energy, as shown in Fig. 4. The authors are grateful to T. Yanagida, Y. Ishii and members of the Single Molecule Processes Project, JST for many meaningful suggestions. They would also like to thank F. Oosawa, T. Yomo, T. S. Komatsu, S. Sasa and T. Shibata for stimulating discussions. This research was supported by Grants-in-Aids for Scientific Research from the Ministry of Education, Science, Sports and Culture of Japan.
warning/0004/hep-ph0004261.html	ar5iv	text	# Four-quark state in QCD ## 1 Introduction The introduction of color and the development of QCD explain the classification and many spectroscopic characteristics of hadrons, they predict the possibility of the existence of exotic hardonic structures such as glueball and multiquark states too. However, in contrast with the normal hardons, the existence of multiquark states is neither forbidden nor required by color confinement, so the identification of multiquark states is an interesting topic. It is possible too that four-quark configurations lead to hadron-hadron potentials, which play an important role in final state interactions. Therefore, even though there exists no multiquark state, it is helpful to study multiquark configurations to understand final state interactions. The search for exotics has gone on for a long time, especially in light hadrons energy region. As we know, there are several ambiguous resonances among the light resonances, such as the $`\sigma (4001200)`$, $`a_0(980)`$, $`f_0(980)`$, $`f_0(1500)`$, $`f_J(1710)`$, etc. There are many analyses of these resonances. However, since we know little about glueball and multiquark states, especially because they are complicated by mixing among these hadron states, we could not identify them at present. These light hadron states need further investigation. Not until we have complete knowledge about the normal hadron, glueball and multiquark state, can we classify the hadron zoo unambiguously. Experimentally, there are several four-quark candidates, such as $`a_0(980)`$, $`f_0(980)`$, $`f_1(1420)`$, $`f_2(1565)`$, and $`\mathrm{\Psi }(4040)`$ etc. The $`a_0(980)`$ and $`f_0(980)`$ lie below the threshold of $`K\overline{K}`$; the rest lie below the threshold of some other meson pairs($`K\overline{K}^{}`$, $`\omega \omega `$, and $`D^{}\overline{D}^{}`$) too. In practice, their characteristics of decay imply intensely their multiquark bound state role. Theoretically, the four-quark system has previously been studied in the frameworks of the bag model, of nonrelativistic potential models,, and from many other points of view, but the conclusions for the existence of four-quark state are quite model dependent, and few investigations are based on QCD. The $`a_0(980)`$ and $`f_0(980)`$ have been interpreted as the four-quark states in many papers,, but none of them has been definitely confirmed as the four-quark state by experiments so far. It is widely believed that the QCD sum rule is a capable nonperturbative method to extract the properties of hadrons, but there exist few works on multiquark states, especially systematic calculations of four-quark states. So we hope to examine the properties of four-quark states with this method in this paper. Since there are four quarks(including antiquarks) in this multibody system, the analysis is much more complicated than that for the normal hadron. Apart from the complication of dynamical interaction among the quarks, the couplings of color, spin, and flavor among the quarks are not unique either. It is obvious that the four-quark state can consist of diquark-antidiquark or $`q\overline{q}`$ pairs(meson-meson like). We will concentrate our analysis on the simple $`0^{++}`$ meson-meson-like four-quark states only. The spin of the quarks in each $`q\overline{q}`$ pair can couple to both singlets and triplets, while color can couple to both singlets and octets, so there are different combinations for the spin and color between quarks. The QCD sum rule is based on the assumption of the existence of bound state and resonance, but it cannot distinguish the interaction between quarks inside the hadron directly. However, we try to construct a different sum rule with different current, in which couplings of the spin and color are different, to detect some information about the four-quark states. For light quark systems, it may be insufficient to predict the characteristics of the interaction inside the hadron. However, in the case of heavy quark systems, it is possible to study the interactions using the sum rule from the analysis of the wave function of this system. For the presence of interactions, different four-quark states with the same overall quantum numbers will mix with each other, which makes our analyses much more difficult. Accordingly, the mixing between these currents hasnot been taken into account either. To find out any difference between the light four-quark states and heavy four-quark states explicitly, we give the sum rules and conclusions for them, respectively. Surely, we have no intention of identifying the $`a_0(980)`$ and $`f_0(980)`$ only through analyses of the spectra of four-quark states, and no attention has been placed on them. This paper is organized as follows. The features of four-quark state have been analyzed in Sec. 2. In Sec. 3, the sum rules for the $`0^{++}`$ light four-quark states have been constructed. The heavy four-quark states with one heavy quark have been analyzed in heavy quark effective theory(HQET) in Sec. 4. We give the numerical results of the spectra in Sec. 5. The last section is reserved for the conclusion and discussion. ## 2 Review of the features of four-quark state In the bag model, once the confinement has been imposed on by the boundary condition and four quarks inside the bound state have been arranged symmetrically about their center of mass, the spectra and dominant decay couplings were calculated. It was found that $`q^2\overline{q}^2`$ resonances were generally too broad and heavy to show up as bumps in mass spectra except the $`0^{++}`$ state, which can be seen in phase shift analyses. The lighter $`0^{++}`$ four-quark states were predicted to couple strongly to two pseudoscalars while the heavier coupled strongly to two vectors. The observed $`a_0(980)`$ and $`f_0(980)`$ have been identified as the isospin-1 and isospin-0 four-quark states, respectively. In the nonrelativistic potential model,, after the introduction of the color dependent confinement force and hyperfine interactions, the existence of four-quark state is predicted as a dynamical solution of the Schr$`\ddot{o}`$dinger equation. In contrast with the prediction of bag model, it was found that normally the ground state of this four-quark system consisted of two free mesons except for the $`K\overline{K}`$ system(named for the $`K\overline{K}`$ molecule), which was in fact a weakly bound $`0^{++}`$ state. There does not exist a rich discrete spectrum of four-quark states either. The $`a_0(980)`$ and $`f_0(980)`$ have been interpreted as this kind of isospin-1 and isospin-0 $`K\overline{K}`$ molecule, respectively. There are some other points of view too. For the combination of four color quarks into a color singlet hadron, there exist different ways. We can combine two color triplet q’s into a color $`6`$ or $`\overline{3}`$; similarly, we can combine two antitriplet $`\overline{q}`$’s into $`3`$ or $`\overline{6}`$. Therefore, there are two ways to combine the four quarks $`qq\overline{q}q`$ into the final color singlet: $`3\overline{3}`$ and $`6\overline{6}`$. The four-quark states can consist of two $`q\overline{q}`$ pairs too; then we can combine the $`q\overline{q}`$ pair into color singlet $`1`$ and color octet $`8`$, and obtain the final color singlet for the four-quark states from color combination: $`11`$ and $`88`$. It is more complicated because the $`3\overline{3}`$ $`6\overline{6}`$ couplings can mix to give the $`11`$-$`88`$ color configurations. Neither $`\stackrel{}{L}`$ nor $`\stackrel{}{S}`$ is conserved in a relativistic quark model. Nevertheless, if we consider only the S-wave sector, the algebra generated by the states and their currents is an $`SU(2)`$. Therefore, the spin couplings between two quarks can be symmetric triplet and asymmetric singlet. For example, the four-quark scalar states can be considered either as bound states of two ordinary pseudoscalar mesons or as bound states of two ordinary vector mesons. As for the flavor combination, the literature has a detailed description. The special character of the four-quark states may lie in the existence of flavor exotics, such as $`Q=2`$ or $`S=2`$, which may be our best chance to find genuine multiquark states. It was suggested too that, similar to the bound states below the $`K\overline{K}`$ threshold, bound states $`D\overline{K}`$ and $`DK`$ should exist near and possibly below the $`DK`$ threshold. However, except for the shortness of experimence, there existed few quantitative calculations of this kind of heavy quark system either. The development of heavy quark effective theory provides us a capable tool to deal with with these systems. In the analyses of the literatures mentioned above, all the interactions were put into theory by hand, so conclusions about the properties of interactions between quarks inside four-quark states are not conclusive either. Taking into account the complication of couplings of spin and color inside the possible physical four-quark states, we try to construct sum rules with different currents to explore them. No matter how much we can extract from these analyses, they are helpful to us from the point of view of the sum rule itself. For flavor, all the following calculations are kept in the unbroken symmetric $`SU(3)`$. ## 3 Sum rules for the light four-quark state First, let us consider the light $`q\overline{q}`$ pairs four-quark states, in which the color couples to a singlet. For the scalar bound states, they can be regarded as either two bound pseudoscalars or two two bound vectors. So we choose the current as $`j_1(x)=(\overline{q}\mathrm{\Gamma }\mathrm{\Lambda }^mq)(\overline{q}\mathrm{\Gamma }\mathrm{\Lambda }^nq)(x),`$ (1) where for the pseudoscalar quark pairs, $`\mathrm{\Gamma }=\gamma _5`$, while $`\mathrm{\Gamma }=\gamma _\mu `$ for the vector pairs. $`\mathrm{\Lambda }^m`$ is the generator of flavor $`SU(3)`$. To obtain the operator product expansion, two kinds of Feynman diagram should be taken into account: one is unbound while the other is bound. For the pseudoscalar current, the contribution of the bound one is only $`1/12`$ of the unbound one; it is the same case for other currents(the suppressed factor may not be $`1/12`$) too. Therefore, we will omit the contributions of bound diagram in the following calculations. All our calculations are in the x representation, and only those terms contributing to the sum rules after Borel transformation are kept in the following formulas. The operator product expansion for the correlation function with pseudoscalar pairs inside the currents can be expressed as $`\mathrm{\Pi }(q^2)`$ $`=`$ $`i{\displaystyle 𝑑xe^{iqx}T\{j_1(x),j_1^{}(0)\}},`$ $`=`$ $`A(q^2)^4\mathrm{ln}(q^2)B(q^2)^2\mathrm{ln}(q^2)Cq^2\mathrm{ln}(q^2),`$ where A, B and C correspond to the perturbative contribution, two-gluon condensate, and four-quark condensate, respectively, while two-quark condensate vanishes. For $`\mathrm{\Gamma }=\gamma _5`$, $`A={\displaystyle \frac{1}{163840\pi ^6}},B={\displaystyle \frac{3\alpha G^2}{2^{11}\pi }},C={\displaystyle \frac{\overline{q}q^2}{64\pi ^2}},`$ (3) while $`A={\displaystyle \frac{1}{40960\pi ^6}},B={\displaystyle \frac{\overline{q}q^2}{32\pi ^2}}`$ (4) in the case of $`\mathrm{\Gamma }=\gamma _\mu `$, where the two-gluon condensate vanishes also. The imaginary part of the correlation functions can be represented as $`Im\mathrm{\Pi }(s)=\pi f_0^2(m^4)^2\delta (sm^2)+\pi (As^4+Bs^2+Cs)\theta (ss_0),`$ (5) where the first term is from the lowest lying bound state or resonance and the second one is from higher resonances or continuum states. Similarly, the $`q\overline{q}`$ pair can couple to a color octet too, and the currents are chosen as $`j_2(x)=f^{ab_1c_1}f^{ab_2c_2}(\overline{q}^{b_1}\mathrm{\Gamma }\mathrm{\Lambda }^mq^{c_1})(\overline{q}^{b_2}\mathrm{\Gamma }\mathrm{\Lambda }^nq^{c_2})(x),`$ (6) where $`\mathrm{\Gamma }`$ is the same as the definition below formula (1). The coefficients of the correlation functions for $`\mathrm{\Gamma }=\gamma _5`$ and $`\mathrm{\Gamma }=\gamma _\mu `$ are $`A={\displaystyle \frac{3}{163840\pi ^6}},B={\displaystyle \frac{9\alpha G^2}{2^{11}\pi }},C={\displaystyle \frac{3\overline{q}q^2}{64\pi ^2}},`$ (7) and $`A={\displaystyle \frac{3}{40960\pi ^6}},C={\displaystyle \frac{3\overline{q}q^2}{8\pi ^2}},`$ (8) respectively. Then after equating the quark sides with the hadron sides with the dispersion relation, we obtain the mass of the four-quark state, $`m^2(s0,\tau )={\displaystyle \frac{R_{k+1}(s_0,\tau )}{R_k(s_0,\tau )}},`$ (9) where $`s_0`$ is the continuum threshold, $`\tau `$ is the Borel transformation variable and $`R_k(\tau ,s_0)`$ $`=`$ $`{\displaystyle \frac{1}{\tau }}\widehat{L}[(q^2)^k\{\mathrm{\Pi }(Q^2){\displaystyle \underset{k=0}{\overset{n1}{}}}a_k(q^2)^k\}]{\displaystyle \frac{1}{\pi }}{\displaystyle _{s_0}^+\mathrm{}}s^ke^{s\tau }Im\mathrm{\Pi }^{\{pert\}}(s)𝑑s`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^{s_0}}s^ke^{s\tau }Im\mathrm{\Pi }(s)𝑑s.`$ ## 4 Sum rules for the heavy four-quark state In this section, we will discuss the scalar four-quark systems with one heavy quark. With the same consideration as the previous section, the interpolating current corresponding to the color singlet of a quark-antiquark pair is chosen as $`j_3(x)=(\overline{q}\mathrm{\Gamma }h_v)(\overline{q}\mathrm{\Gamma }\mathrm{\Lambda }^mq)`$ (11) where $`q(x)`$ is the light quark field, $`h_v(x)`$ is the heavy quark effective field, and $`v`$ is the velocity of the heavy quark. Then, we construct the correlation function as $`\mathrm{\Pi }(\omega )`$ $`=`$ $`i{\displaystyle d^4xe^{iqx}0\|T\{j_3(x),j_3^{}(0)\}\|0},`$ (12) where $`\omega =2qv.`$ (13) After twice suitable Borel transformations, the imaginary part of it is obtained $`Im\mathrm{\Pi }(\tau )=A\tau ^8+B\overline{q}q\tau ^5+C\alpha _sG^2\tau ^4+D\overline{q}q^2\tau ^2.`$ (14) When $`\mathrm{\Gamma }=\gamma _5`$, $`A={\displaystyle \frac{9}{2^98!\pi ^5}},B={\displaystyle \frac{3}{2^75!\pi ^3}},C={\displaystyle \frac{9}{2^{10}4!\pi ^4}},D={\displaystyle \frac{1}{2^6\pi }}.`$ (15) For $`\mathrm{\Gamma }=\gamma _\mu `$, $`A={\displaystyle \frac{9}{2^78!\pi ^5}},B={\displaystyle \frac{3}{2^65!\pi ^3}},D={\displaystyle \frac{1}{2^5\pi }},`$ (16) and C vanishes. To obtain the results above, we have taken use of the infinite heavy quark mass limit. The current with the color octet quark-antiquark pair is chosen as $`j_4(x)=f^{ab1c1}f^{ab2c2}(\overline{q}^{b1}\mathrm{\Gamma }h_v^{c1})(\overline{q}^{b2}\mathrm{\Gamma }\mathrm{\Lambda }^mq^{c2}).`$ (17) The coefficients of $`A`$, $`B`$, $`C`$ and $`D`$ for $`\mathrm{\Gamma }=\gamma _5,\gamma _\mu `$ are $`A={\displaystyle \frac{27}{2^98!\pi ^5}},B={\displaystyle \frac{9}{2^75!\pi ^3}},C={\displaystyle \frac{27}{2^{10}4!\pi ^4}},D={\displaystyle \frac{3}{2^6\pi }},`$ (18) and $`A={\displaystyle \frac{27}{2^78!\pi ^5}},B={\displaystyle \frac{9}{2^65!\pi ^3}},D={\displaystyle \frac{3}{2^5\pi }},`$ (19) respectively. On the phenomenal side, the correlation function is expressed as $`\mathrm{\Pi }(\omega )={\displaystyle \frac{F_{H^+}^2}{(2\mathrm{\Lambda }\omega )}}+{\displaystyle _{\omega _c}^{\mathrm{}}}𝑑\omega ^{}{\displaystyle \frac{Im\mathrm{\Pi }_s(\omega ^{})}{\omega ^{}\omega }},`$ (20) where the first term on the right side is the dominant pole term resulting from the lowest lying resonance contribution, the second term represents the contribution of the continuum state and higher resonances, and $`\omega _c`$ is the continuum threshold. Making use of the dispersion relations for the correlation functions to equate the quark sides with hadron sides, we obtain $`{\displaystyle \frac{F_{H^+}^2}{(2\mathrm{\Lambda }\omega )}}=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\omega _c}}𝑑\omega ^{}{\displaystyle \frac{Im\mathrm{\Pi }(\omega ^{})}{\omega ^{}\omega }};`$ (21) after Borel transformation, they are turned into $`F_{H^\pm }^2e^{2\mathrm{\Lambda }/T}=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\omega _c}}𝑑\omega ^{}Im\mathrm{\Pi }(\omega ^{})e^{\omega ^{}/T},`$ (22) where $`T`$ is the Borel transformation variable. So the $`\mathrm{\Lambda }`$ can be determined as $`2\mathrm{\Lambda }={\displaystyle \frac{_0^{\omega _c}𝑑\omega ^{}\omega ^{}Im\mathrm{\Pi }(\omega ^{})e^{\omega ^{}/T}}{_0^{\omega _c}𝑑\omega ^{}Im\mathrm{\Pi }(\omega ^{})e^{\omega ^{}/T}}}`$ (23) ## 5 Numerical results of the spectra for the four-quark states In this section, we will give the numerical results of the spectra for the four-quark states. To proceed with the calculation, the mass of the b and c quarks are chosen as $`4.7`$ GeV and $`1.3`$ GeV, respectively. The condensates are chosen as $`0\|\overline{q}q\|0=(0.24GeV)^3,0\|\alpha _sG^2\|0=0.06GeV^4.`$ (24) A few words should be given about the technical details first. In the light quark case, we tried the calculation with $`s_0=1.0,1.5`$ and $`2.5`$ GeV, respectively. The mass square of them is displayed in Figs. 1 and 2 though the platform is not satisfactory. For the heavy four-quark states, the $`2\mathrm{\Lambda }`$ obtained by us is shown in Figs. 3 and 4, where the $`\omega _c`$ are chosen as $`2.0`$, $`3.0`$, and $`4.0`$ GeV, respectively. To find which $`s_0`$ or $`\omega _c`$ is the suitable one for our sum rules, the ordinary criteria of the determination of the threshold and platform are taken into account. All the results are collected in Table I. 1. From the results obtained here, the light four-quark states are found to be light. When the mass of the s quark is taken into account, the conclusion will not change. So it is reasonable to search for four-quark states in the light hadron regions, while the heavy four-quark states with one $`c`$ or $`b`$ quark are found to lie below $`2.0`$ GeV or $`5.5`$ GeV, respectively. The mass prediction of the $`D\overline{K}`$ or $`DK`$ four-quark states is higher than the result here. As for the effect of interactions between quarks on the mass of four-quark states, in principle, it cannot be studied through sum rule methods directly. There is no correspondence between the relative interpolating currents and the bound states. Nevertheless, we can proceed with the sum rules process with different currents as we did previously. Especially, we believe that sum rules may be capable of finding some information about the interactions in heavy quark systems, where there exists a nonrelative limit. It is found that the couplings of spin and color inside the $`q\overline{q}`$ pair have little effect on the mass determination of them. The four different combinations between the spin and color inside the currents result in similar masses in both light and heavy quark systems. Through the lessons from analyses of mesons and glueballs, we have learned that perturbative contributions play a dominant role in the sum rules. However, in the present case, the four-quark condensate plays the dominant role, and all the other contributions(including the perturbative term) are negligible. The results obtained here infer that either the interactions inside the four-quark states are some what special or it is not suitable to deal with these states using QCD sum rules. For systems with one heavy quark, there exists some symmetries and a nonrelativistic model may work well. Therefore, sum rule analyses with different combinations of spin and color in the interpolating currents in heavy quark effective theory may explore some physical information about the four-quark states. It can be seen that the platform for these systems is much better than that for light quark systems. Especially, in contrast with the light quark case, the perturbative contribution dominates the sum rule here, which means that application of the HQET sum rule to heavy four-quark states is suitable. ## 6 Conclusion and discussion The spectra of $`(q\overline{q})(q\overline{q})`$ meson-meson like four-quark states have been calculated using QCD sum rules. Four-quark states with light quarks are found below $`1.0`$ GeV. The mass of heavy four-quark states with a c or b quark has been predicted too. In other models, to extract the properties of hadrons, the characteristics of interactions were put into the theory by hand. Unlikely, both the hadron properties and the interactions inside are predicted on basis of QCD in the framework of sum rules. For four-quark states, whether they are compact states such as the ordinary mesons or weakly bound states such as molecules is of the key concern with us. Moreover, the complicated couplings of spin and color between the quarks inside the bound states need not be speculated about either. A possible choice is to construct suitable currents corresponding to the topic with which we are concerned, and the analysis of the reasonableness for choosing currents is important to us. In this paper, we proceeded with sum rule analyses with different currents, where different combinations of spin and color inside the $`q\overline{q}`$ pair have been tried. In fact, it may be a good way to detect the physical couplings of spin and color inside the bound states in heavy quark systems. Even though there is no correspondence between the currents and physical bound states at all, it is necessary to find which interpolating current is the right one for the special question. Previous results confirm that different combinations of spin and color inside the currents have little effect on the mass determination of four-quark states. Experimentally, we cannot confirm the four-quark state yet; we have observed only some final strong coupling meson-meson channel. For this reason, we consider only the $`(q\overline{q})(q\overline{q})`$ four-quark currents in our sum rules. For a complete sum rule analysis, the $`(qq)(\overline{q}\overline{q})`$ four-quark states should be taken into account too. Moreover, besides the mixing between ordinary hadrons and four-quark states, there exists mixing between the four-quark states with the same overall quantum number also. Accordingly, the mixing between different currents should be considered. Unfortunately, we know little about the mixing, both in QCD and in some models. This problem has remained beyond the scope of this paper. Acknowledgment: The author is grateful to Professor Yuan Ben Dai for useful discussions. Figure caption Fig. 1: Mass square of $`0^{++}`$ four-quark state from current $`j_1(x)`$ versus Borel variable $`\tau `$. Fig. 2: Mass square of $`0^{++}`$ four-quark state from current $`j_2(x)`$ versus Borel variable $`\tau `$. Fig. 3: $`2\mathrm{\Lambda }`$ of $`0^{++}`$ heavy-light four-quark state from current $`j_3(x)`$ versus Borel variable T. Fig. 4: $`2\mathrm{\Lambda }`$ of $`0^{++}`$ heavy-light four-quark state from current $`j_4(x)`$ versus Borel variable T.
warning/0004/astro-ph0004369.html	ar5iv	text	# Global Spiral Modes in NGC 1566: Observations and Theory ## 1 INTRODUCTION There is a general consensus that the regular spiral arms observed in disk galaxies are a manifestation of the global density waves propagating in galactic disks (Binney & Tremaine 1987). In the approach advocated for many years by C. C. Lin and his collaborators (e.g. Bertin et al. 1989), such global waves arise from the gravitational instability of exponentially growing internal modes which are stabilized at some “quasi-steady” level due to nonlinear effects. Another explanation assumes that the global spiral arms arise from swing-amplified chaotic fluctuations originating from N-particle noise (James & Sellwood 1978), or from the gravitational drag of the co-orbiting mass clumps within the disk (Toomre 1990). A long-standing controversy over the physical mechanisms responsible for generation of the global density waves has recently been resolved by Shu et al (1999). Those authors showed that global structures arising from the internal normal modes of a gravitating galactic disk, and the spiral patterns driven by external tidal influences, are both manifestations of the fundamental underlying process of modal selection and growth. The discovery of the protoplanetary disks around T Tauri stars initiated a new wave of interest in global instabilities in self-gravitating disks. A number of authors have developed various techniques to study linear density perturbations in gaseous and the collisionless self-gravitating disks (Adams, Ruden, & Shu 1989; Noh, Vishniac, & Cochran 1991; Laughlin & Różyczka 1996; Vauterin & Dejonghe 1996; Kikuchi, Korchagin, & Miyama 1997). Linear stability analysis methods for self-gravitating disks are well complemented by direct two-dimensional simulations (e.g., Tomley et al 1994; Miyama et al. 1994; Nelson et al. 1998). Zhang (1996), and independently Laughlin, Korchagin, & Adams (1997, 1998) have examined nonlinear self-interaction of a single strong global mode. They found that nonlinear self-interaction is responsible for the nonlinear saturation of the mode, and concluded that the saturation mechanism provides a key for understanding the long-term behavior of spiral density waves. An ultimate aim of any theory of spiral structure is the construction of models of particular galaxies on the basis of the available observations. A number of authors have undertaken such comparisons (e.g., Lin, Yuan, & Shu 1969; Roberts, Roberts, & Shu 1975; Mishurov et al 1979; Elmegreen & Elmegreen 1990). The procedure in these studies has relied on empirical estimates of the positions of the corotation and Lindblad resonances using “optical tracers” in the galactic images. Local dispersion relations were then used to calculate the spiral response under additional assumptions regarding the radial behavior of the stability parameter $`Q`$ derived by Toomre (1964). These two foregoing assumptions, and especially the stipulated choice of Toomre’s $`Q`$-parameter, considerably reduce the predictive power of the spiral density wave theory. Long-slit spectroscopic observations of spiral galaxies can provide a stronger basis for the comparison of density wave theory with observations. Knowledge of the $`z`$-velocity dispersion in a galactic disk provides a quantitative estimate of the surface density distribution under the realistic assumption that disks are self-gravitating in the $`z`$-direction (Bottema 1992, hereafter RB). The velocity dispersion, together with the known rotation curve and the surface density distribution, uniquely determine the axisymmetric background state of the galactic disk, and can be used for linear and nonlinear analyses aimed at the modeling of individual spiral galaxies. The aim of our paper is to model the spiral structure in the nearby grand design spiral galaxy NGC 1566. The galaxy is apparently isolated, and it has been the subject of many photometric and spectroscopic studies. In particular, the rotation curve and the radial profile of the stellar velocity dispersion have been accurately measured (RB). NGC 1566 represents one of the best candidates for the comparison of spiral density wave theory with observations. In contrast to previous comparisons between theory and observations, our study models the spiral arms in the disk of NGC 1566 from “scratch”, without making any additional assumptions with regards to the positions of the principal resonances in its disk or the behavior of Q-parameter. Basing our work on both our own, as well as previously published observations, we undertake the linear global modal analysis of this system, and extend the analysis of the spiral structure into the nonlinear regime using 2D numerical simulations. We then compare the results of our linear, and nonlinear simulations with the observed properties of the spiral arms in NGC 1566. The purpose of this comparison is to determine the extent to which the internal global modes emerging from random perturbations imposed on the gravitating disk of NGC 1566 can be responsible for the observed properties of its spiral arms. In Section 2 we present the optical maps in $`B`$\- and $`I`$-bands for NGC 1566, along with the results of our measurements of the amplitude variations in the spiral arms. In Section 3, we use the available observational data to determine the background equilibrium properties in the disk of NGC 1566. In sections 4 and 5 we describe the basic equations and our computational methods. Section 6 gives the results of the linear global modal analysis. Section 7 describes the results of the nonlinear 2D simulations and discusses our comparison of the theoretical results with observations. ## 2 OBSERVATIONS ### 2.1 Photographic Observations and Data Reduction The data consists of 13 B-band images and 9 I-band images, each with an exposure time of 300 seconds. The images were obtained with the 30in telescope at Mt. Stromlo Observatory (35.32 degrees south latitude, 149.00 degrees east longitude, elevation 770 meters) by the RAPT observers on the night of 12 November 1999. The detector used was a liquid nitrogen cooled, thinned SITe CCD with 1024x1024 square 24 micron (0.36 arc-sec) pixels. A correction for a bad column was made by interpolating the intensities of the adjacent columns. Each image was processed by subtracting a bias image, dividing by a flat field image for the appropriate band, spatially filtering to remove cosmic rays, and subtracting the mean sky. The images were then aligned, and those for each band were combined by taking the median value in each pixel. Observations were made of Landholt CCD standards on the same night. ### 2.2 Inner Spiral Structure in NGC 1566 The morphology of NGC 1566 has been studied observationally in different bands by number of authors (Vaucouleurs 1973; Hackwell & Schweizer 1983; Pence, Taylor, & Atherton 1990). The most prominent morphological feature of NGC 1566 is a near-perfect two-armed spiral pattern observed in its inner parts ($`3.5^{}\times 3.5^{}`$). Pence et al. (1990) report that the faint outer spiral arms of NGC 1566 are $`7.5^{}\times 7.0^{}`$ in size. The broader outer spirals are a continuation of the bright inner spiral pattern, and form a pseudoring at the periphery of the galaxy. Infrared photometry (Hackwell & Schweizer 1983) indicates the presence of a short bar structure. Figure 1 shows the superimposed $`B`$\- and $`I`$-band images of the inner spiral structure of NGC 1566. The blue HII regions are arrayed along the inside edge of the inner spiral arms. Outside of the narrow ridge of blue light, is a broad yellow arm of older stars which decreases its brightness while moving radially outward. The digitized image in $`I`$-band (Figure 2) was used for the measurements of the radial dependence of surface brightness variations in spiral arms. Figure 3 shows the azimuthal profiles of logarithm of the number of counts averaged in six rings at radii $`20^{\prime \prime }`$, $`30^{\prime \prime }`$, $`50^{\prime \prime }`$, $`60^{\prime \prime }`$, $`80^{\prime \prime }`$ and $`100^{\prime \prime }`$ with thickness $`5.4^{\prime \prime }`$ each. This value, proportional to the surface brightness, gives the azimuthal variations of the surface brightness associated with the spiral arms. The amplitude variations measured from the mean surface brightness in a given ring increase with radius. The result concurs with observations of Hackwell & Schweizer (1983), and with the analysis of the optical image of NGC 1566 made by Elmegreen & Elmegreen (1990) who also found an increase of the amplitude variations in spiral arms with radius. The surface brightness variations are about 7 percent at radius $`50^{\prime \prime }`$, and increase up to 57 percent at radius $`100^{\prime \prime }`$. The overall spiral amplitude variations in I-band agree with observational data of Hackwell & Schweizer (1983) who found in the infrared spiral amplitude variations about $`\pm 33`$ percent at $`50^{\prime \prime }`$ radius. ## 3 EQUILIBRIUM MODEL Observational data do not provide a unique set of background equilibrium functions for NGC 1566. We therefore discuss the stability properties of a family of disks allowed within the observational error bars. For numerical purposes, it is convenient to use units in which the gravitational constant $`G`$ is equal to unity, the unit of mass is equal to $`10^{10}M_{}`$, and unit of length is 2 kpc. The velocity and time units are then equal to 149.1 km/sec and $`1.34\times 10^7`$ years. We use these units throughout the paper except specially mentioned cases. The units we use are different from the galactic units in common use, e.g. those defined by Mihalas and Routly (1968) ( 1km/sec, 1 kpc, $`2.32\times 10^5M_{}`$). Our units are more advantageous in the numerical simulations. With these units, the typical dimensionless rotational velocity of the disk and its total mass are of order of unity which is essential in choice of the timestep in the 2D integration. Similar units are used in the 2D numerical simulations of gravitating protoplanetary disks (e.g. Laughlin & Różyczka 1996). Surface Density Distribution. In the $`H`$-band, the surface brightness distribution in NGC 1566 is represented by an exponential law with a radial scalelength $`h_{SB}=15.5^{\prime \prime }`$, corresponding to 1.3 kpc for the adopted distance 17.4 Mpc. The $`z`$-component of the stellar velocity dispersion in NGC 1566 is well fitted by an exponential law with a radial scalelength approximately twice as large as the scalelength of the surface brightness distribution (RB), or with the $`z`$-component of the velocity dispersion proportional to the square root from disk’s surface brightness. Exponential distributions of the velocity dispersion with its value proportional to the square root of the surface brightness were found also in other disk galaxies (Bottema 1993). This fact is consistent with the assumption that the galactic disks can be represented by the locally isothermal, self-gravitating stellar sheets with the surface density $`\sigma (r)`$, and the vertical velocity dispersion $`c_z`$ related as: $$\sigma (r)=c_z^2/\pi Gz_0$$ (1) Here $`G`$ is the gravitational constant, and $`z_0`$ is the effective thickness of the disk. Van der Kruit & Searle (1981) in a study of edge-on disk galaxies found that, to rather high accuracy, the parameter $`z_0`$ does not depend on the radius. We use this assumption in our modeling. RB found that equation (1) provides a good fit to the observational profile of the velocity dispersion with the constant effective thickness of the disk $`z_0=0.7`$ kpc, and with the mass-to-light ratio $`(M/L)_H=0.45\pm 0.15`$. The value of the central velocity dispersion can be considered separately. Bottema finds that the core of NGC 1566 has a constant velocity dispersion $`115\pm 10`$ km/sec. This value does not fit the overall exponential distribution of the velocity dispersion within the NGC 1566 disk, which has a known scalelength $`31^{\prime \prime }`$. Furthermore, the 115 km/sec central velocity dispersion does not match the $`80\pm 15`$ km/sec value measured at radius $`22^{\prime \prime }`$ by van der Kruit & Freeman (1984) and reconfirmed by RB. Bottema thus assumed that the anomalous velocity dispersion in the core is related to the small central bulge. In our model, we choose the central value of the velocity dispersion to be 155 km/sec, which gives a satisfactory fit to the observed radial distribution of the velocity dispersion in the NGC 1566 disk. The exponentially decreasing stellar velocity dispersion is about 3 km/sec at radius $`120^{\prime \prime }`$, corresponding to 10 kpc for the adopted distance. We therefore choose a 10 kpc outer radius for the disk of NGC 1566. Once the $`z`$-component of the stellar velocity dispersion is fixed, equation (1) provides the surface density distribution of the NGC 1566 disk. For the numerical reasons, we assume that the surface density distribution goes to zero at the outer boundary of the disk, and use a surface density distribution of the form: $$\sigma (r)=\sigma _0\mathrm{exp}(r/h_\sigma )[1(r/R_{out})^2]^5\text{,}$$ (2) This distribution incorporates a dimensionless factor $`[1(r/R_{out})^2]^5`$ which causes the surface density at the outer edge of the disk to vanish in a smooth way. The functional form (2) provides a good fit to the observed exponential distribution of the surface brightness in the disks of galaxies. Note, however, that the global stability properties are insensitive to the particular choice of the functional forms of the equilibrium functions. The corresponding $`z`$-component of stellar velocity dispersion is thus given by the equation: $$c_z(r)=c_{z0}\mathrm{exp}(r/2h_\sigma )[1(r/R_{out})^2]^{2.5}$$ (3) The assumed value 155 km/sec of the velocity dispersion at the center of the disk yields a central surface density of $`\sigma _0=2.44\times 10^9M_{}/\mathrm{kpc}^2`$, and a total disk mass equal to $`1.79\times 10^{10}M_{}`$. Rotation Curve. With an inclination of $`28^{}`$, NGC 1566 is not well suited for an accurate determination of the rotation curve. The medium curve on Figure 4 indicated by pluses shows the observed rotational velocity of NGC 1566 in units 149.1 km/sec taken from RB. Errors indicated by the upper and lower dotted curves were also taken from RB, and arise from an uncertain knowledge of the inclination of NGC 1566, which was estimated by RB to be $`\pm 5^{}`$. In modeling NGC 1566’s dynamical properties, we have adopted the rotation curve given by the equation: $$v_0(r)=\frac{V_1r}{(r^2+R_1^2)^{3/4}}+\frac{V_2r}{(r^2+R_2^2)^{3/4}}$$ (4) This rotation curve exhibits the basic features of galactic rotation. It goes to zero for small $`r`$, is “flat” in the outer regions of the disk, and presents a Keplerian decline in the regions beyond the galactic halo. The rotation law (4) differs from the model rotation curve used by Bertin et al (1989). We found, however, that the model curve (4) better reproduces observed rotation in NGC 1566. Figure 4 superimposes this rotation curve on the observational data. The parameters $`V_1,V_2,R_1,R_2`$, determining the “low”, “medium”, and “upper” rotation curves are given in Table 1. Radial Velocity Dispersion. Theoretical arguments based on the results of N-body experiments (Villumsen 1985) suggest that the ratio of vertical and azimuthal velocity dispersions is constant throughout the galactic disks, and maintains a value of roughly 0.6. RB found this value to be adequate in his analysis of the stability properties of NGC 1566. We therefore adopt the ratio $`c_z/c_r=0.6`$ as a “standard”, while also considering higher values $`c_z/c_r=0.8`$, and $`c_z/c_r=1.0`$. Radial velocity dispersion together with an epicyclic frequency and a surface density determines the Toomre $`Q`$ parameter which is an important measure of the stability of a gravitating disk. For a stellar disk, the $`Q`$ parameter is given by $$Q=\frac{c_r(r)\kappa (r)}{3.36\sigma (r)}$$ (5) Here $`\kappa (r)`$ is an epicyclic frequency and $`c_r`$ is the radial velocity dispersion in the disk. Fluid Approximation. In our analysis of the stability properties of NGC 1566, we use the hydrodynamic approximation, modeling the stellar disk as a fluid with a polytropic equation of state with the polytropic index $`\gamma `$: $$P_s=K_s\sigma (r)^\gamma $$ (6) Here $`P_s`$ is the vertically integrated pressure and $`K_s`$ is the polytropic constant. The radial velocity dispersion is then related to the surface density by: $$c_r=(\gamma K_s\sigma (r)^{\gamma 1})^{1/2}$$ (7) In our simulations we choose the value of polytropic constant $`\gamma =2.0`$ which naturally produces the empirical “square root” proportionality between the velocity dispersion and the surface density found in galactic disks. There are some additional arguments justifying the hydrodynamic approximation in the description of stellar disks. In the papers of Marochnik (1966), Hunter (1979), and Sygnet, Pellat, & Tagger (1987) it was shown that in some particular models the behavior of perturbations in the collisionless disks can be described by introducing an isotropic pressure with polytropic constant $`\gamma =2`$, or alternately, a pressure tensor. Using the hydrodynamic approximation, Kikuchi et al. (1997) made a direct comparison of the global stability properties of gravitating disks with solutions of the collisionless Boltzmann equation obtained by Vauterin & Dejonghe (1996). Kikuchi et al. found good qualitative, and in most cases good quantitative agreement between the results obtained in the hydrodynamic approximation, and those found by the direct solution of Boltzmann equation. Halo Potential. The equilibrium rotation of the disk $`v_0(r)`$ is balanced by the external gravitational potential of the rigid halo potential $`\mathrm{\Psi }_H`$, the self-gravity of the disk $`\mathrm{\Psi }`$, and the radial pressure gradient: $$\frac{v_0^2}{r}=\frac{1}{\sigma }\frac{dP_s}{dr}+\frac{d}{dr}\left(\mathrm{\Psi }+\mathrm{\Psi }_H\right)$$ (8) Since the surface density distribution and the equilibrium rotation curve are fixed by equations (2) and (4), the self-gravity of the disk can be found by numerically solving Poisson’s equation, and equation (8) can be used to calculate the gradient of the halo potential $`d\mathrm{\Psi }_H/dr`$. The total mass of a spherically symmetric rigid halo within a radius $`r`$ is thus determined by the expression: $$M_H(r)=r^2\frac{d\mathrm{\Psi }_H(r)}{dr}$$ (9) ## 4 BASIC EQUATIONS Galactic disks are the multi-component systems containing collisionless stars and gas in different phases. The common approximation in studying the spiral structure of galaxies is an assumption that the galactic disks can be considered as a one-component “gas” which can be described by a polytropic equation of state. Previous studies showed that such an approximation adequately describes the mechanism of the growth of global modes, but the effects of the multi-component nature of the galactic disks, and especially the presence of the cold gas component might be important (Lin & Shu 1966; Sellwood & Carlberg 1984; Bertin & Romeo 1988). In this paper we study the dynamics of the disk of NGC 1566 using both one-component, and multi-component models taking the effects of star formation and the self-gravity of gas into account. One-Component Disk. The behavior of our one-component model is described by the standard set of the continuity equation, the momentum equations, and Poisson’s equation in polar coordinates: $$\frac{u}{t}+u\frac{u}{r}+\frac{v}{r}\frac{u}{\varphi }\frac{v^2}{r}=\frac{1}{\sigma }\frac{P_s}{r}\frac{}{r}\left(\mathrm{\Psi }+\mathrm{\Psi }_H\right),$$ (10) $$\frac{v}{t}+u\frac{v}{r}+\frac{v}{r}\frac{v}{\varphi }+\frac{vu}{r}=\frac{1}{\sigma r}\frac{P_s}{\varphi }\frac{1}{r}\frac{}{\varphi }\left(\mathrm{\Psi }+\mathrm{\Psi }_H\right),$$ (11) $$\frac{\sigma }{t}+\frac{1}{r}\frac{}{r}\left(r\sigma u\right)+\frac{1}{r}\frac{}{\varphi }\left(\sigma v\right)=0,$$ (12) and $$\mathrm{\Psi }(r,\varphi )=_{R_{in}}^{R_{out}}_0^{2\pi }\frac{\sigma (r^{},\varphi ^{})r^{}dr^{}d\varphi ^{}}{\sqrt{r^2+r^22rr^{}\mathrm{cos}(\varphi \varphi ^{})}},$$ (13) Here, $`u`$ and $`v`$ are the radial and azimuthal velocities within the disk, and $`\sigma `$ is the surface density. The self-gravity of the gas $`\mathrm{\Psi }`$, the explicit contribution $`\mathrm{\Psi }_H`$ from the rigid halo, and the pressure gradient determine the behavior of perturbations in the disk. All of the dependent variables are functions of the radial coordinate $`r`$, the azimuthal angle $`\varphi `$, and the time $`t`$. Multi-Component Disk. To study how the cold gaseous component and the presence of star formation might affect the growth of spirals in the NGC 1566 disk, we consider a multi-component model introduced by Korchagin & Theis (1999). This model splits the galactic disk into the gas, “active” stars, and “inactive” stellar remnants, considering all three components as fluids, coupled by nonlinear interchange processes and by the common gravitational potential. The nonlinear interchange processes bear the basic features of the “chemo-dynamical” approach, developed by Theis, Burkert, & Hensler (1992). We assume that the disk of NGC 1566 was gaseous in the past. Its initial mass, surface density distribution, and rotation were taken from today’s background distributions for the stellar disk in NGC 1566. We then assume that the growth of spiral perturbations is accompanied by the transformations of the gaseous disk into the “stellar” one given by the right-hand-sides of the continuity equations. In cylindrical coordinates they are: $$\frac{D_g\sigma _g}{Dt}=C_2\sigma _g^2+\eta \frac{\sigma _s}{\tau }$$ (14) $$\frac{D_s\sigma _s}{Dt}=\zeta C_2\sigma _g^2\frac{\sigma _s}{\tau }$$ (15) $$\frac{D_r\sigma _r}{Dt}=(1\zeta )C_2\sigma _g^2+(1\eta )\frac{\sigma _s}{\tau }$$ (16) Here $`D_{g,s,r}/Dt`$ are the corresponding “material” time derivatives written in the cylindrical coordinates $$\frac{D_{g,s,r}}{Dt}=\frac{}{t}+\frac{1}{r}\frac{}{r}ru_{g,s,r}+\frac{1}{r}\frac{}{\varphi }v_{g,s,r},$$ (17) $`\sigma _{g,s,r}`$ are the surface densities, and $`u_{g,s,r}`$ and $`v_{g,s,r}`$ are the radial and azimuthal components of the velocities of gas, stars and remnants. The terms on the right-hand-side of the equations (14)–(16) model the birth of stars according to the Schmidt law ( Schmidt 1959) with the rate proportional to the constant $`C_2`$, return of a fraction $`\eta `$ of stellar mass into the ISM by the ejection, and the formation of stars and remnants with the efficiency $`\zeta `$ and $`1\zeta `$ correspondingly. The parameter $`\tau `$ represents the mean stellar lifetime of massive stars. With the help of the equations (14)–(17), the corresponding momentum equations can be written as: $$\sigma _g\frac{D𝒗_g}{Dt}+P_g+\sigma _g\left(\mathrm{\Psi }+\mathrm{\Psi }_H\right)=C_2\sigma _g^2𝒗_g+\eta \frac{\sigma _s}{\tau }𝒗_s$$ (18) $$\sigma _s\frac{D𝒗_s}{Dt}+P_s+\sigma _s\left(\mathrm{\Psi }+\mathrm{\Psi }_H\right)=\zeta C_2\sigma _g^2𝒗_g\frac{\sigma _s}{\tau }𝒗_s$$ (19) $$\sigma _r\frac{D𝒗_r}{Dt}+P_r+\sigma _r\left(\mathrm{\Psi }+\mathrm{\Psi }_H\right)=(1\zeta )C_2\sigma _g^2𝒗_g+(1\eta )\frac{\sigma _s}{\tau }𝒗_s$$ (20) The gravitational potential $`\mathrm{\Phi }`$ is determined by the total density of all components, and can be written in the form of a Poisson integral as: $$\mathrm{\Psi }(r,\varphi )=_{R_{in}}^{R_{out}}_0^{2\pi }\frac{(\sigma _g(r^{},\varphi ^{})+\sigma _s(r^{},\varphi ^{})+\sigma _r(r^{},\varphi ^{}))r^{}dr^{}d\varphi ^{}}{\sqrt{r^2+r^22rr^{}\mathrm{cos}(\varphi \varphi ^{})}}$$ (21) The equations of state for the “pressures” $`P_{g,s,r}`$ of the three components close the system of equations (14)–(21): $$P_{g,s,r}=K_{g,s,r}\sigma _{g,s,r}^{\gamma _{g,s,r}}$$ (22) The parameters used in the numerical simulations are listed in the Table 1. ## 5 COMPUTATIONAL METHODS ### 5.1 Linear Modal Analysis A linear modal analysis starts by considering non-axisymmetric perturbations of the equations (10)-(13), in the form of a Fourier decomposition: $$f(r)+f_1(r)e^{im\varphi i\omega t}.$$ (23) Here, $`f(r)`$ denotes any unperturbed quantity, $`m`$ is the azimuthal wave number and $`\omega `$ is the complex frequency of the perturbation. Following the commonly used procedure (e.g., Lin & Lau 1979; Adams et al. 1989), a governing integro-differential equation can be obtained after the substituting expression (23) into the equations (10)–(13). To eliminate the singularity at the the origin $`r=0`$, we introduce the coordinate transformations of the perturbed enthalpy $`w_1`$, potential $`\psi _1`$, and the surface density $`\sigma _1`$ as: $$w_1=r^m\stackrel{~}{w}_1,\mathrm{\Psi }_1=r^m\stackrel{~}{\mathrm{\Psi }}_1,\sigma _1=r^m\stackrel{~}{\sigma }_1.$$ After linearizing equations (10)–(13) with respect to the perturbed quantities, and eliminating perturbed velocities, one obtains an equation of the form: $$\frac{d^2}{dr^2}(\stackrel{~}{w}_1+\stackrel{~}{\mathrm{\Psi }}_1)+A\frac{d}{dr}(\stackrel{~}{w}_1+\stackrel{~}{\mathrm{\Psi }}_1)+B(\stackrel{~}{w}_1+\stackrel{~}{\mathrm{\Psi }}_1)\frac{D}{c_r^2}\stackrel{~}{w}_1=0.$$ (24) Here, the coefficients $`A`$, $`B`$, and $`D`$ are given by the expressions: $$A=\frac{2m+1}{r}+\frac{1}{\sigma }\frac{d\sigma }{dr}\frac{1}{D}\frac{dD}{dr},$$ (25) $$B=\frac{m}{r}\left[\left(\frac{1}{\sigma }\frac{d\sigma }{dr}\frac{1}{D}\frac{dD}{dr}\right)\left(1\frac{2\mathrm{\Omega }}{\omega m\mathrm{\Omega }}\right)\frac{2}{\omega m\mathrm{\Omega }}\frac{d\mathrm{\Omega }}{dr}\right],$$ (26) $$D=\kappa ^2(\omega m\mathrm{\Omega })^2.$$ (27) In these equations, $`\mathrm{\Omega }`$ is the angular velocity and $`\kappa `$ is the epicyclic frequency defined as $$\kappa ^2=\frac{1}{r^3}\frac{d}{dr}(r^4\mathrm{\Omega }^2).$$ The perturbed potential $`\stackrel{~}{\mathrm{\Psi }}_1`$ is expressed by the Poisson equation in integral form: $`\stackrel{~}{\mathrm{\Psi }}_1(r)=2\pi G{\displaystyle \frac{\mathrm{\Gamma }(m+1/2)}{\mathrm{\Gamma }(m+1)\mathrm{\Gamma }(1/2)}}`$ $`[{\displaystyle _0^r}\stackrel{~}{\sigma }_1(r^{})\left({\displaystyle \frac{r^{}}{r}}\right)^{2m+1}F({\displaystyle \frac{1}{2}},m+{\displaystyle \frac{1}{2}};m+1;{\displaystyle \frac{r^2}{r^2}})dr^{}`$ $`+{\displaystyle _r^{R_{out}}}\stackrel{~}{\sigma }_1(r^{})F({\displaystyle \frac{1}{2}},m+{\displaystyle \frac{1}{2}};m+1;{\displaystyle \frac{r^2}{r^2}})dr^{}],`$ (28) where $`F`$ is a hypergeometric function. The perturbed enthalpy $`w_1`$ and surface density $`\sigma _1`$ are related by $`w_1=(c_r^2/\sigma )\sigma _1`$. The inner and outer boundary conditions, together with the equations (24) and (28) formulate the eigenvalue problem. The inner boundary condition follows from the regularity of solutions at the center of the disk: $$\frac{d}{dr}(\stackrel{~}{w}_1+\stackrel{~}{\mathrm{\Psi }}_1)+\frac{m}{2m+1}\left[\left(\frac{1}{\sigma }\frac{d\sigma }{dr}\frac{1}{D}\frac{dD}{dr}\right)\left(1\frac{2\mathrm{\Omega }}{\omega m\mathrm{\Omega }}\right)\frac{2}{\omega m\mathrm{\Omega }}\frac{d\mathrm{\Omega }}{dr}\right](\stackrel{~}{w}_1+\stackrel{~}{\mathrm{\Psi }}_1)=0.$$ (29) Likewise, at the outer boundary where $`c_r^2=0`$, we require $$\frac{1}{\sigma }\frac{dP_s}{dr}\left[\frac{d}{dr}(\stackrel{~}{w}_1+\stackrel{~}{\mathrm{\Psi }}_1)+\frac{m}{r}(\stackrel{~}{w}_1+\stackrel{~}{\mathrm{\Psi }}_1)\right]D\stackrel{~}{w}_1=0.$$ (30) The eigenvalue problem formulated by the equations (24)–(28), and the boundary conditions (29) and (30) is solved numerically by means of a matrix method. In this method, the perturbed enthalpy $`\stackrel{~}{w}_1`$ is represented by the $`(N+1)`$-dimensional vector with its value taken at $`(N+1)`$ radial grid points. When the derivatives in equation (24) are expressed as finite differences, and the integral in the equation (28) is expressed by a finite sum, the problem is reduced to an $`(N+1)\times (N+1)`$ matrix equation: $$\underset{k=0}{\overset{N}{}}M_{ik}\stackrel{~}{w}_{1,k}=0,\mathrm{for}i=0,1,\mathrm{},N$$ (31) In practice, the first and last rows in the equation (31) are replaced by the boundary conditions expressed in proper differenced form. Equation (31) has a non-trivial vector solution if the condition $$\mathrm{det}M(\omega )=0$$ (32) is satisfied, which yields an eigenfrequency $`\omega `$ of the global mode. The growth rate and the pattern speed of the mode are then written as $`\mathrm{Im}(\omega )`$ and by $`\mathrm{Re}(\omega )/m`$ respectively. ### 5.2 2D Numerical Simulations To solve the one-component, and multi-component hydrodynamical equations we use two-dimensional numerical codes based on a second order Van Leer advection scheme implemented by Stone & Norman (1992) in a general purpose fluid dynamics code, called ZEUS-2D. This code was designed for modeling astrophysical systems in two spatial dimensions, and it can be used for the simulations in a variety of astrophysical processes. The ZEUS-2D code uses accurate enough hydrodynamical algorithms and allows to include complex physical effects in a self-consistent fashion. This code thus provides a good basis for the implementation of the nonlinear mass transfer processes into the multi–phase hydrodynamics. In our code, the mass and momentum interchange processes between disk components are computed at the first sub-step of the ZEUS-type code. To advance the solutions due to interchange processes given by the right-hand sides of the equations (14)–(16) we use a fifth order Cash-Karp Runge-Kutta routine with the time step limitation imposed by the Courant-Friedrichs-Levy criterion and the values of the parameters $`\tau `$ and $`C_2`$ in the mass and momentum interchange processes. In both codes, the Poisson equation is solved by applying the two-dimensional Fourier convolution theorem in polar coordinates (Binney & Tremaine 1987). By introducing the new variable $`u=\mathrm{ln}r`$, the Poisson integral in equation (21) can be rewritten as: $$\mathrm{\Psi }^{}(u,\varphi )=_{\mathrm{ln}R_{in}}^{\mathrm{ln}R_{out}}_0^{2\pi }\sigma ^{}(u^{},\varphi ^{})K(uu^{},\varphi \varphi ^{})𝑑u^{}𝑑\varphi ^{},$$ (33) where $`\mathrm{\Psi }^{}=\mathrm{\Psi }e^{u/2}`$ is the reduced potential, $`\sigma ^{}=\sigma e^{3u/2}`$ is the reduced surface density, and the Poisson kernel is defined by the expression $`K(u,\varphi )=1/\sqrt{2(\mathrm{cosh}u\mathrm{cos}\varphi )}`$. The Poisson kernel in equation (33) depends on the differences $`uu^{}`$ and $`\varphi \varphi ^{}`$, and therefore the Fourier convolution theorem can be applied to obtain the reduced potential $`\mathrm{\Psi }^{}`$. The codes solve hydrodynamical equations using equally spaced azimuthal, and logarithmically spaced radial zones. For the one-component simulations a grid with $`512\times 512`$ zones was employed. The simulations of the dynamics of multi-component disk were done with a grid of $`256\times 256`$ zones. All simulations were performed on the parallel supercomputer VPP300/16R at the National Astronomical Observatory of Japan. ## 6 LINEAR MODES Figure 5 shows the equilibrium curves of our “standard” model for NGC 1566 disk. These curves were computed for the surface density distribution, velocity dispersion, and the rotation curve in the disk given by the expressions (2), (3) and (4), and for the ratio of the $`z`$-velocity dispersion to the radial one equal to 0.6. The parameters of the rotation curve were chosen to imitate the most probable “medium” rotation curve of NGC 1566. With these parameters, the minimum value of the Toomre $`Q`$-parameter defined by the equation (5) is equal to 2.45. The procedure of searching of the linear global modes, described in the section 5.1, did not reveal any unstable modes. This result is in apparent contradiction with the existence of the spirals in NGC 1566, and we have searched therefore for a possible unstable configuration within observational errors. The unstable equilibrium configuration can be achieved if the disk has a slower rotation and/or a smaller radial velocity dispersion. We calculated the equilibrium properties of the disk with the “low” rotation curve with other equilibrium parameters being unaltered. In this disk, the $`Q`$-parameter has minimum value about 1.65. Nevertheless, the linear modal analysis still did not reveal any unstable global modes. This result concurs with the stability analysis by Vauterin & Dejonghe (1996) of self-gravitating disks, who found that a high enough value of the central velocity dispersion stabilizes the disk. Comparison of their Figure 1 and Figure 12 shows that the disk is stable when the ratio of the central velocity dispersion to the value of the rotational velocity in its “flat” part is $`0.8`$. With the “low” rotation curve, the ratio of the central velocity dispersion to the maximum value of the rotational velocity is about 1.0, which stabilizes the system. Another way to achieve an unstable configuration is to increase the ratio $`c_z/c_r`$, and hence decrease the radial velocity dispersion. N-body experiments by Lacey (1984) predict the value of the ratio of the velocity dispersions $`c_z/c_r`$ to be about 0.8 which is larger than the value $`c_z/c_r0.6`$ found by Villumsen (1985). We studied the stability properties of the disk which has a “medium” rotation curve and a ratio of the velocity dispersions $`c_z/c_r=0.8`$. The minimum value of $`Q`$-parameter in such a model is lower than that in the “standard” disk ($`1.8`$), but the linear stability analysis shows that such a disk is stable. The situation changes, however, if we choose the “low” rotation curve. The matrix method described in the previous section yields an $`m=2`$ global mode with eigenvalues $`\mathrm{Re}(\omega _2)=0.893`$, $`\mathrm{Im}(\omega _2)=0.347`$ as the most unstable global mode. Its nearest competitor, a three-armed spiral, has eigenvalues equal to $`\mathrm{Re}(\omega _3)=1.426`$, and $`\mathrm{Im}(\omega _3)=0.156`$. Similar behavior occurs in a disk with a velocity dispersion ratio increased to $`c_z/c_r=1.0`$. We found that this disk is unstable with the most probable “medium” rotation curve. Figure 6 presents the equilibrium properties of such a disk, and Figures 7 and 8 show contour plots for dominant $`m=2`$, and its nearest competitor, $`m=3`$ global modes. The contour levels in Figures 7 and 8 are logarithmically spaced between the maximum value of the perturbed density, and one-hundredth of the maximum perturbed density. The eigenvalues of these two modes found by the matrix linear method are $`\mathrm{Re}(\omega _2)=1.567`$, $`\mathrm{Im}(\omega _2)=0.280`$, and $`\mathrm{Re}(\omega _3)=2.859`$, $`\mathrm{Im}(\omega _3)=0.217`$, respectively. In summary, the results of our linear modal analysis allow us to conclude that the disk of NGC 1566 is unstable if the ratio of the vertical and the azimuthal velocity dispersions is close to unity. With the ratio $`c_z/c_r`$ equal to unity, the disk is unstable with the observationally most probably “medium” rotation curve. If the disk of NGC 1566 has a lower value of this ratio, the actual rotation curve of the disk of NGC 1566 must be close to the lower limit determined by the observational errors. The $`m=2`$ spiral, and the $`m=3`$ spiral are the most unstable in the disk, and the dominating $`m=2`$ global mode resembles open spiral arms in the inner regions of NGC 1566. Linear modal analysis is unable, however, to determine the relative amplitudes of the two competing unstable modes at their nonlinear saturated stage. Furthermore, it can not predict the radial dependence of the amplitude actually observed in spirals. These questions are addressed by the nonlinear two-dimensional simulations which we discuss in the following sections. ## 7 NONLINEAR SIMULATIONS ### 7.1 Nonlinear Dynamics of a “Standard” Disk The results of the linear modal analysis are well complemented by direct two-dimensional simulations of disk dynamics. The time dependence of the global modes in the disk can be expressed in terms of the global Fourier perturbation amplitudes, which are defined as: $$A_m\frac{1}{M_d}\left\|_0^{2\pi }_{R_{in}}^{R_{out}}\sigma (r,\varphi )r𝑑re^{im\varphi }𝑑\varphi \right\|.$$ (34) Here $`M_d`$ is the mass of the disk. Figure 9 plots the time dependence of the global Fourier amplitudes for modes $`m=16`$ resulting from a random initial density perturbation seeded into every grid cell of the “standard” disk with a ratio of dispersions $`c_z/c_r=0.6`$. Within each cell, the perturbation can have a value of up to one part in a thousand of the equilibrium density. Figure 9 confirms the results of the linear modal analysis. The “standard” disk with a ratio of dispersions $`c_z/c_r=0.6`$ is stable for the “medium” rotation curve. Similar stable behavior of the global modes is observed in the disk with “low” rotation. Note that random perturbations at the level $`10^6`$ exist in the disk during the entire computational timeframe. This low level of departure from the equilibrium state does not lead to the growth of the swing-amplified spirals. Our previous consideration of the stability properties of galaxy NGC 1566 was based on a one-component approximation. An additional cold gas component may play an important role in the destabilization of self-gravitating disks. In linear approximation, this effect has been studied by number of authors, beginning with the work of Lin & Shu (1966) and Lynden-Bell (1967) (see also Jog & Solomon 1984; Sellwood & Carlberg 1984; Bertin & Romeo 1988). Recently, Korchagin & Theis (1999) extended the realization of the destabilizing role of a cold gas component onto multi-phase star-forming disks. They demonstrated that spiral structure will grow faster on a non-stationary star-forming environment in comparison to a one-component system with the same total surface density distribution and rotation curve. It might be possible therefore, that the presence of gas will change the stability properties of a “standard” disk. To study this possibility, we simulated the dynamics of the multi-component model described in the Section 4. We started our simulations by assuming that the multicomponent disk is initially in a gaseous phase with a small admixture of stars, and that the gaseous disk has a surface density distribution which corresponds to the observed density profile in NGC 1566. The rotation curve of a gaseous disk was assumed to be the most probable “medium” curve for the galaxy NGC 1566. Figure 10 shows the equilibrium properties of the gaseous component computed at the initial moment of time. The minimum value of the $`Q`$-parameter for the gaseous phase $`Q=1.22`$, and hence, the multi-component disk is unstable to spiral perturbations. This disk was also seeded with random initial density perturbations of amplitude up to one-thousandth of the equilibrium density. The subsequent behavior of this disk is determined by the mass and momentum transformations given by the right-hand-sides of equations (14)–(16), and by the development of the random initial perturbations. The two-dimensional simulations demonstrate that the multi-component disk is indeed unstable, which is illustrated by the global amplitude diagram (Figure 11) plotted for the density perturbations growing in the stellar component. The growth of the Fourier harmonics is accompanied by the transformation of the gas phase into the stellar one as shown in Figure 12. The main unstable mode has an azimuthal wavenumber $`m=2`$, but the growth of perturbations is stabilized at a rather low level $`\mathrm{log}_{10}A_m2`$. More importantly, the growing perturbations do not resemble the spiral arms observed in the disk of NGC 1566. Figure 13 displays contour plots of the density perturbations in the gaseous phase taken at a time $`t=20`$. In the central region we observe a short open spiral outlined by the shock fronts, which is embedded in a system of weak ring-like structures. At later stages of evolution, when most of the gas is transformed into the stellar component, the density perturbation in the stellar component resembles a short bar-like perturbation rather then a global spiral structure (Figure 14). ### 7.2 “Cold” Disks The results of the linear modal analysis described in Section 6 show that NGC 1566 is unstable with respect to the growth of a two armed spiral if the ratio of the velocity dispersions, $`c_z/c_r`$, is close to unity. Figure 15 plots the behavior of the global amplitudes for the Fourier harmonics $`m=16`$ developing from random perturbations in the disk with the ratio of the velocity dispersions $`c_z/c_r`$ equal to unity. In agreement with the linear modal analysis, the nonlinear simulations show that the $`m=2`$ global mode is the primary instability of the disk. The dominant two-armed mode experiences exponential growth until time $`t=40`$, when the exponential growth phase merges into a lingering saturation phase, at which point the $`m=2`$ mode saturates at an amplitude level of $`\mathrm{log}_{10}A_20.5`$. The modal growth rates for the $`m=2`$ perturbation (as well as its nearest competitor, the $`m=3`$ mode) observed in the nonlinear simulations are equal to $`\mathrm{Im}(\omega _2)=0.27`$ and $`\mathrm{Im}(\omega _3)=0.21`$, which is in a good agreement with the results of the linear modal analysis $`\mathrm{Im}(\omega _2)=0.280`$, $`\mathrm{Im}(\omega _3)=0.217`$. Figure 16 shows a contour plot of the surface density perturbation taken at time $`t=30`$ while the perturbations were experiencing the linear phase of exponential growth. Similarly to Figures 7 and 8, contours are logarithmically spaced between the maximum value of perturbed density and one-hundredth of the maximum density perturbations. Figure 16 further demonstrates that the two-armed open spiral, which looks similar to the one obtained in the linear global analysis (see Figure 7), emerges from the random perturbations seeded in the unstable disk with $`c_z/c_r=1`$. During the saturation phase, when the exponential growth of the spiral pattern is stopped, the spiral pattern deforms from its originally sinusoidal shape. Figure 17 shows contour plots for the density perturbations in the disk taken at time $`t=45`$. These contour plots depict the formation of the shock front at the concave edge of the spiral arms. However, the perturbation maintains itself as a two-armed spiral during the entire simulation duration. In Section 2, we discussed the observed surface brightness variations in the spiral arms of galaxy NGC 1566. To compare the amplitude variations of the perturbations seen in the numerical simulations with the observed surface brightness variations of the NGC 1566 spirals, we plotted in Figure 18 the azimuthal variations of the function defined as $`2.5\mathrm{log}_{10}\sigma (t,R_N,\varphi )`$, which was taken at time $`t=45`$ at the fixed radii $`R_N`$ equal to 0.25, 0.5, 1.0, 1.5 and 2.5. Figure 18 should be compared directly with Figure 3. The overall behavior seen in both Figures is quite similar. Near the center of the disk, the spiral perturbation has small amplitude variations. The “surface brightness” variations shown in Figure 18 increase with radius out to approximately $`r=2`$, but at larger radii $`r2.0`$, the “surface brightness” variations tend to decrease. The surface brightness variations in the spiral arms of NGC 1566 do not show a tendency to decrease within a 8.5 kpc radius. This fact reflects a basic discrepancy between the theoretical predictions and the observations. The spiral arms found in linear, and nonlinear simulations are considerably shorter in comparison to the observed spirals in NGC 1566. The growth of unstable global modes is accompanied by momentum and surface density redistribution of the background axisymmetric disk structure. Nonlinear self-interaction of global modes (Laughlin et al. 1998), and spiral shocks funnel matter towards the center of the disk, resulting in a steepening of the surface density profile. Figure 19 illustrates how the long-term evolution of spiral perturbations affects the azimuthally averaged surface density distribution within the disk, making it steeper than the observed density distribution. The detailed kinematic study of NGC 1566 undertaken by Pence et al. (1990) found considerable discrepancy between the observed spiral velocity field and the theoretical predictions. In constructing a theoretical velocity field, Pence et al. (1990) used results of calculations by Roberts & Hausman (1984), who studied the evolution of colliding gas clouds moving in a fixed sinusoidal spiral potential. Pence et al. (1990) found that the observed velocity residuals across the spiral arm reach a maximum at the edge of the arm, contrary to the predictions following from the calculations of Roberts & Hausman (1994). The observed gradient of the velocity field was found to be approximately five times larger than predicted. Figure 20 plots the azimuthal profile of the residual velocity obtained in our numerical simulations superimposed on the surface density profile sampled at radius 1.0 at time $`t=45`$. The results of the numerical simulations show qualitative agreement with observations. The velocity field has a strong velocity gradient associated with the shock front. The maximum of the residual velocity is located near the minimum of the density distribution, but the velocity shift found in the numerical simulations is considerably larger than the values reported by Pence et al. (1990). ## 8 CONCLUSIONS This paper has focused on the observational study of the spiral arms in galaxy NGC 1566, and on the theoretical study of the stability properties of this galaxy. We have applied a global modal approach and have used two-dimensional one-component and multi-component simulations to study the dynamics of the self-gravitating disk in the galaxy NGC 1566 and we have extended our analysis to the nonlinear stage using 2D numerical simulations. In our theoretical analysis of the spiral structure in NGC 1566 we have not made any additional assumptions used in previous comparisons, and we have followed the development of the spiral arms emerging from random perturbations all the way through to the nonlinear saturation stage. The general conclusions which can be made from our work are as follows: As C.C. Lin and his collaborators envisioned many years ago, the most unstable linear global mode emerging from stochastic noise in a galactic disk determines the appearance of the spiral structure on linear, and on nonlinear stages of the evolution of the spiral pattern. We have found a good agreement between the linear global modal analysis of the stability of the disk of galaxy NGC 1566, and direct 2D numerical simulations. The theoretical spiral pattern obtained from the linear global modal analysis, and the pattern emerging from the random perturbations in the nonlinear simulations are in agreement, at least qualitatively, with the observed two-armed spiral structure in the disk of NGC 1566. We did not aim to discriminate between the two principal physical mechanisms explaining the physics of spiral instabilities in gravitating disks. We concur however with the conclusion of Shu et al (1999) that “…there is a coexistence of the two principal mechanisms that produce the beautiful structures that astronomers observe in the universe of spiral galaxies”. The specific results of our analysis are specified below. 1. The CCD images of NGC 1566 in $`B`$ and $`I`$-bands obtained with the Australian National University 30in telescope were used for measurements of the radial dependence of the amplitude variations in the spiral arms of NGC 1566. The azimuthal variations of the surface brightness in the $`I`$-band increase with radius up to $`57`$ % at $`100^{\prime \prime }`$. Our results concur with the previous measurements of the amplitude variations within the spiral arms of NGC 1566. 2. The linear stability analysis of the disk of NGC 1566 made under the standard assumption used in the galactic dynamics, namely that the ratio of the vertical and the radial velocity dispersions is equal to 0.6, shows that disk is stable towards spiral perturbations within observational error bars. We confirm this conclusion with help of 2D simulations of the disk evolution seeded with random perturbations. 3. By increasing the $`c_z/c_r`$ ratio up to $`0.81.0`$, the disk, when seeded with random perturbations, becomes unstable with respect to $`m=2`$, and $`m=3`$ spiral modes. The growth rates and the shapes of the global modes found in the linear analysis, and those seen in the nonlinear simulations are in good agreement. The two-armed spiral mode prevails over its competitors, and thus determines the behavior of perturbations during the linear, and during the nonlinear phases (see Figures 16 and 17). At the nonlinear stage, the two-armed spiral saturates at an amplitude $`\mathrm{log}_{10}A_20.5`$. The surface density, and the velocity residual variations in the arms are in a qualitative agreement with observations. 4. The nonlinear phase of instability is characterized by the transport of angular momentum towards the disk center making the surface density distribution more steep than the observed surface brightness profile in the disk of NGC 1566. The spiral arms found in the linear modal analysis, and seen in the nonlinear simulations, are considerably shorter than the observed spiral arms in the disk of NGC 1566. We argue therefore, that the surface density distribution in the disk of the galaxy NGC 1566 was different in the past when spiral structure in NGC 1566 was growing linearly. We thank the RAPT Group of amateur astronomers (E. Pozza, A. Brakel, B. Crooke, S. McKeown, G. Wyper, K. Ward, D. Baines, P. Purcell, T. Leach, J. Howard, D. McDowell, A. Salmon, A. Gurtierrz) for providing the images from the 30in telescope at the ANU’s Mt. Stromlo Observatory. VK acknowledges Prof. S. Miyama for hospitality, and National Astronomical Observatory of Japan for providing COE fellowship. BAP acknowledges the hospitality of the National Astronomical Observatory of Japan. The computations were performed on the Fujitsu VPP300/16R at the Astronomical Data Analysis Center of the National Astronomical Observatory, Japan.
warning/0004/math-ph0004023.html	ar5iv	text	# Untitled Document e to the A, in a New Way, Some More to Say Paul Federbush Department of Mathematics University of Michigan Ann Arbor, MI 48109-1109 (pfed@math.lsa.umich.edu) Abstract Expressions are given for the exponential of a hermitian matrix, $`A`$. Replacing $`A`$ by $`iA`$ these are explicit formulas for the Fourier transform of $`e^{iA}`$. They extend to any size $`A`$ the previous results for the $`2\times 2,3\times 3`$, and $`4\times 4`$ cases. The expressions are elegant and should prove useful. The support of the Fourier transform of $`e^{iA}`$ was established by E. Nelson in . (That is the Fourier transform of each entry of the matrix $`e^{iA}`$ in terms of the entries of the hermitian matrix $`A`$.) But I believe this result of E. Nelson is very little known in the mathematical community at large. In further work , , the transform was exhibited in the $`2\times 2`$ case, and presented in some unwieldy forms in higher dimensions. In a previous paper,, explicit formulas were obtained for $`2\times 2,3\times 3`$ and $`4\times 4`$ matrices. We here treat the general case. Let $`A`$ be an $`r\times r`$ hermitian matrix. We write $$\mathrm{Det}(1A)=\underset{j=0}{\overset{r}{}}P_j(A)$$ (1) where $`P_j(A)`$ is homogeneous of degree $`j`$ in the entries of $`A`$. The formulas we obtain for the exponential of $`A`$ are as follows: $$(e^A)_{\alpha \beta }=\frac{1}{\mathrm{\Gamma }(r)}\underset{j=0}{\overset{r}{}}P_j(A)\frac{d^{rj}}{ds^{rj}}\left(𝑑\mathrm{\Omega }e^{sTr(AW)}W_{\alpha \beta }s^r\right)\|_{s=1}$$ (2) or $$(e^A)_{\alpha \beta }=\frac{1}{\mathrm{\Gamma }(r)}\underset{j=0}{\overset{r}{}}P_j(A)\frac{d^{rj}}{ds^{rj}}\left(𝑑\mathrm{\Omega }e^{s<A\stackrel{}{n},\stackrel{}{\overline{n}}>}n_\alpha \overline{n}_\beta s^r\right)\|_{s=1}$$ (3) Here $`\stackrel{}{n}`$ is a unit vector in $`\text{ }\mathrm{C}^r`$, and $`W_{ij}=n_i\overline{n}_j`$, a rank one hermitian matrix. $`𝑑\mathrm{\Omega }`$ denotes a normalized integral over all such $`\stackrel{}{n}`$, an integral over the unit sphere in $`\text{ }\mathrm{C}^r`$ with unitary-invariant measure. That the support of the Fourier transform lies on the complex projective space of such $`W`$ is the content of Nelson’s theorem. In fact the formulas in (2) and (3) do not coincide with formulas in when $`r=2,3`$ or 4, but formulas of such type are not unique. We do not know the full scope of such non-uniqueness. We first sketch a derivation/proof of formulas (2) and (3), especially emphasizing the ideas. We note the relation between gaussian integrals in $`d=2r`$ real dimensions, and integrals over the corresponding unit sphere $`S^{d1}`$. $`{\displaystyle \frac{1}{𝒩}}{\displaystyle 𝑑x_ie^{\mathrm{\Sigma }x_i^2}\stackrel{2N}{}x_{\alpha (i)}}`$ $`=`$ $`{\displaystyle \frac{1}{𝒩}}{\displaystyle r^{d1}r^{2N}e^{r^2}𝑑r𝑑\mathrm{\Omega }^{}\stackrel{2N}{}n_{\alpha (i)}}`$ (4) $`=`$ $`{\displaystyle _{S^{d1}}}𝑑\mathrm{\Omega }{\displaystyle \stackrel{2N}{}}n_{\alpha (i)}{\displaystyle \frac{𝑑re^{r^2}r^{2N+d1}}{𝑑re^{r^2}r^{d1}}}`$ (5) $`=`$ $`{\displaystyle _{S^{d1}}}𝑑\mathrm{\Omega }{\displaystyle \stackrel{2N}{}}n_{\alpha (i)}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{2N+d}{2}\right)}{\mathrm{\Gamma }\left(\frac{d}{2}\right)}}.`$ (6) Here $`n_i`$ is the unit vector parallel to $`x_i`$, and $`𝑑\mathrm{\Omega }^{}`$ is integral over the sphere in its usual measure and $`𝑑\mathrm{\Omega }`$ the normalized spherical measure. From (6) we see the integral over a unit sphere of a homogeneous polynomial of degree $`2N`$ is “approximately” $`1/N!`$ the gaussian integral of the same polynomial. We note that multiplying the term in $`A^k`$ by $`\frac{1}{k!}`$ induces a transform (formally) as follows $$1+A+A^2+\mathrm{}=\frac{1}{1A}e^A.$$ (7) We consider the gaussian integral formula (for $`\|A\|<1`$): $$\frac{\mathrm{Det}(1A)}{𝒩}𝑑x_ie^{\mathrm{\Sigma }\|x_i^2\|+<A\stackrel{}{x},\stackrel{}{\overline{x}}>}x_\alpha \overline{x}_\beta =\left(\frac{1}{1A}\right)_{\alpha \beta }.$$ (8) In the expansion of the integrand on the left side of (8) each power of $`A`$ has associated to it two powers of $`x`$. Thus converting from a gaussian integral to an integral over a unit sphere approximately multiplies each power of $`A^N`$ by $`\frac{1}{N!}`$, which would convert $`\frac{1}{1A}`$ to $`e^A`$. The wrong formula we get putting these ideas together would yield: $$\mathrm{`}\mathrm{`}\mathrm{Det}(1A)𝑑\mathrm{\Omega }e^{<A\stackrel{}{n},\stackrel{}{\overline{n}}>}n_\alpha \overline{n}_\beta =(e^A)_{\alpha \beta }\mathrm{"}$$ (9) We turn to the easy task of converting the above careless argument leading to the wrong formula (9), to the detailed correct computation that turns (9) into (2). (Hitherto we have trodden a path redolent with the creative epiphanies of mathematical research, we now segue to the ineluctable concomitant consecration to inferential syntax.) We start from the right side of equation (3), and assume for the moment $`\|A\|<1`$. $$\frac{1}{\mathrm{\Gamma }(r)}\underset{j=0}{\overset{r}{}}P_j(A)\frac{d^{rj}}{ds^{rj}}\left(𝑑\mathrm{\Omega }e^{s<A\stackrel{}{n},\stackrel{}{\overline{n}}>}n_\alpha \overline{n}_\beta s^r\right)\|_{s=1}$$ (10) We expand the exponent and perform the operations on $`s`$, getting $$\underset{j=0}{\overset{r}{}}P_j(A)d\mathrm{\Omega }\underset{k=0}{\overset{\mathrm{}}{}}(<A\stackrel{}{n},\stackrel{}{\overline{n}}>)^kn_\alpha \overline{n}_\beta \frac{1}{\mathrm{\Gamma }(r)}\frac{1}{k!}\frac{(r+k)!}{(k+j)!}$$ (11) Now we use the equality of equations (4), (5), (6) to convert (11) to $$\underset{j=0}{\overset{r}{}}P_j(A)\frac{1}{𝒩}dx_ie^{\mathrm{\Sigma }x_i^2}\underset{k=0}{\overset{\mathrm{}}{}}(<A\stackrel{}{x},\stackrel{}{\overline{x}}>)^kx_\alpha \overline{x}_\beta \frac{1}{\mathrm{\Gamma }(r)}\frac{1}{k!}\frac{(r+k)!}{(k+j)!}\frac{\mathrm{\Gamma }(\frac{d}{2})}{\mathrm{\Gamma }(\frac{2(k+1)+d}{2})}$$ (12) We rewrite the equality of equation (8) in expanded form $$\underset{j=0}{\overset{r}{}}P_j(A)\frac{1}{𝒩}dx_ie^{\mathrm{\Sigma }x_i^2}\underset{k=0}{\overset{\mathrm{}}{}}(<A\stackrel{}{x},\stackrel{}{\overline{x}}>)^kx_\alpha \overline{x}_\beta \frac{1}{k!}=1+A+A^2+\mathrm{}$$ (13) We mark the fact that equation (13) is a separate equality for each homogeneous degree in powers of $`A`$. On each side of the equation we multiply terms homogeneous of degree $`\mathrm{}`$ by $`\frac{1}{\mathrm{}!}`$, arriving at $$\underset{j=0}{\overset{r}{}}P_j(A)\frac{1}{𝒩}dx_ie^{\mathrm{\Sigma }x_i^2}\underset{k=0}{\overset{\mathrm{}}{}}(<A\stackrel{}{x},\stackrel{}{\overline{x}}>)^kx_\alpha \overline{x}_\beta \frac{1}{k!}\frac{1}{(j+k)!}=1+A+\frac{A^2}{2!}+\mathrm{}=e^A.$$ (14) The equality of the left side of (14) with (12) follows from $$\frac{1}{\mathrm{\Gamma }(r)}\frac{1}{k!}\frac{(r+k)!}{(k+j)!}\frac{\mathrm{\Gamma }(\frac{d}{2})}{\mathrm{\Gamma }(\frac{2(k+1)+d}{2})}=\frac{1}{k!}\frac{1}{(j+k)!}$$ (15) using $`r=2d`$. We have thus established our equalities of equations (2) and (3) for $`\|A\|<1`$. But each side of these equations is analytic in the elements of $`A`$, so the equalities hold for all hermitian $`A`$. Acknowledgment: I would like to thank Alexander Barvinok for an all important discussion on evaluating integrals over the unit sphere. References * \] E. Nelson, Operants: A functional calculus for non-commuting operators, Functional Analysis and Related Fields, Proceedings of a conference in honor of Professor Marshal Stone (Univ. of Chicago, May 1968) (F.E. Browder, ed.), Springer-Verlag, Berlin, Heidelberg, and New York, 1970, pp. 172-187. MR 54:978. * \] B. Jefferies, “The Weyl Calculus for Hermitian Matrices”, Proc. A.M.S. 124 (96) p. 121-128. * \] M.E. Taylor “Functions of Several Self-Adjoint Operators”, Proc. A.M.S. 19 (1968), 91-98. MR 36:3149. * \] R.F.V. Anderson, “The Weyl Functional Calculus”, J. Func. Anal. 4 (1969) 240-267. MR 58:30405. * \] P. Federbush, “e to the A, in a New Way”, math-ph/9903006, to be published in the Michigan Math. Journal.
warning/0004/math0004129.html	ar5iv	text	# Contents ## 1 Introduction An orbifold is a topological space locally modeled on the quotient of a smooth manifold by a finite group. Therefore, orbifolds belong to one of the simplest kinds of singular spaces. Orbifolds appear naturally in many branches of mathematics. For example, symplectic reduction often gives rise to orbifolds. An algebraic 3-fold with terminal singularities can be deformed into a symplectic orbifold. Orbifold also appears naturally in string theory, where many known Calabi-Yau 3-folds are the so called crepant resolutions of a Calabi-Yau orbifold. The physicists even attempted to formulate string theories on Calabi-Yau orbifolds which are expected to be “equivalent” to the string theories on its crepant resolutions \[DHVW\]. As a consequence of this orbifold string theory consideration, one has the following prediction that “orbifold quantum cohomology” is “isomorphic” to the ordinary quantum cohomology of its crepant resolutions. At this moment, even the physical idea around this subject is still vague and incomplete, particularly for the possible isomorphism. However, it seems that there are interesting new mathematical structures that are behind such orbifold string theories. This article is the first paper of a program to understand these new mathematical treasures behind orbifold string theory. We introduce orbifold cohomology groups of an almost complex orbifold, and orbifold Dolbeault cohomology groups of a complex orbifold. The main result of this paper is the construction of orbifold cup products on orbifold cohomology groups and orbifold Dolbeault cohomology groups, which make the corresponding total orbifold cohomology into a ring with unit. We will call the resulting rings orbifold cohomology ring or orbifold Dolbeault cohomology ring. (See Theorems 4.1.5 and 4.1.7 for details). In the case when the almost complex orbifold is closed and symplectic, the orbifold cohomology ring corresponds to the “classical part” of the orbifold quantum cohomology ring constructed in \[CR\]. Originally, this article is a small part of the much longer paper \[CR\] regarding the theory of orbifold quantum cohomology. However, we feel that the classical part (i.e. the orbifold cohomology) of the orbifold quantum cohomology is interesting in its own right, and technically, it is also much simpler to construct. Therefore, we decided to put it in a separate paper. A brief history is in order. In the case of Gorenstein global quotients, orbifold Euler characteristic-Hodge numbers have been extensively studied in the literature (see \[RO\],\[BD\],\[Re\] for a more complete reference). However, we would like to point out that (i) our orbifold cohomology is well-defined for any almost complex orbifold which may or may not be Gorenstein. Furthermore, it has an interesting feature that an orbifold cohomology class of a non-Gorenstein orbifold could have a rational degree (See examples in section 5); (ii) Even in the case of Gorensterin orbifolds, orbifold cohomology ring contains much more information than just orbifold Betti-Hodge numbers. In the case of global quotients, some constructions of this paper are already known to physicists. A notable exception is the orbifold cup product. On the other hand, many interesting orbifolds are not global quotients in general. For examples, most of Calabi-Yau hypersurfaces of weighted projective spaces are not global quotients. In this article, we systematically developed the theory (including the construction of orbifold cup products) for general orbifolds. Our construction of orbifold cup products is motivated by the construction of orbifold quantum cohomology. The second author would like to thank R. Dijkgraaf for bringing the orbifold string theory to his attention and E. Zaslow for valuable discussions. ## 2 Recollections on Orbifold In this section, we review basic definitions in the theory of orbifold. A systematic treatment of various aspects of differential geometry on orbifolds is contained in our forthcoming paper \[CR\]. The notion of orbifold was first introduced by Satake in \[S\], where a different name, V-manifold, was used. Let $`U`$ be a connected topological space, $`V`$ be a connected n-dimensional smooth manifold with a smooth action by a finite group $`G`$. Here we assume throughout that the fixed-point set of each element of the group is either the whole space or of codimension at least two. In particular, the action of $`G`$ does not have to be effective. This is the case, for example, when the action is orientation-preserving. This requirement has a consequence that the non-fixed-point set is locally connected if it is not empty. We will call the subgroup of $`G`$, which consists of elements fixing the whole space $`V`$, the kernel of the action. An n-dimensional uniformizing system of $`U`$ is a triple $`(V,G,\pi )`$, where $`\pi :VU`$ is a continuous map inducing a homeomorphism between the quotient space $`V/G`$ and $`U`$. Two uniformizing systems $`(V_i,G_i,\pi _i)`$, $`i=1,2`$, are isomorphic if there is a diffeomorphism $`\varphi :V_1V_2`$ and an isomorphism $`\lambda :G_1G_2`$ such that $`\varphi `$ is $`\lambda `$-equivariant, and $`\pi _2\varphi =\pi _1`$. If $`(\varphi ,\lambda )`$ is an automorphism of $`(V,G,\pi )`$, then there is $`gG`$ such that $`\varphi (x)=gx`$ and $`\lambda (a)=gag^1`$ for any $`xV`$ and $`aG`$. Note that here $`g`$ is unique iff the action of $`G`$ on $`V`$ is effective. Let $`i:U^{}U`$ be a connected open subset of $`U`$, and $`(V^{},G^{},\pi ^{})`$ be a uniformizing system of $`U^{}`$. We say that $`(V^{},G^{},\pi ^{})`$ is induced from $`(V,G,\pi )`$ if there is a monomorphism $`\tau :G^{}G`$ which is an isomorphism restricted to the kernels of the action of $`G^{}`$ and $`G`$ respectively, and a $`\tau `$-equivariant open embedding $`\psi :V^{}V`$ such that $`i\pi ^{}=\pi \psi `$. The pair $`(\psi ,\tau ):(V^{},G^{},\pi ^{})(V,G,\pi )`$ is called an injection. Two injections $`(\psi _i,\tau _i):(V_i^{},G_i^{},\pi _i^{})(V,G,\pi )`$, $`i=1,2`$, are isomorphic if there is an isomorphism $`(\varphi ,\lambda )`$ between $`(V_1^{},G_1^{},\pi _1^{})`$ and $`(V_2^{},G_2^{},\pi _2^{})`$, and an automorphism $`(\overline{\varphi },\overline{\lambda })`$ of $`(V,G,\pi )`$ such that $`(\overline{\varphi },\overline{\lambda })(\psi _1,\tau _1)=(\psi _2,\tau _2)(\varphi ,\lambda )`$. One can easily verify that for any connected open subset $`U^{}`$ of $`U`$, a uniformizing system $`(V,G,\pi )`$ of $`U`$ induces a unique isomorphism class of uniformizing systems of $`U^{}`$. Let $`U`$ be a connected and locally connected topological space. For any point $`pU`$, we can define the germ of uniformizing systems at $`p`$ in the following sense. Let $`(V_1,G_1,\pi _1)`$ and $`(V_2,G_2,\pi _2)`$ be uniformizing systems of neighborhoods $`U_1`$ and $`U_2`$ of $`p`$. We say that $`(V_1,G_1,\pi _1)`$ and $`(V_2,G_2,\pi _2)`$ are equivalent at $`p`$ if they induce isomorphic uniformizing systems for a neighborhood $`U_3`$ of $`p`$. Definition 2.1: 1. Let $`X`$ be a Hausdorff, second countable topological space. An n-dimensional orbifold structure on $`X`$ is given by the following data: for any point $`pX`$, there is a neighborhood $`U_p`$ and a n-dimensional uniformizing system $`(V_p,G_p,\pi _p)`$ of $`U_p`$ such that for any point $`qU_p`$, $`(V_p,G_p,\pi _p)`$ and $`(V_q,G_q,\pi _q)`$ are equivalent at $`q`$ (i.e., defining the same germ at $`q`$). With a given germ of orbifold structures, $`X`$ is called an orbifold. An open subset $`U`$ of $`X`$ is called a uniformized open subset if it is uniformized by a $`(V,G,\pi )`$ such that for each $`pU`$, $`(V,G,\pi )`$ defines the same germ with $`(V_p,G_p,\pi _p)`$ at $`p`$. We may assume that each $`V_p`$ is a n-ball centered at origin $`o`$ and $`\pi _p^1(p)=o`$. In particular, the origin $`o`$ is fixed by $`G_p`$. If $`G_p`$ acts effectively for every $`p`$, we call $`X`$ a reduced orbifold. 2. The notion of orbifold with boundary, in which we allow the uniformizing systems to be smooth manifolds with boundary, with a finite group action preserving the boundary, can be similarly defined. If $`X`$ is an orbifold with boundary, then it is easily seen that the boundary $`X`$ inherits an orbifold structure from $`X`$ and becomes an orbifold. $`\mathrm{}`$ In a certain sense, Satake’s definition of orbifold is less intrinsic than ours, although they are equivalent. In \[S\], an orbifold structure on $`X`$ is given by an open cover $`𝒰`$ of $`X`$ satisfying the following conditions: * Each element $`U`$ in $`𝒰`$ is uniformized, say by $`(V,G,\pi )`$. * If $`U^{}U`$, then there is a collection of injections $`(V^{},G^{},\pi ^{})(V,G,\pi )`$. * For any point $`pU_1U_2`$, $`U_1,U_2𝒰`$, there is a $`U_3𝒰`$ such that $`pU_3U_1U_2`$. It clearly defines an orbifold structure on $`X`$ in the sense of Definition 2.1. We will call such a cover of an orbifold $`X`$ a compatible cover if it gives rise to the same germ of orbifold structures on $`X`$. We remark that the orbifolds considered by Satake in \[S\] are all reduced. Now we consider a class of continuous maps between two orbifolds which respect the orbifold structures in a certain sense. Let $`U`$ be uniformized by $`(V,G,\pi )`$ and $`U^{}`$ by $`(V^{},G^{},\pi ^{})`$, and $`f:UU^{}`$ be a continuous map. A $`C^l`$ lifting, $`0l\mathrm{}`$, of $`f`$ is a $`C^l`$ map $`\stackrel{~}{f}:VV^{}`$ such that $`\pi ^{}\stackrel{~}{f}=f\pi `$, and for any $`gG`$, there is $`g^{}G^{}`$ so that $`g^{}\stackrel{~}{f}(x)=\stackrel{~}{f}(gx)`$ for any $`xV`$. Two liftings $`\stackrel{~}{f}_i:(V_i,G_i,\pi _i)(V_i^{},G_i^{},\pi _i^{})`$, $`i=1,2`$, are isomorphic if there exist isomorphisms $`(\varphi ,\tau ):(V_1,G_1,\pi _1)(V_2,G_2,\pi _2)`$ and $`(\varphi ^{},\tau ^{}):(V_1^{},G_1^{},\pi _1^{})(V_2^{},G_2^{},\pi _2^{})`$ such that $`\varphi ^{}\stackrel{~}{f}_1=\stackrel{~}{f}_2\varphi `$. Let $`pU`$ be any point. Then for any uniformized neighborhood $`U_p`$ of $`p`$ and uniformized neighborhood $`U_{f(p)}`$ of $`f(p)`$ such that $`f(U_p)U_{f(p)}`$, a lifting $`\stackrel{~}{f}`$ of $`f`$ will induce a lifting $`\stackrel{~}{f}_p`$ for $`f\|_{U_p}:U_pU_{f(p)}`$ as follows: For any injection $`(\varphi ,\tau ):(V_p,G_p,\pi _p)(V,G,\pi )`$, consider the map $`\stackrel{~}{f}\varphi :V_pV^{}`$. Observe that the inclusion $`\pi ^{}\stackrel{~}{f}\varphi (V_p)U_{f(p)}`$ implies that $`\stackrel{~}{f}\varphi (V_p)`$ lies in $`(\pi ^{})^1(U_{f(p)})`$. Therefore there is an injection $`(\varphi ^{},\tau ^{}):(V_{f(p)},G_{f(p)},\pi _{f(p)})(V^{},G^{},\pi ^{})`$ such that $`\stackrel{~}{f}\varphi (V_p)\varphi ^{}(V_{f(p)})`$. We define $`\stackrel{~}{f}_p=(\varphi ^{})^1\stackrel{~}{f}\varphi `$. In this way we obtain a lifting $`\stackrel{~}{f}_p:(V_p,G_p,\pi _p)(V_{f(p)},G_{f(p)},\pi _{f(p)})`$ for $`f\|_{U_p}:U_pU_{f(p)}`$. We can verify that different choices give isomorphic liftings. We define the germ of liftings as follows: two liftings are equivalent at $`p`$ if they induce isomorphic liftings on a smaller neighborhood of $`p`$. Let $`f:XX^{}`$ be a continuous map between orbifolds $`X`$ and $`X^{}`$. A lifting of $`f`$ consists of following data: for any point $`pX`$, there exist charts $`(V_p,G_p,\pi _p)`$ at $`p`$ and $`(V_{f(p)},G_{f(p)},\pi _{f(p)})`$ at $`f(p)`$ and a lifting $`\stackrel{~}{f}_p`$ of $`f\|_{\pi _p(V_p)}:\pi _p(V_p)\pi _{f(p)}(V_{f(p)})`$ such that for any $`q\pi _p(V_p)`$, $`\stackrel{~}{f}_p`$ and $`\stackrel{~}{f}_q`$ induce the same germ of liftings of $`f`$ at $`q`$. We can define the germ of liftings in the sense that two liftings of $`f`$, $`\{\stackrel{~}{f}_{p,i}:(V_{p,i},G_{p,i},\pi _{p,i})(V_{f(p),i},G_{f(p),i},\pi _{f(p),i}):pX\}`$, $`i=1,2`$, are equivalent if for each $`pX`$, $`\stackrel{~}{f}_{p,i},i=1,2`$, induce the same germ of liftings of $`f`$ at $`p`$. Definition 2.2: A $`C^l`$ map ($`0l\mathrm{}`$) between orbifolds $`X`$ and $`X^{}`$ is a germ of $`C^l`$ liftings of a continuous map between $`X`$ and $`X^{}`$. $`\mathrm{}`$ We denote by $`\stackrel{~}{f}`$ a $`C^l`$ map which is a germ of liftings of a continuous map $`f`$. Our definition of $`C^l`$ maps corresponds to the notion of $`V`$-maps in \[S\]. Next we describe the notion of orbifold vector bundle, which corresponds to the notion of smooth vector bundle over manifolds. When there is no confusion, we will simply call it an orbifold bundle. We begin with local uniformizing systems for orbifold bundles. Given a uniformized topological space $`U`$ and a topological space $`E`$ with a surjective continuous map $`pr:EU`$, a uniformizing system of rank $`k`$ orbifold bundle for $`E`$ over $`U`$ consists of the following data: * A uniformizing system $`(V,G,\pi )`$ of $`U`$. * A uniformizing system $`(V\times 𝐑^k,G,\stackrel{~}{\pi })`$ for $`E`$. The action of $`G`$ on $`V\times 𝐑^k`$ is an extension of the action of $`G`$ on $`V`$ given by $`g(x,v)=(gx,\rho (x,g)v)`$ where $`\rho :V\times GAut(𝐑^k)`$ is a smooth map satisfying: $$\rho (gx,h)\rho (x,g)=\rho (x,hg),g,hG,xV.$$ * The natural projection map $`\stackrel{~}{pr}:V\times 𝐑^kV`$ satisfies $`\pi \stackrel{~}{pr}=pr\stackrel{~}{\pi }`$. We can similarly define isomorphisms between uniformizing systems of orbifold bundle for $`E`$ over $`U`$. The only additional requirement is that the diffeomorphisms between $`V\times 𝐑^k`$ are linear on each fiber of $`\stackrel{~}{pr}:V\times 𝐑^kV`$. Moreover, for each connected open subset $`U^{}`$ of $`U`$, we can similarly prove that there is a unique isomorphism class of induced uniformizing systems of orbifold bundle for $`E^{}=pr^1(U^{})`$ over $`U^{}`$. The germ of uniformizing systems of orbifold bundle at a point $`pU`$ can be also similarly defined. Definition 2.3: 1. Let $`X`$ be an orbifold and $`E`$ be a topological space with a surjective continuous map $`pr:EX`$. A rank $`k`$ orbifold bundle structure on $`E`$ over $`X`$ consists of following data: For each point $`pX`$, there is a uniformized neighborhood $`U_p`$ and a uniformizing system of rank $`k`$ orbifold bundle for $`pr^1(U_p)`$ over $`U_p`$ such that for any $`qU_p`$, the uniformizing systems of orbifold bundle over $`U_p`$ and $`U_q`$ define the same germ at $`q`$. The topological space $`E`$ with a given germ of orbifold bundle structures becomes an orbifold ($`E`$ is obviously Hausdorff and second countable) and is called an orbifold bundle over $`X`$. Each chart $`(V_p\times 𝐑^k,G_p,\stackrel{~}{\pi }_p)`$ is called a local trivialization of $`E`$. At each point $`pX`$, the fiber $`E_p=pr^1(p)`$ is isomorphic to $`𝐑^k/G_p`$. It contains a linear subspace $`E^p`$ of fixed points of $`G_p`$. 2. The notion of orbifold bundle over an orbifold with boundary is similarly defined. One can easily verify that if $`pr:EX`$ is an orbifold bundle over an orbifold with boundary $`X`$, then the restriction to the boundary $`X`$, $`E_X=pr^1(X)`$, is an orbifold bundle over $`X`$. 3. One can define fiber orbifold bundle in the same vein. $`\mathrm{}`$ A $`C^l`$ map $`\stackrel{~}{s}`$ from $`X`$ to an orbifold bundle $`pr:EX`$ is called a $`C^l`$ section if locally $`\stackrel{~}{s}`$ is given by $`\stackrel{~}{s}_p:V_pV_p\times 𝐑^k`$ where $`\stackrel{~}{s}_p`$ is $`G_p`$-equivariant and $`\stackrel{~}{pr}\stackrel{~}{s}_p=Id`$ on $`V_p`$. We observe that * For each point $`p`$, $`s(p)`$ lies in $`E^p`$, the linear subspace of fixed points of $`G_p`$. * The space of all $`C^l`$ sections of $`E`$, denoted by $`C^l(E)`$, has a structure of vector space over $`𝐑`$ (or $`𝐂`$) as well as a $`C^l(X)`$-module structure. * The $`C^l`$ sections $`\stackrel{~}{s}`$ are in $`1:1`$ correspondence with the underlying continuous maps $`s`$. Orbifold bundles are more conveniently described by transition maps, e.g. as in \[S\]. More precisely, an orbifold bundle over an orbifold $`X`$ can be constructed from the following data: A compatible cover $`𝒰`$ of $`X`$ such that for any injection $`i:(V^{},G^{},\pi ^{})(V,G,\pi )`$, there is a smooth map $`g_i:V^{}Aut(𝐑^k)`$ giving an open embedding $`V^{}\times 𝐑^kV\times 𝐑^k`$ by $`(x,v)(i(x),g_i(x)v)`$, and for any composition of injections $`ji`$, we have $$g_{ji}(x)=g_j(i(x))g_i(x).$$ $`(2.1)`$ Two collections of maps $`g^{(1)}`$ and $`g^{(2)}`$ define isomorphic orbifold bundles if there are maps $`\delta _V:VAut(𝐑^k)`$ such that for any injection $`i:(V^{},G^{},\pi ^{})(V,G,\pi )`$, we have $$g_i^{(2)}(x)=\delta _V(i(x))g_i^{(1)}(x)(\delta _V^{}(x))^1,xV^{}.$$ $`(2.2)`$ Since the equation (2.1) behaves naturally under constructions of vector spaces such as tensor product, exterior product, etc., we can define the corresponding constructions for orbifold bundles. Example 2.4: For an orbifold $`X`$, the tangent bundle $`TX`$ can be constructed because the differential of any injection satisfies the equation (2.1). Likewise, we define cotangent bundle $`T^{}X`$, the bundles of exterior power or tensor product. The $`C^{\mathrm{}}`$ sections of these bundles give us vector fields, differential forms or tensor fields on $`X`$. We remark that if $`\omega `$ is a differential form on $`X^{}`$ and $`\stackrel{~}{f}:XX^{}`$ is a $`C^{\mathrm{}}`$ map, then there is a pull-back form $`\stackrel{~}{f}^{}\omega `$ on $`X`$. Let $`U`$ be an open subset of an orbifold $`X`$ with an orbifold structure $`\{(V_p,G_p,\pi _p):pX\}`$, then $`\{(V_p^{},G_p^{},\pi _p^{}):pU\}`$ is an orbifold structure on $`U`$ where $`(V_p^{},G_p^{},\pi _p^{})`$ is a uniformizing system of $`\pi _p(V_p)U`$ induced from $`(V_p,G_p,\pi _p)`$. Likewise, let $`pr:EX`$ be an orbifold bundle and $`U`$ an open subset of $`X`$, then $`pr:E_U=pr^1(U)U`$ inherits a unique germ of orbifold bundle structures from $`E`$, called the restriction of $`E`$ over $`U`$. When $`U`$ is a uniformized open set in $`X`$, say uniformized by $`(V,G,\pi )`$, then there is a smooth vector bundle $`E_V`$ over $`V`$ with a smooth action of $`G`$ such that $`(E_V,G,\stackrel{~}{\pi })`$ uniformizes $`E_U`$. This is seen as follows: We first take a compatible cover $`𝒰`$ of $`U`$, fine enough so that the preimage under $`\pi `$ is a compatible cover of $`V`$. Let $`E_U`$ be given by a set of transition maps with respect to $`𝒰`$ satisfying (2.1), then the pull-backs under $`\pi `$ form a set of transition maps with respect to $`\pi ^1(𝒰)`$ with an action of $`G`$ by permutations, also satisfying (2.1), so that it defines a smooth vector bundle over $`V`$ with a compatible smooth action of $`G`$. Any $`C^l`$ section of $`E`$ on $`X`$ restricts to a $`C^l`$ section of $`E_U`$ on $`U`$, and when $`U`$ is a uniformized open set by $`(V,G,\pi )`$, it lifts to a $`G`$-equivariant $`C^l`$ section of $`E_V`$ on $`V`$. Integration over orbifolds is defined as follows. Let $`U`$ be a connected n-dimensional orbifold, which is uniformized by $`(V,G,\pi )`$, with the kernel of the action of $`G`$ on $`V`$ denoted by $`K`$. For any compact supported differential n-form $`\omega `$ on $`U`$, which is, by definition, a $`G`$-equivariant compact supported n-form $`\stackrel{~}{\omega }`$ on $`V`$, the integration of $`\omega `$ on $`U`$ is defined by $$_U^{orb}\omega :=\frac{1}{\|G\|}_V\stackrel{~}{\omega },$$ $`(2.3)`$ where $`\|G\|`$ is the order of the group $`G`$. In general, let $`X`$ be an orbifold. Fix a $`C^{\mathrm{}}`$ partition of unity $`\{\rho _i\}`$ subordinated to $`\{U_i\}`$ where each $`U_i`$ is a uniformized open set in $`X`$. Then the integration over $`X`$ is defined by $$_X^{orb}\omega :=\underset{i}{}_{U_i}^{orb}\rho _i\omega ,$$ $`(2.4)`$ which is independent of the choice of the partition of unity $`\{\rho _i\}`$. We remark that it is important throughout this paper that we adopt the integration over orbifolds as in $`(2.3)`$ and $`(2.4)`$, where we divide the integral over the uniformizer $`V`$ by the group order $`\|G\|`$ instead of $`\|G\|/\|K\|`$ ($`K`$ is the kernel of the action). As a result, the fundamental class of an orbifold is rational in general. The integration $`^{orb}`$ coincides with the usual measure-theoretic integration if and only if the orbifold is reduced. The de Rham cohomology groups of an orbifold are defined similarly through differential forms, which are naturally isomorphic to the singular cohomology groups with real coefficients. For an oriented, closed orbifold, the singular cohomology groups are naturally isomorphic to the intersection homology groups, both with rational coefficients, for which the Poincaré duality is valid \[GM\]. Characteristic classes (Euler class for oriented orbifold bundles, Chern classes for complex orbifold bundles, and Pontrjagin classes for real orbifold bundles) are well-defined for orbifold bundles. One way to define them is through Chern-Weil theory, so that the characteristic classes take values in the deRham cohomology groups. Another way to define them is through the transgressions in the Serre spectral sequences with rational coefficients of the associated Stiefel orbifold bundles, so that these characteristic classes are defined over the rationals \[K1\]. ## 3 Orbifold Cohomology Groups In this section, we introduce the main object of study, the orbifold cohomology groups of an almost complex orbifold. ### 3.1 Twisted sectors Let $`X`$ be an orbifold. For any point $`pX`$, let $`(V_p,G_p,\pi _p)`$ be a local chart at $`p`$. Consider the set of pairs: $$\stackrel{~}{X}=\{(p,(g)_{G_p})\|pX,gG_p\},$$ $`(\mathrm{3.1.1})`$ where $`(g)_{G_p}`$ is the conjugacy class of $`g`$ in $`G_p`$. If there is no confusion, we will omit the subscript $`G_p`$ to simplify the notation. There is a surjective map $`\pi :\stackrel{~}{X}X`$ defined by $`(p,(g))p`$. Lemma 3.1.1 (Kawasaki,\[K1\]): The set $`\stackrel{~}{X}`$ is naturally an orbifold (not necessarily connected) with an orbifold structure given by $$\{\pi _{p,g}:(V_p^g,C(g))V_p^g/C(g):pX,gG_p.\},$$ where $`V_p^g`$ is the fixed-point set of $`g`$ in $`V_p`$, $`C(g)`$ is the centralizer of $`g`$ in $`G_p`$. Moreover, if $`X`$ is closed, so is $`\stackrel{~}{X}`$. Under this orbifold structure, the map $`\pi :\stackrel{~}{X}X`$ is a $`C^{\mathrm{}}`$ map. Proof: First we identify a point $`(q,(h))`$ in $`\stackrel{~}{X}`$ as a point in $`_{\{(g),gG_p\}}V_p^g/C(g)`$ if $`qU_p`$ for some $`pX`$. Pick a representative $`yV_p`$ such that $`\pi _p(y)=q`$. Then this gives rise to a monomorphism $`\lambda _y:G_qG_p`$. Pick a representative $`hG_q`$ for $`(h)`$ in $`G_q`$, we let $`g=\lambda _y(h)`$. Then $`yV_p^g`$. So we have a map $`\mathrm{\Phi }:(q,h)(y,g)(V_p^g,G_p)`$. If we change $`h`$ by a $`h^{}=a^1haG_q`$ for $`aG_q`$, then $`g`$ is changed to $`\lambda _y(a^1ha)=\lambda _y(a)^1g\lambda _y(a)`$. So we have $`\mathrm{\Phi }:(q,a^1ha)(y,\lambda _y(a)^1g\lambda _y(a))(V_p^{\lambda _y(a)^1g\lambda _y(a)},G_p)`$. (Note that $`\lambda _y`$ is determined up to conjugacy by an element in $`G_q`$.) If we take a different representative $`y^{}V_p`$ such that $`\pi _p(y^{})=q`$, and suppose $`y^{}=by`$ for some $`bG_p`$. Then we have a different identification $`\lambda _y^{}:G_qG_p`$ of $`G_q`$ as a subgroup of $`G_p`$ where $`\lambda _y^{}=b\lambda _yb^1`$. In this case, we have $`\mathrm{\Phi }:(q,h)(y^{},bgb^1)(V_p^{bgb^1},G_p)`$. If $`g=bgb^1`$, then $`bC(g)`$. In any event, $`\mathrm{\Phi }`$ induces a map $`\varphi `$ sending $`(q,(h))`$ to a point in $`_{\{(g),gG_p\}}V_p^g/C(g)`$. It is one to one because if $`\varphi (q_1,(h_1))=\varphi (q_2,(h_2))`$, then we may assume that $`\mathrm{\Phi }(q_1,h_1)=\mathrm{\Phi }(q_2,h_2)`$ after applying conjugations. But this means that $`(q_1,h_1)=(q_2,h_2)`$. It is easily seen that this map $`\varphi `$ is also onto. Hence we have shown that $`\stackrel{~}{X}`$ is covered by $`_{\{pX\}}_{\{(g),gG_p\}}V_p^g/C(g)`$. We define a topology on $`\stackrel{~}{X}`$ so that each $`V_p^g/C(g)`$ is an open subset for any $`(p,g)`$ where $`pX`$ and $`gG_p`$. We also uniformize $`V_p^g/C(g)`$ by $`(V_p^g,C(g))`$. It remains to show that these charts fit together to form an orbifold structure on $`\stackrel{~}{X}`$. Let $`xV_p^g/C(g)`$ and take a representative $`\stackrel{~}{x}`$ in $`V_p^g`$. Let $`H_x`$ be the isotropy subgroup of $`\stackrel{~}{x}`$ in $`C(g)`$. Then $`(V_p^g,C(g))`$ induces a germ of uniformizing system at $`x`$ as $`(B_x,H_x)`$ where $`B_x`$ is a small ball in $`V_p^g`$ centered at $`\stackrel{~}{x}`$. Let $`\pi _p(\stackrel{~}{x})=q`$. We need to write $`(B_x,H_x)`$ as $`(V_q^h,C(h))`$ for some $`hG_q`$. We let $`\lambda _x:G_qG_p`$ be an induced monomorphism resulted from choosing $`\stackrel{~}{x}`$ as the representative of $`q`$ in $`V_p`$. We define $`h=\lambda _x^1(g)`$ ($`g`$ is in $`\lambda _x(G_q)`$ since $`\stackrel{~}{x}V_p^g`$ and $`\pi _p(\stackrel{~}{x})=q`$.) Then we can identify $`B_x`$ as $`V_q^h`$. We also see that $`H_x=\lambda _x(C(h))`$. Therefore $`(B_x,H_x)`$ is identified as $`(V_q^h,C(h))`$. The map $`\pi :\stackrel{~}{X}X`$ is obviously continuous with the given topology of $`\stackrel{~}{X}`$, and actually is a $`C^{\mathrm{}}`$ map with the given orbifold structure on $`\stackrel{~}{X}`$ with the local liftings given by embeddings $`V_p^gV_p`$. We finish the proof by showing that $`\stackrel{~}{X}`$ is Hausdorff and second countable with the given topology. Let $`(p,(g))`$ and $`(q,(h))`$ be distinct two points in $`\stackrel{~}{X}`$. When $`pq`$, there are $`U_p`$, $`U_q`$ such that $`U_pU_q=\mathrm{}`$ since $`X`$ is Hausdorff. It is easily seen that in this case $`(p,(g))`$ and $`(q,(h))`$ are separated by disjoint neighborhoods $`\pi ^1(U_p)`$ and $`\pi ^1(U_q)`$, where $`\pi :\stackrel{~}{X}X`$. When $`p=q`$, we must then have $`(g)(h)`$. In this case, $`(p,(g))`$ and $`(q,(h))`$ lie in different open subsets $`V_p^g/C(g)`$ and $`V_q^h/C(h)`$ respectively. Hence $`\stackrel{~}{X}`$ is Hausdorff. The second countability of $`\stackrel{~}{X}`$ follows from the second countability of $`X`$ and the fact that $`\pi ^1(U_p)`$ is a finite union of open subsets of $`\stackrel{~}{X}`$ for each $`pX`$ and a uniformized neighborhood $`U_p`$ of $`p`$. $`\mathrm{}`$ Next, we would like to describe the connected components of $`\stackrel{~}{X}`$. Recall that every point $`p`$ has a local chart $`(V_p,G_p,\pi _p)`$ which gives a local uniformized neighborhood $`U_p=\pi _p(V_p)`$. If $`qU_p`$, up to conjugation, there is an injective homomorphism $`G_qG_p`$. For $`gG_q`$, the conjugacy class $`(g)_{G_p}`$ is well-defined. We define an equivalence relation $`(g)_{G_q}(g)_{G_p}`$. Let $`T`$ be the set of equivalence classes. To abuse the notation, we often use $`(g)`$ to denote the equivalence class which $`(g)_{G_q}`$ belongs to. It is clear that $`\stackrel{~}{X}`$ is decomposed as a disjoint union of connected components $$\stackrel{~}{X}=\underset{(g)T}{}X_{(g)},$$ $`(\mathrm{3.1.2})`$ where $$X_{(g)}=\{(p,(g^{})_{G_p})\|g^{}G_p,(g^{})_{G_p}(g)\}.$$ $`(\mathrm{3.1.3})`$ Definition 3.1.2: $`X_{(g)}`$ for $`g1`$ is called a twisted sector. Furthermore, we call $`X_{(1)}=X`$ the nontwisted sector. Example 3.1.3: Consider the case that the orbifold $`X=Y/G`$ is a global quotient. We will show that $`\stackrel{~}{X}`$ can be identified with $`_{\{(g),gG\}}Y^g/C(g)`$ where $`Y^g`$ is the fixed-point set of element $`gG`$. Let $`\pi :\stackrel{~}{X}X`$ be the surjective map defined by $`(p,(g))p`$. Then for any $`pX`$, the preimage $`\pi ^1(p)`$ in $`\stackrel{~}{X}`$ has a neighborhood described by $`W_p=_{\{(g),gG_p\}}V_p^g/C(g)`$, which is uniformized by $`\widehat{W}_p=_{\{(g),gG_p\}}V_p^g`$. For each $`pX`$, pick a $`yY`$ that represents $`p`$, and an injection $`(\varphi _p,\lambda _p):(V_p,G_p)(Y,G)`$ whose image is centered at $`y`$. This induces an open embedding $`\stackrel{~}{f}_p:\widehat{W}_p_{\{(\lambda _p(g)),\lambda _p(g)G\}}Y^{\lambda _p(g)}_{\{(g),gG\}}Y^g`$, which induces a homeomorphism $`f_p`$ from $`W_p`$ into $`_{\{(g),gG\}}Y^g/C(g)`$ that is independent of the choice of $`y`$ and $`(\varphi _p,\lambda _p)`$. These maps $`\{f_p;pX\}`$ fit together to define a map $`f:\stackrel{~}{X}_{\{(g),gG\}}Y^g/C(g)`$ which we can verify to be a homeomorphism. $`\mathrm{}`$ Remark 3.1.4: There is a natural $`C^{\mathrm{}}`$ map $`I:\stackrel{~}{X}\stackrel{~}{X}`$ defined by $$I((p,(g)_{G_p}))=(p,(g^1)_{G_p}).$$ $`(\mathrm{3.1.4})`$ The map $`I`$ is an involution (i.e., $`I^2=Id`$) which induces an involution on the set $`T`$ of equivalence classes of relations $`(g)_{G_q}(g)_{G_p}`$. We denoted by $`(g^1)`$ the image of $`(g)`$ under this induced map. ### 3.2 Degree shifting and orbifold cohomology group For the rest of the paper, we will assume that $`X`$ is an almost complex orbifold with an almost complex structure $`J`$. Recall that an almost complex structure $`J`$ on $`X`$ is a smooth section of the orbifold bundle $`End(TX)`$ such that $`J^2=Id`$. Observe that $`\stackrel{~}{X}`$ naturally inherits an almost complex structure from the one on $`X`$, and the map $`\pi :\stackrel{~}{X}X`$ defined by $`(p,(g)_{G_p})p`$ is naturally pseudo-holomorphic, i.e., its differential commutes with the almost complex structures on $`\stackrel{~}{X}`$ and $`X`$. An important feature of orbifold cohomology groups is degree shifting, which we shall explain now. Let $`p`$ be any point of $`X`$. The almost complex structure on $`X`$ gives rise to a representation $`\rho _p:G_pGL(n,𝐂)`$ (here $`n=dim_𝐂X`$). For any $`gG_p`$, we write $`\rho _p(g)`$ as a diagonal matrix $$diag(e^{2\pi im_{1,g}/m_g},\mathrm{},e^{2\pi im_{n,g}/m_g}),$$ where $`m_g`$ is the order of $`\rho _p(g)`$, and $`0m_{i,g}<m_g`$. This matrix depends only on the conjugacy class $`(g)_{G_p}`$ of $`g`$ in $`G_p`$. We define a function $`\iota :\stackrel{~}{X}𝐐`$ by $$\iota (p,(g)_{G_p})=\underset{i=1}{\overset{n}{}}\frac{m_{i,g}}{m_g}.$$ It is straightforward to show the following Lemma 3.2.1: The function $`\iota :X_{(g)}𝐐`$ is constant. Its constant value, which will be denoted by $`\iota _{(g)}`$, satisfies the following conditions: * $`\iota _{(g)}`$ is integral if and only if $`\rho _p(g)SL(n,𝐂)`$. * $$\iota _{(g)}+\iota _{(g^1)}=rank(\rho _p(g)I),$$ $`(\mathrm{3.2.1})`$ which is the “complex codimension” $`dim_𝐂Xdim_𝐂X_{(g)}=ndim_𝐂X_{(g)}`$ of $`X_{(g)}`$ in $`X`$. As a consequence, $`\iota _{(g)}+dim_𝐂X_{(g)}<n`$ when $`\rho _p(g)I`$. Definition 3.2.2: $`\iota _{(g)}`$ is called a degree shifting number. In the definition of orbifold cohomology groups, we will shift up the degree of cohomology classes of $`X_{(g)}`$ by $`2\iota _{(g)}`$. The reason for such a degree shifting will become clear after we discuss the dimension of moduli space of ghost maps (see formula (4.2.14)). An orbifold $`X`$ is called a $`SL`$-orbifold if $`\rho _p(g)SL(n,𝐂)`$ for all $`pX`$ and $`gG_p`$, and called a $`SP`$-orbifold if $`\rho _p(g)SP(n,𝐂)`$. In particular, a Calabi-Yau orbifold is a $`SL`$-orbifold, and a holomorphic symplectic orbifold or hyperkahler orbifold is a $`SP`$-orbifold. By Lemma 3.2.1, $`\iota _{(g)}`$ is integral if and only if $`X`$ is a $`SL`$-orbifold. We observe that although the almost complex structure $`J`$ is involved in the definition of degree shifting numbers $`\iota _{(g)}`$, they do not depend on $`J`$ because locally the parameter space of almost complex structures, which is the coset $`SO(2n,𝐑)/U(n,𝐂)`$, is connected. Definition 3.2.3: We define the orbifold cohomology groups $`H_{orb}^d(X)`$ of $`X`$ by $$H_{orb}^d(X)=_{(g)T}H^{d2\iota _{(g)}}(X_{(g)})$$ $`(\mathrm{3.2.2})`$ and orbifold Betti numbers $`b_{orb}^d=_{(g)}dimH^{d2\iota _{(g)}}(X_{(g)})`$. Here each $`H^{}(X_{(g)})`$ is the singular cohomology of $`X_{(g)}`$ with real coefficients, which is isomorphic to the corresponding de Rham cohomology group. As a consequence, the cohomology classes can be represented by closed differential forms on $`X_{(g)}`$. Note that, in general, orbifold cohomology groups are rationally graded. Suppose $`X`$ is a complex orbifold with an integrable complex structure $`J`$. Then each twisted sector $`X_{(g)}`$ is also a complex orbifold with the induced complex structure. We consider the Čech cohomology groups on $`X`$ and each $`X_{(g)}`$ with coefficients in the sheaves of holomorphic forms (in the orbifold sense). These Čech cohomology groups are identified with the Dolbeault cohomology groups of $`(p,q)`$-forms (in the orbifold sense). When $`X`$ is closed, the harmonic theory \[Ba\] can be applied to show that these groups are finite dimensional, and there is a Kodaira-Serre duality between them. When $`X`$ is a closed Kahler orbifold (so is each $`X_{(g)}`$), these groups are then related to the singular cohomology groups of $`X`$ and $`X_{(g)}`$ as in the smooth case, and the Hodge decomposition theorem holds for these cohomology groups. Definition 3.2.4: Let $`X`$ be a complex orbifold. We define, for $`0p,qdim_𝐂X`$, orbifold Dolbeault cohomology groups $$H_{orb}^{p,q}(X)=_{(g)}H^{p\iota _{(g)},q\iota _{(g)}}(X_{(g)}).$$ $`(\mathrm{3.2.3})`$ We define orbifold Hodge numbers by $`h_{orb}^{p,q}(X)=dimH_{orb}^{p,q}(X)`$. Remark 3.2.5: We can define compact supported orbifold cohomology groups $`H_{orb,c}^{}(X),H_{orb,c}^,(X)`$ in the obvious fashion. ### 3.3 Poincaré duality Recall that there is a natural $`C^{\mathrm{}}`$ map $`I:X_{(g)}X_{(g^1)}`$ defined by $`(p,(g))(p,(g^1))`$, which is an automorphism of $`\stackrel{~}{X}`$ as an orbifold and $`I^2=Id`$ (Remark 3.1.4). Proposition 3.3.1: (Poincaré duality) For any $`0d2n`$, the pairing $$<>_{orb}:H_{orb}^d(X)\times H_{orb,c}^{2nd}(X)𝐑$$ defined by the direct sum of $$<>_{orb}^{(g)}:H^{d2\iota _{(g)}}(X_{(g)})\times H_c^{2nd2\iota _{(g^1)}}(X_{(g^1)})𝐑$$ where $$<\alpha ,\beta >_{orb}^{(g)}=_{X_{(g)}}^{orb}\alpha I^{}(\beta )$$ $`(\mathrm{3.3.4})`$ for $`\alpha H^{d2\iota _{(g)}}(X_{(g)}),\beta H_c^{2nd2\iota _{(g^1)}}(X_{(g^1)})`$ is nondegenerate. Here the integral in the right hand side of $`(\mathrm{3.3.4})`$ is defined using $`(2.4)`$. Note that $`<>_{orb}`$ equals the ordinary Poincaré pairing when restricted to the nontwisted sectors $`H^{}(X)`$. Proof: By (3.2.1), we have $$2nd2\iota _{(g^1)}=dimX_{(g)}d2\iota _{(g)}.$$ Furthermore, $`I\|_{X_{(g)}}:X_{(g)}X_{(g^1)}`$ is a homeomorphism. Under this homeomorphism, $`<>_{orb}^{(g)}`$ is isomorphic to the ordinary Poincaré pairing on $`X_{(g)}`$. Hence $`<>_{orb}`$ is nondegenerate. $`\mathrm{}`$ For the case of orbifold Dolbeault cohomology, the following proposition is straightforward. Proposition 3.3.2: Let $`X`$ be an $`n`$-dimensional complex orbifold. There is a Kodaira-Serre duality pairing $$<>_{orb}:H_{orb}^{p,q}(X)\times H_{orb,c}^{np,nq}(X)𝐂$$ similarly defined as in the previous proposition. When $`X`$ is closed and Kahler, the following is true: * $`H_{orb}^r(X)𝐂=_{r=p+q}H_{orb}^{p,q}(X)`$ * $`H_{orb}^{p,q}(X)=\overline{H_{orb}^{q,p}(X)}`$, and the two pairings (Poincaré and Kodaira-Serre) coincide. ## 4 Orbifold Cup Product and Orbifold Cohomology Ring ### 4.1 Orbifold cup product In this section, we give an explicit definition of the orbifold cup product. Its interpretation in terms of Gromov-Witten invariants and the proof of associativity of the product will be given in subsequent sections. Let $`X`$ be an orbifold, and $`(V_p,G_p,\pi _p)`$ be a uniformizing system at point $`pX`$. We define the $`k`$-multi-sector of $`X`$, which is denoted by $`\stackrel{~}{X}_k`$, to be the set of all pairs $`(p,(𝐠))`$, where $`pX`$, $`𝐠=(g_1,\mathrm{},g_k)`$ with each $`g_iG_p`$, and $`(𝐠)`$ stands for the conjugacy class of $`𝐠=(g_1,\mathrm{},g_k)`$. Here two $`k`$-tuple $`(g_1^{(i)},\mathrm{},g_k^{(i)})`$, $`i=1,2`$, are conjugate if there is a $`gG_p`$ such that $`g_j^{(2)}=gg_j^{(1)}g^1`$ for all $`j=1,\mathrm{},k`$. Lemma 4.1.1: The $`k`$-multi-sector $`\stackrel{~}{X}_k`$ is naturally an orbifold, with the orbifold structure given by $$\{\pi _{p,𝐠}:(V_p^𝐠,C(𝐠))V_p^𝐠/C(𝐠)\},$$ $`(\mathrm{4.1.1})`$ where $`V_p^𝐠=V_p^{g_1}V_p^{g_2}\mathrm{}V_p^{g_k}`$, $`C(𝐠)=C(g_1)C(g_2)\mathrm{}C(g_k)`$. Here $`𝐠=(g_1,\mathrm{},g_k)`$, $`V_p^g`$ stands for the fixed-point set of $`gG_p`$ in $`V_p`$, and $`C(g)`$ for the centralizer of $`g`$ in $`G_p`$. For each $`i=1,\mathrm{},k`$, there is a $`C^{\mathrm{}}`$ map $`e_i:\stackrel{~}{X}_k\stackrel{~}{X}`$ defined by sending $`(p,(𝐠))`$ to $`(p,(g_i))`$ where $`𝐠=(g_1,\mathrm{},g_k)`$. When $`X`$ is almost complex, $`\stackrel{~}{X}_k`$ inherits an almost complex structure from $`X`$, and when $`X`$ is closed, $`\stackrel{~}{X}_k`$ is a finite disjoint union of closed orbifolds. Proof: The proof is parallel to the proof of Lemma 3.1.1 where $`\stackrel{~}{X}`$ is shown to be an orbifold. First we identify a point $`(q,(𝐡))`$ in $`\stackrel{~}{X}_k`$ as a point in $`_{\{(p,(𝐠))\stackrel{~}{X}_k\}}V_p^𝐠/C(𝐠)`$ if $`qU_p`$. Pick a representative $`yV_p`$ such that $`\pi _p(y)=q`$. Then this gives rise to a monomorphism $`\lambda _y:G_qG_p`$. Pick a representative $`𝐡=(h_1,\mathrm{},h_k)G_q\times \mathrm{}\times G_q`$ for $`(𝐡)`$, we let $`𝐠=\lambda _y(𝐡)`$. Then $`yV_p^𝐠`$. So we have a map $`\theta :(q,𝐡)(y,𝐠)`$. If we change $`𝐡`$ by $`𝐡^{}=a^1𝐡a`$ for some $`aG_q`$, then $`𝐠`$ is changed to $`\lambda _y(a^1𝐡a)=\lambda _y(a)^1𝐠\lambda _y(a)`$. So we have $`\theta :(q,a^1𝐡a)(y,\lambda _y(a)^1𝐠\lambda _y(a))`$ where $`y`$ is regarded as a point in $`V_p^{\lambda _y(a)^1𝐠\lambda _y(a)}`$. (Note that $`\lambda _y`$ is determined up to conjugacy by an element in $`G_q`$.) If we take a different representative $`y^{}V_p`$ such that $`\pi _p(y^{})=q`$, and suppose $`y^{}=by`$ for some $`bG_p`$. Then we have a different identification $`\lambda _y^{}:G_qG_p`$ of $`G_q`$ as a subgroup of $`G_p`$ where $`\lambda _y^{}=b\lambda _yb^1`$. In this case, we have $`\theta :(q,𝐡)(y^{},b𝐠b^1)`$ where $`y^{}V_p^{b𝐠b^1}`$. If $`𝐠=b𝐠b^1`$, then $`bC(𝐠)`$. Therefore we have shown that $`\theta `$ induces a map sending $`(q,(𝐡))`$ to a point in $`_{\{(p,(𝐠))\stackrel{~}{X}_k\}}V_p^𝐠/C(𝐠)`$, which can be similarly shown to be one to one and onto. Hence we have shown that $`\stackrel{~}{X}_k`$ is covered by $`_{\{(p,(𝐠))\stackrel{~}{X}_k\}}V_p^𝐠/C(𝐠)`$. We define a topology on $`\stackrel{~}{X}_k`$ so that each $`V_p^𝐠/C(𝐠)`$ is an open subset for any $`(p,𝐠)`$. We also uniformize $`V_p^𝐠/C(𝐠)`$ by $`(V_p^𝐠,C(𝐠))`$. It remains to show that these charts fit together to form an orbifold structure on $`\stackrel{~}{X}_k`$. Let $`xV_p^𝐠/C(𝐠)`$ and take a representative $`\stackrel{~}{x}`$ in $`V_p^𝐠`$. Let $`H_x`$ be the isotropy subgroup of $`\stackrel{~}{x}`$ in $`C(𝐠)`$. Then $`(V_p^𝐠,C(𝐠))`$ induces a germ of uniformizing system at $`x`$ as $`(B_x,H_x)`$ where $`B_x`$ is a small ball in $`V_p^𝐠`$ centered at $`\stackrel{~}{x}`$. Let $`\pi _p(\stackrel{~}{x})=q`$. We need to write $`(B_x,H_x)`$ as $`(V_q^𝐡,C(𝐡))`$ for some $`𝐡G_q\times \mathrm{}\times G_q`$. We let $`\lambda _x:G_qG_p`$ be an induced monomorphism resulted from choosing $`\stackrel{~}{x}`$ as the representative of $`q`$ in $`V_p`$. We define $`𝐡=\lambda _x^1(𝐠)`$ (each $`g_i`$ is in $`\lambda _x(G_q)`$ since $`\stackrel{~}{x}V_p^𝐠`$ and $`\pi _p(\stackrel{~}{x})=q`$.) Then we can identify $`B_x`$ as $`V_q^𝐡`$. We also see that $`H_x=\lambda _x(C(𝐡))`$. Therefore $`(B_x,H_x)`$ is identified as $`(V_q^𝐡,C(𝐡))`$. Hence we proved that $`\stackrel{~}{X}_k`$ is naturally an orbifold with the orbifold structure described above ($`\stackrel{~}{X}_k`$ is Hausdorff and second countable with the given topology for similar reasons). The rest of the lemma is obvious. $`\mathrm{}`$ We can also describe the components of $`\stackrel{~}{X}_k`$ in the same fashion. Using the conjugacy class of monomorphisms $`\pi _{pq}:G_qG_p`$ in the patching condition, we can define an equivalence relation $`(𝐠)_{G_q}(\pi _{pq}(𝐠))_{G_p}`$ similarly. Let $`T_k`$ be the set of equivalence classes. We will write a general element of $`T_k`$ as $`(𝐠)`$. Then $`\stackrel{~}{X}_k`$ is decomposed as a disjoint union of connected orbifolds $$\stackrel{~}{X}_k=\underset{(𝐠)T_k}{}X_{(𝐠)},$$ $`(\mathrm{4.1.2})`$ where $$X_{(𝐠)}=\{(p,(𝐠^{})_{G_p})\|(𝐠^{})_{G_p}(𝐠)\}.$$ $`(\mathrm{4.1.3})`$ There is a map $`o:T_kT`$ induced by the map $`o:(g_1,g_2,\mathrm{},g_k)g_1g_2\mathrm{}g_k`$. We set $`T_k^o=o^1((1))`$. Then $`T_k^oT_k`$ is the subset of equivalence classes $`(𝐠)`$ such that $`𝐠=(g_1,\mathrm{},g_k)`$ satisfies the condition $`g_1\mathrm{}g_k=1`$. Finally, we set $$\stackrel{~}{X}_k^o:=\underset{(𝐠)T_k^o}{}X_{(𝐠)}.$$ $`(\mathrm{4.1.4})`$ In order to define the orbifold cup product, we need a digression on a few classical results about reduced 2-dimensional orbifolds (cf. \[Th\], \[Sc\]). Every closed orbifold of dimension 2 is complex, whose underlying topological space is a closed Riemann surface. More concretely, a closed, reduced 2-dimensional orbifold consists of the following data: a closed Riemann surface $`\mathrm{\Sigma }`$ with complex structure $`j`$, a finite subset of distinct points $`𝐳=(z_1,\mathrm{},z_k)`$ on $`\mathrm{\Sigma }`$, each with a multiplicity $`m_i2`$ (let $`𝐦=(m_1,\mathrm{},m_k)`$), such that the orbifold structure at $`z_i`$ is given by the ramified covering $`zz^{m_i}`$. We will also call a closed, reduced 2-dimensional orbifold a complex orbicurve when the underlying complex analytic structure is emphasized. A $`C^{\mathrm{}}`$ map $`\stackrel{~}{\pi }`$ between two reduced connected 2-dimensional orbifolds is called an orbifold covering if the local liftings of $`\stackrel{~}{\pi }`$ are either a diffeomorphism or a ramified covering. It is shown that the universal orbifold covering exists, and its group of deck transformations is defined to be the orbifold fundamental group of the orbifold. In fact, given a reduced 2-orbifold $`\mathrm{\Sigma }`$, with orbifold fundamental group denoted by $`\pi _1^{orb}(\mathrm{\Sigma })`$, for any subgroup $`\mathrm{\Gamma }`$ of $`\pi _1^{orb}(\mathrm{\Sigma })`$, there is a reduced 2-orbifold $`\stackrel{~}{\mathrm{\Sigma }}`$ and an orbifold covering $`\stackrel{~}{\pi }:\stackrel{~}{\mathrm{\Sigma }}\mathrm{\Sigma }`$ such that $`\stackrel{~}{\pi }`$ induces an injective homomorphism $`\pi _1^{orb}(\stackrel{~}{\mathrm{\Sigma }})\pi _1^{orb}(\mathrm{\Sigma })`$ with image $`\mathrm{\Gamma }\pi _1^{orb}(\mathrm{\Sigma })`$. The orbifold fundamental group of a reduced, closed 2-orbifold $`(\mathrm{\Sigma },𝐳,𝐦)`$ has a presentation $$\pi _1^{orb}(\mathrm{\Sigma })=\{x_i,y_i,\lambda _j,i=1,\mathrm{},g,j=1,\mathrm{},k\|\underset{i}{}x_iy_ix_i^1y_i^1\underset{j}{}\lambda _j=1,\lambda _j^{m_j}=1\},$$ where $`g`$ is the genus of $`\mathrm{\Sigma }`$, $`𝐳=(z_1,\mathrm{},z_k)`$ and $`𝐦=(m_1,\mathrm{},m_k)`$. The remaining ingredient is to construct an “obstruction bundle” $`E_{(𝐠)}`$ over each component $`X_{(𝐠)}`$ where $`(𝐠)T_3^o`$. For this purpose, we consider the Riemann sphere $`S^2`$ with three distinct marked points $`𝐳=(0,1,\mathrm{})`$. Suppose $`(𝐠)`$ is represented by $`𝐠=(g_1,g_2,g_3)`$ and the order of $`g_i`$ is $`m_i`$ for $`i=1,2,3`$. We give a reduced orbifold structure on $`S^2`$ by assigning $`𝐦=(m_1,m_2,m_3)`$ as the multiplicity of $`𝐳`$. The orbifold fundamental group $`\pi _1^{orb}(S^2)`$ has the following presentation $$\pi _1^{orb}(S^2)=\{\lambda _1,\lambda _2,\lambda _3\|\lambda _i^{m_i}=1,\lambda _1\lambda _2\lambda _3=1\},$$ where each generator $`\lambda _i`$ is geometrically represented by a loop around the marked point $`z_i`$ (here recall that $`(z_1,z_2,z_3)=(0,1,\mathrm{})`$). Now for each point $`(p,(𝐠)_{G_p})X_{(𝐠)}`$, fix a representation $`𝐠`$ of $`(𝐠)_{G_p}`$ where $`𝐠=(g_1,g_2,g_3)`$, we define a homomorphism $`\rho _{p,𝐠}:\pi _1^{orb}(S^2)G_p`$ by sending $`\lambda _i`$ to $`g_i`$, which is possible since $`g_1g_2g_3=1`$. Let $`GG_p`$ be the image of $`\rho _{p,𝐠}`$. There is a reduced 2-orbifold $`\mathrm{\Sigma }`$ and an orbifold covering $`\stackrel{~}{\pi }:\mathrm{\Sigma }S^2`$, which induces the following short exact sequence $$1\pi _1(\mathrm{\Sigma })\pi _1^{orb}(S^2)G1.$$ The group $`G`$ acts on $`\mathrm{\Sigma }`$ as the group of deck transformations, whose finiteness implies that $`\mathrm{\Sigma }`$ is closed. Moreover, $`\mathrm{\Sigma }`$ actually has a trivial orbifold structure (i.e. $`\mathrm{\Sigma }`$ is a Riemann surface) since each map $`\lambda _ig_i`$ is injective, and we can assume $`G`$ acts on $`\mathrm{\Sigma }`$ holomorphically. At end, we obtained a uniformizing system $`(\mathrm{\Sigma },G,\stackrel{~}{\pi })`$ of $`(S^2,𝐳,𝐦)`$, which depends on $`(p,𝐠)`$, but is locally constant. The “obstruction bundle” $`E_{(𝐠)}`$ over $`X_{(𝐠)}`$ is constructed as follows. On the local chart $`(V_p^𝐠,C(𝐠))`$ of $`X_{(g)}`$, $`E_{(𝐠)}`$ is given by $`(H^1(\mathrm{\Sigma })TV_p)^G\times V_p^𝐠V_p^𝐠`$, where $`(H^1(\mathrm{\Sigma })TV_p)^G`$ is the invariant subspace of $`G`$. We define an action of $`C(𝐠)`$ on $`H^1(\mathrm{\Sigma })TV_p`$, which is trivial on the first factor and the usual one on $`TV_p`$, then it is clear that $`C(𝐠)`$ commutes with $`G`$, hence $`(H^1(\mathrm{\Sigma })TV_p)^G`$ is invariant under $`C(𝐠)`$. In summary, we have obtained an action of $`C(𝐠)`$ on $`(H^1(\mathrm{\Sigma })TV_p)^G\times V_p^𝐠V_p^𝐠`$, extending the usual one on $`V_p^𝐠`$, and it is easily seen that these trivializations fit together to define the bundle $`E_{(𝐠)}`$ over $`X_{(𝐠)}`$. If we set $`e:X_{(𝐠)}X`$ to be the $`C^{\mathrm{}}`$ map $`(p,(𝐠)_{G_p})p`$, one may think of $`E_{(𝐠)}`$ as $`(H^1(\mathrm{\Sigma })e^{}TX)^G`$. Since we do not assume that $`X`$ is compact, $`X_{(𝐠)}`$ could be a non-compact orbifold in general. The Euler class of $`E_{(𝐠)}`$ depends on a choice of connection on $`E_{(𝐠)}`$. Let $`e_A(E_{(𝐠)})`$ be the Euler form computed from connection $`A`$ by Chern-Weil theory. Definition 4.1.2: For $`\alpha ,\beta H_{orb}^{}(X)`$, and $`\gamma H_{orb,c}^{}(X)`$, we define a 3-point function $$<\alpha ,\beta ,\gamma >_{orb}=\underset{(𝐠)T_3^0}{}_{X_{(𝐠)}}^{orb}e_1^{}\alpha e_2^{}\beta e_3^{}\gamma e_A(E_{(𝐠)}),$$ $`(\mathrm{4.1.5})`$ where each $`e_i:X_{(𝐠)}\stackrel{~}{X}`$ is the $`C^{\mathrm{}}`$ map defined by $`(p,(𝐠)_{G_p})(p,(g_i)_{G_p})`$ for $`𝐠=(g_1,g_2,g_3)`$. Integration over orbifolds is defined by equation $`(2.4)`$. Note that since $`\gamma `$ is compact supported, each integral is finite, and the summation is over a finite subset of $`T_3^o`$. Moreover, if we choose different connection $`A^{}`$, $`e_A(E_{(𝐠)}),e_A^{}(E_{(𝐠)})`$ differ by an exact form. Hence the 3-point function is independent of the choice of the connection $`A`$. Definition 4.1.3: We define the orbifold cup product on $`H_{orb}^{}(X)`$ by the relation $$<\alpha _{orb}\beta ,\gamma >_{orb}=<\alpha ,\beta ,\gamma >_{orb}.$$ $`(\mathrm{4.1.6})`$ Next we shall give a decomposition of the orbifold cup product $`\alpha _{orb}\beta `$ according to the decomposition $`H_{orb}^{}(X)=_{(g)T}H^{2\iota _{(g)}}(X_{(g)})`$, when $`\alpha ,\beta `$ are homogeneous, i.e. $`\alpha H^{}(X_{(g_1)})`$ and $`\beta H^{}(X_{(g_2)})`$ for some $`(g_1),(g_2)T`$. We need to introduce some notation first. Given $`(g_1),(g_2)T`$, let $`T((g_1),(g_2))`$ be the subset of $`T_2`$ which consists of $`(𝐡)`$ where $`𝐡=(h_1,h_2)`$ satisfies $`(h_1)=(g_1)`$ and $`(h_2)=(g_2)`$. Recall that there is map $`o:T_kT`$ defined by sending $`(g_1,g_2,\mathrm{},g_k)`$ to $`g_1g_2\mathrm{}g_k`$. We define a map $`\delta :𝐠(𝐠,o(𝐠)^1)`$, which clearly induces a one to one correspondence between $`T_k`$ and $`T_{k+1}^o`$. We also denote by $`\delta `$ the resulting isomorphism $`\stackrel{~}{X}_k\stackrel{~}{X}_{k+1}^o`$. Finally, we set $`\delta _i=e_i\delta `$. Decomposition Lemma 4.1.4: For any $`\alpha H^{}(X_{(g_1)})`$, $`\beta H^{}(X_{(g_2)})`$, $$\alpha _{orb}\beta =\underset{(𝐡)T((g_1),(g_2))}{}(\alpha _{orb}\beta )_{(𝐡)},$$ $`(\mathrm{4.1.7})`$ where $`(\alpha _{orb}\beta )_{(𝐡)}H^{}(X_{o((𝐡))})`$ is defined by the relation $$<(\alpha _{orb}\beta )_{o((𝐡))},\gamma >_{orb}=_{X_{(𝐡)}}^{orb}\delta _1^{}\alpha \delta _2^{}\beta \delta _3^{}\gamma e_A(\delta ^{}E_{\delta (𝐡)}),$$ $`(\mathrm{4.1.8})`$ for $`\gamma H_c^{}(X_{(o(𝐡)^1)})`$. In the subsequent sections, we shall describe the 3-point function and orbifold cup product in terms of Gromov-Witten invariants. In fact, we will prove the following Theorem 4.1.5: Let $`X`$ be an almost complex orbifold with almost complex structure $`J`$ and $`dim_𝐂X=n`$. The orbifold cup product preserves the orbifold grading, i.e., $$_{orb}:H_{orb}^p(X)\times H_{orb}^q(X)H_{orb}^{p+q}(X)$$ for any $`0p,q2n`$ such that $`p+q2n`$, and has the following properties: 1. The total orbifold cohomology group $`H_{orb}^{}(X)=_{0d2n}H_{orb}^d(X)`$ is a ring with unit $`e_X^0H^0(X)`$ under $`_{orb}`$, where $`e_X^0`$ is the Poincaré dual to the fundamental class $`[X]`$. In particular, $`_{orb}`$ is associative. 2. When $`X`$ is closed, for each $`H_{orb}^d(X)\times H_{orb}^{2nd}(X)H_{orb}^{2n}(X)`$, we have $$_X^{orb}\alpha _{orb}\beta =<\alpha ,\beta >_{orb}.$$ $`(\mathrm{4.1.9})`$ 3. The cup product $`_{orb}`$ is invariant under deformation of $`J`$. 4. When $`X`$ is of integral degree shifting numbers, the total orbifold cohomology group $`H_{orb}^{}(X)`$ is integrally graded, and we have supercommutativity $$\alpha _1_{orb}\alpha _2=(1)^{\mathrm{deg}\alpha _1\mathrm{deg}\alpha _2}\alpha _2_{orb}\alpha _1.$$ 5. Restricted to the nontwisted sectors, i.e., the ordinary cohomologies $`H^{}(X)`$, the cup product $`_{orb}`$ equals the ordinary cup product on $`X`$. When $`X`$ is a complex orbifold, the definition of orbifold cup product $`_{orb}`$ on the total orbifold Dolbeault cohomology group of $`X`$ is completely parallel. We observe that in this case all the objects we have been dealing with are holomorphic, i.e., $`\stackrel{~}{X}_k`$ is a complex orbifold, the “obstruction bundles” $`E_{(𝐠)}X_{(𝐠)}`$ are holomorphic orbifold bundles, and the evaluation maps $`e_i`$ are holomorphic. Definition 4.1.6: For any $`\alpha _1H_{orb}^{p,q}(X)`$, $`\alpha _2H_{orb}^{p^{},q^{}}(X)`$, we define a 3-point function and orbifold cup product in the same fashion as in Definitions 4.1.2, 4.1.3. $`\mathrm{}`$ Note that since the top Chern class of a holomorphic orbifold bundle can be represented by a closed $`(r,r)`$-form where $`r`$ is the (complex) rank of the bundle, it follows that the orbifold cup product preserves the orbifold bi-grading, i.e., $`_{orb}:H_{orb}^{p,q}(X)\times H_{orb}^{p^{},q^{}}(X)H_{orb}^{p+p^{},q+q^{}}(X)`$. The following theorem can be similarly proved. Theorem 4.1.7: Let $`X`$ be a n-dimensional complex orbifold with complex structure $`J`$. The orbifold cup product $$_{orb}:H_{orb}^{p,q}(X)\times H_{orb}^{p^{},q^{}}(X)H_{orb}^{p+p^{},q+q^{}}(X)$$ has the following properties: 1. The total orbifold Dolbeault cohomology group is a ring with unit $`e_X^0H_{orb}^{0,0}(X)`$ under $`_{orb}`$, where $`e_X^0`$ is the class represented by the equaling-one constant function on $`X`$. 2. When $`X`$ is closed, for each $`H_{orb}^{p,q}(X)\times H_{orb}^{np,nq}(X)H_{orb}^{n,n}(X)`$, the integral $`_X\alpha _{orb}\beta `$ equals the Kodaira-Serre pairing $`<\alpha ,\beta >_{orb}`$. 3. The cup product $`_{orb}`$ is invariant under deformation of $`J`$. 4. When $`X`$ is of integral degree shifting numbers, the total orbifold Dolbeault cohomology group of $`X`$ is integrally graded, and we have supercommutativity $$\alpha _1_{orb}\alpha _2=(1)^{\mathrm{deg}\alpha _1\mathrm{deg}\alpha _2}\alpha _2_{orb}\alpha _1.$$ 5. Restricted to the nontwisted sectors, i.e., the ordinary Dolbeault cohomologies $`H^,(X)`$, the cup product $`_{orb}`$ coincides with the ordinary wedge product on $`X`$. 6. When $`X`$ is Kahler and closed, the cup product $`_{orb}`$ coincides with the orbifold cup product on the total orbifold cohomology group $`H_{orb}^{}(X)`$ under the relation $$H_{orb}^r(X)𝐂=_{p+q=r}H_{orb}^{p,q}(X).$$ ### 4.2 Moduli space of ghost maps We first give a classification of rank-n complex orbifold bundles over a closed, reduced, 2-dimensional orbifold. Let $`(\mathrm{\Sigma },𝐳,𝐦)`$ be a closed, reduced, 2-dimensional orbifold, where $`𝐳=(z_1,\mathrm{},z_k)`$ and $`𝐦=(m_1,\mathrm{},m_k)`$. Let $`E`$ be a complex orbifold bundle of rank $`n`$ over $`(\mathrm{\Sigma },𝐳,𝐦)`$. Then at each singular point $`z_i`$, $`i=1,\mathrm{},k`$, $`E`$ determines a representation $`\rho _i:𝐙_{m_i}Aut(𝐂^n)`$ so that over a disc neighborhood of $`z_i`$, $`E`$ is uniformized by $`(D\times 𝐂^n,𝐙_{m_i},\pi )`$ where the action of $`𝐙_{m_i}`$ on $`D\times 𝐂^n`$ is given by $$e^{2\pi i/m_i}(z,w)=(e^{2\pi i/m_i}z,\rho _i(e^{2\pi i/m_i})w)$$ $`(\mathrm{4.2.1})`$ for any $`w𝐂^n`$. Each representation $`\rho _i`$ is uniquely determined by a $`n`$-tuple of integers $`(m_{i,1},\mathrm{},m_{i,n})`$ with $`0m_{i,j}<m_i`$, as it is given by matrix $$\rho _i(e^{2\pi i/m_i})=diag(e^{2\pi im_{i,1}/m_i},\mathrm{},e^{2\pi im_{i,n}/m_i}).$$ $`(\mathrm{4.2.2})`$ Over the punctured disc $`D_i\{0\}`$ at $`z_i`$, $`E`$ inherits a specific trivialization from $`(D\times 𝐂^n,𝐙_{m_i},\pi )`$ as follows: We define a $`𝐙_{m_i}`$-equivariant map $`\mathrm{\Psi }_i:D\{0\}\times 𝐂^nD\{0\}\times 𝐂^n`$ by $$(z,w_1,\mathrm{},w_n)(z^{m_i},z^{m_{i,1}}w_1,\mathrm{},z^{m_{i,n}}w_n),$$ $`(\mathrm{4.2.3})`$ where $`Z_{m_i}`$ acts trivially on the second $`D\{0\}\times 𝐂^n`$. Hence $`\mathrm{\Psi }_i`$ induces a trivialization $`\psi _i:E_{D_i\{0\}}D_i\{0\}\times 𝐂^n`$. We can extend the smooth complex vector bundle $`E_{\mathrm{\Sigma }𝐳}`$ over $`\mathrm{\Sigma }𝐳`$ to a smooth complex vector bundle over $`\mathrm{\Sigma }`$ by using these trivializations $`\psi _i`$. We call the resulting complex vector bundle the de-singularization of $`E`$, and denote it by $`\|E\|`$. Proposition 4.2.1: The space of isomorphism classes of complex orbifold bundles of rank $`n`$ over a closed, reduced, 2-dimensional orbifold $`(\mathrm{\Sigma },𝐳,𝐦)`$ where $`𝐳=(z_1,\mathrm{},z_k)`$ and $`𝐦=(m_1,\mathrm{},m_k)`$, is in 1:1 correspondence with the set of $`(c,(m_{1,1},\mathrm{},m_{1,n}),\mathrm{},(m_{k,1},\mathrm{},m_{k,n}))`$ for $`c𝐐`$, $`m_{i,j}𝐙`$, where $`c`$ and $`m_{i,j}`$ are confined by the following condition: $$0m_{i,j}<m_i\text{ and }c\underset{i=1}{\overset{k}{}}\underset{j=1}{\overset{n}{}}\frac{m_{i,j}}{m_i}(mod𝐙).$$ $`(\mathrm{4.2.4})`$ In fact, $`c`$ is the first Chern number of the orbifold bundle and $`c(_{i=1}^k_{j=1}^n\frac{m_{i,j}}{m_i})`$ is the first Chern number of its de-singularization. Proof: We only need to show the relation: $$c_1(E)([\mathrm{\Sigma }])=c_1(\|E\|)([\mathrm{\Sigma }])+\underset{i=1}{\overset{k}{}}\underset{j=1}{\overset{n}{}}\frac{m_{i,j}}{m_i}.$$ $`(\mathrm{4.2.5})`$ We take a connection $`_0`$ on $`\|E\|`$ which equals $`d`$ on a disc neighborhood $`D_i`$ of each $`z_i𝐳`$ so that $`c_1(\|E\|)([\mathrm{\Sigma }])=_\mathrm{\Sigma }c_1(_0)`$. We use $`_0^{}`$ to denote the pull-back connection $`br_i^{}_0`$ on $`D\{0\}\times 𝐂^n`$ via $`br_i:DD_i`$ by $`zz^{m_i}`$. On the other hand, on each uniformizing system $`(D\times 𝐂^n,𝐙_{m_i},\pi )`$, we take the trivial connection $`_i=d`$ which is obvious $`𝐙_{m_i}`$-equivariant. Furthermore, we take a $`𝐙_{m_i}`$-equivariant cut-off function $`\beta _i`$ on $`D`$ which equals one in a neighborhood of the boundary $`D`$. We are going to paste these connections together to get a connection $``$ on $`E`$. We define $``$ on each uniformizing system $`(D\times 𝐂^n,𝐙_{m_i},\pi )`$ by $$_vu=(1\beta _i)(_i)_vu+\beta _i\overline{\psi }_i^1(_0)_{\overline{\psi }_iv}\overline{\psi }_iu,$$ $`(\mathrm{4.2.6})`$ where $`\overline{\psi }_i:D\{0\}\times 𝐂^nD\{0\}\times 𝐂^n`$ is given by $$(z,w_1,\mathrm{},w_n)(z,z^{m_{i,1}}w_1,\mathrm{},z^{m_{i,n}}w_n).$$ $`(\mathrm{4.2.7})`$ One easily verifies that $`F()=F(_0)`$ on $`\mathrm{\Sigma }(_iD_i)`$ and $$F()=diag(d(\beta _im_{i,1}dz/z),\mathrm{},d(\beta _im_{i,n}dz/z))$$ on each uniformizing system $`(D,𝐙_{m_i},\pi )`$. So $`c_1(E)([\mathrm{\Sigma }])`$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }^{orb}}c_1()`$ $`=`$ $`{\displaystyle _{\mathrm{\Sigma }(_iD_i)}}c_1(_0)+{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{1}{m_i}}{\displaystyle _D}c_1()`$ $`=`$ $`c_1(\|E\|)([\mathrm{\Sigma }])+{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{m_{i,j}}{m_i}}.`$ Here the integraton over $`\mathrm{\Sigma }`$, $`_\mathrm{\Sigma }^{orb}`$, should be understood as in $`(2.4)`$. $`\mathrm{}`$ We will need the following index formula. Proposition 4.2.2: Let $`E`$ be a holomorphic orbifold bundle of rank $`n`$ over a complex orbicurve $`(\mathrm{\Sigma },𝐳,𝐦)`$ of genus $`g`$. Then $`𝒪(E)=𝒪(\|E\|)`$, where $`𝒪(E),𝒪(\|E\|)`$ are sheaves of holomorphic sections of $`E,\|E\|`$. Hence, $$\chi (𝒪(E))=\chi (𝒪(\|E\|))=c_1(\|E\|)([\mathrm{\Sigma }])+n(1g).$$ $`(\mathrm{4.2.9})`$ If $`E`$ corresponds to $`(c,(m_{1,1},\mathrm{},m_{1,n}),\mathrm{},(m_{k,1},\mathrm{},m_{k,n}))`$ (cf. Proposition 4.2.1), then we have $$\chi (𝒪(E))=n(1g)+c_1(E)([\mathrm{\Sigma }])\underset{i=1}{\overset{k}{}}\underset{j=1}{\overset{n}{}}\frac{m_{i,j}}{m_i}.$$ Proof: By construction, we have $`𝒪(E)=𝒪(\|E\|)`$. Hence $$\chi (𝒪(E))=\chi (𝒪(\|E\|))=c_1(\|E\|)([\mathrm{\Sigma }])+n(1g).$$ $`(\mathrm{4.2.10})`$ By proposition 4.2.1, we have $$\chi (𝒪(E))=n(1g)+c_1(E)([\mathrm{\Sigma }])\underset{i=1}{\overset{k}{}}\underset{j=1}{\overset{n}{}}\frac{m_{i,j}}{m_i},$$ if $`E`$ corresponds to $`(c,(m_{1,1},\mathrm{},m_{1,n}),\mathrm{},(m_{k,1},\mathrm{},m_{k,n}))`$. $`\mathrm{}`$ Now we come to the main issue of this section. Recall that suppose $`f:XX^{}`$ is a $`C^{\mathrm{}}`$ map between manifolds and $`E`$ is a smooth vector bundle over $`X^{}`$, then there is a smooth pull-back vector bundle $`f^{}E`$ over $`X`$ and a bundle morphism $`\overline{f}:f^{}EE`$ which covers the map $`f`$. However, if instead, we have a $`C^{\mathrm{}}`$ map $`\stackrel{~}{f}`$ between orbifolds $`X`$ and $`X^{}`$, and an orbifold bundle $`E`$ over orbifold $`X^{}`$, the question whether there is a pull-back orbifold bundle $`E^{}`$ over $`X^{}`$ and an orbifold bundle morphism $`\overline{f}:E^{}E`$ covering the map $`\stackrel{~}{f}`$ is a quite complicated issue: (1) What is the precise meaning of pull-back orbifold bundle $`E^{}`$, (2) $`E^{}`$ might not exist, or even if it exists, it might not be unique. Understanding this question is the first step in our establishment of an orbifold Gromov-Witten theory in \[CR\]. In the present case, given a constant map $`f:\mathrm{\Sigma }X`$ from a marked Riemann surface $`\mathrm{\Sigma }`$ with marked-point set $`𝐳`$ into an almost complex orbifold $`X`$, we need to settle the existence and classification problem of pull-back orbifold bundles via $`f`$, with some reduced orbifold structure on $`\mathrm{\Sigma }`$, whose set of orbifold points is contained in the given marked-point set $`𝐳`$. Let $`(S^2,𝐳)`$ be a genus-zero Riemann surface with k-marked points $`𝐳=(z_1,\mathrm{},z_k)`$, $`pX`$ any point in an almost complex orbifold $`X`$ with $`dim_𝐂X=n`$, and $`(V_p,G_p,\pi _p)`$ a local chart at $`p`$. Then for any k-tuple $`𝐠=(g_1,\mathrm{},g_k)`$ where $`g_iG_p`$, $`i=1,\mathrm{},k`$, there is an orbifold structure on $`S^2`$ so that it becomes a complex orbicurve $`(S^2,𝐳,𝐦)`$ where $`𝐦=(\|g_1\|,\mathrm{},\|g_k\|)`$ (here $`\|g\|`$ stands for the order of $`g`$). If further assuming that $`o(𝐠)=g_1g_2\mathrm{}g_k=1_{G_p}`$, one can construct a rank-n holomorphic orbifold bundle $`E_{p,𝐠}`$ over $`(S^2,𝐳,𝐦)`$, together with an orbifold bundle morphism $`\mathrm{\Phi }_{p,𝐠}:E_{p,𝐠}TX`$ covering the constant map from $`S^2`$ to $`pX`$, as we shall see next. Denote $`\mathrm{𝟏}_{G_p}=(1_{G_p},\mathrm{},1_{G_p})`$. The case $`𝐠=\mathrm{𝟏}_{G_p}`$ is trivial; we simply take the rank-n trivial holomorphic bundle over $`S^2`$. Hence in what follows, we assume that $`𝐠\mathrm{𝟏}_{G_p}`$. We recall that the orbifold fundamental group of $`(S^2,𝐳,𝐦)`$ is given by $$\pi _1^{orb}(S^2)=\{\lambda _1,\lambda _2,\mathrm{},\lambda _k\|\lambda _i^{\|g_i\|}=1,\lambda _1\lambda _2\mathrm{}\lambda _k=1\},$$ where each generator $`\lambda _i`$ is geometrically represented by a loop around the marked point $`z_i`$. We define a homomorphism $`\rho :\pi _1^{orb}(S^2)G_p`$ by sending each $`\lambda _i`$ to $`g_iG_p`$ (note that we assumed that $`g_1g_2\mathrm{}g_k=1_{G_p}`$). There is a closed Riemann surface $`\mathrm{\Sigma }`$ and a finite group $`G`$ acting on $`\mathrm{\Sigma }`$ holomorphically, such that $`(\mathrm{\Sigma },G)`$ uniformizes $`(S^2,𝐳,𝐦)`$ and $`\pi _1(\mathrm{\Sigma })=\mathrm{ker}\rho `$ with $`G=Im\rho G_p`$. We identify $`(TV_p)_p`$ with $`𝐂^n`$ and denote the rank-n trivial holomorphic vector bundle on $`\mathrm{\Sigma }`$ by $`\underset{¯}{𝐂^n}`$. The representation $`GAut((TV_p)_p)`$ defines a holomorphic action on the holomorphic vector bundle $`\underset{¯}{𝐂}^n`$. We take $`E_{p,𝐠}`$ to be the corresponding holomorphic orbifold bundle uniformized by $`(\underset{¯}{𝐂}^n,G,\stackrel{~}{\pi })`$ where $`\stackrel{~}{\pi }:\underset{¯}{𝐂}^n\underset{¯}{𝐂}^n/G`$ is the quotient map. There is a natural orbifold bundle morphism $`\mathrm{\Phi }_{p,𝐠}:E_{p,𝐠}TX`$ sending $`\mathrm{\Sigma }`$ to the point $`p`$. By the nature of construction, if $`𝐠=(g_1,\mathrm{},g_k)`$ and $`𝐠^{}=(g_1^{},\mathrm{},g_k^{})`$ are conjugate, i.e., there is an element $`gG_p`$ such that $`g_i^{}=g^1g_ig`$, then there is an isomorphism $`\psi :E_{p,𝐠}E_{p,𝐠^{}}`$ such that $`\mathrm{\Phi }_{p,𝐠}=\mathrm{\Phi }_{p,𝐠^{}}\psi `$. If there is an isomorphism $`\psi :E_{p,𝐠}E_{p,𝐠^{}}`$ such that $`\mathrm{\Phi }_{p,𝐠}=\mathrm{\Phi }_{p,𝐠^{}}\psi `$, then there is a lifting $`\stackrel{~}{\psi }:\stackrel{~}{E}_{p,𝐠}\stackrel{~}{E}_{p,𝐠^{}}`$ of $`\psi `$ and an automorphism $`\varphi :TV_pTV_p`$, such that $`\varphi \stackrel{~}{\mathrm{\Phi }}_{p,𝐠}=\stackrel{~}{\mathrm{\Phi }}_{p,𝐠^{}}\stackrel{~}{\psi }`$. If $`\varphi `$ is given by the action of an element $`gG_p`$, then we have $`gg_ig^1=g_i^{}`$ for all $`i=1,\mathrm{},k`$. Lemma 4.2.3: Let $`E`$ be a rank-n holomorphic orbifold bundle over $`(S^2,𝐳,𝐦)`$ (for some $`𝐦`$). Suppose that there is an orbifold bundle morphism $`\mathrm{\Phi }:ETX`$ covering a constant map from $`S^2`$ into $`X`$. Then there is a $`(p,𝐠)`$ such that $`(E,\mathrm{\Phi })=(E_{p,𝐠},\mathrm{\Phi }_{p,𝐠})`$. Proof: Let $`E`$ be a rank-n holomorphic orbifold bundle over $`(S^2,𝐳,𝐦)`$ with a morphism $`\mathrm{\Phi }:ETX`$ covering the constant map to a point $`p`$ in X. We will find a $`𝐠`$ so that $`(E,\mathrm{\Phi })=(E_{p,𝐠},\mathrm{\Phi }_{p,𝐠})`$. For this purpose, we again consider the uniformizing system $`(\mathrm{\Sigma },G,\pi )`$ of $`(S^2,𝐳,𝐦)`$ where $`\mathrm{\Sigma }`$ is a closed Riemann surface with a holomorphic action by a finite group $`G`$. Then there is a holomorphic vector bundle $`\stackrel{~}{E}`$ over $`\mathrm{\Sigma }`$ with a compatible action of $`G`$ so that $`(\stackrel{~}{E},G)`$ uniformizes the holomorphic orbifold bundle $`E`$. Moreover, there is a vector bundle morphism $`\stackrel{~}{\mathrm{\Phi }}:\stackrel{~}{E}TV_p`$, which is a lifting of $`\mathrm{\Phi }`$ so that for any $`aG`$, there is a $`\stackrel{~}{\lambda }(a)`$ in $`G_p`$ such that $`\stackrel{~}{\mathrm{\Phi }}a=\stackrel{~}{\lambda }(a)\mathrm{\Phi }`$. In fact, $`a\stackrel{~}{\lambda }(a)`$ defines a homomorphism $`\stackrel{~}{\lambda }:GG_p`$. Since $`\stackrel{~}{\mathrm{\Phi }}`$ covers a constant map from $`\mathrm{\Sigma }`$ into $`V_p`$, the holomorphic vector bundle $`\stackrel{~}{E}`$ is in fact a trivial bundle. Recall that $`G`$ is the quotient group of $`\pi _1^{orb}(S^2)`$ by the normal subgroup $`\pi _1(\mathrm{\Sigma })`$. Let $`\lambda `$ be the induced homomorphism $`\pi _1^{orb}(S^2)G_p`$, and let $`g_i=\lambda (\gamma _i)`$. Then we have $`g_1g_2\mathrm{}g_k=1_{G_p}`$. We simply define $`𝐠=(g_1,g_2,\mathrm{},g_k)`$. It is easily seen that $`(E,\mathrm{\Phi })=(E_{p,𝐠},\mathrm{\Phi }_{p,𝐠})`$. $`\mathrm{}`$ Definition 4.2.4: Given a genus-zero Riemann surface with k-marked points $`(\mathrm{\Sigma },𝐳)`$, where $`𝐳=(z_1,\mathrm{},z_k)`$, we call each equivalence class $`[E_{p,𝐠},\mathrm{\Phi }_{p,𝐠}]`$ of pair $`(E_{p,𝐠},\mathrm{\Phi }_{p,𝐠})`$ a ghost map from $`(\mathrm{\Sigma },𝐳)`$ into $`X`$. A ghost map $`[E,\mathrm{\Phi }]`$ from $`(\mathrm{\Sigma },𝐳)`$ is said to be equivalent to a ghost map $`[E^{},\mathrm{\Phi }^{}]`$ from $`(\mathrm{\Sigma }^{},𝐳^{})`$ ($`𝐳^{}=(z_1^{},\mathrm{},z_k^{})`$) if there is a holomorphic orbifold bundle morphism $`\stackrel{~}{\psi }:EE^{}`$ covering a biholomorphism $`\psi :\mathrm{\Sigma }\mathrm{\Sigma }^{}`$ such that $`\psi (z_i)=z_i^{}`$ and $`\mathrm{\Phi }=\mathrm{\Phi }^{}\stackrel{~}{\psi }`$. An equivalence class of ghost maps is called a ghost curve (with k-marked points). We denote by $`_k`$ the moduli space of ghost curves with k-marked points. $`\mathrm{}`$ As a consequence, we obtain Proposition 4.2.5: Let $`X`$ be an almost complex orbifold. For any $`k0`$, the moduli space of ghost curves with $`k`$-marked points $`_k`$ is naturally an almost complex orbifold. When $`k4`$, $`_k`$ can be identified with $`_{0,k}\times \stackrel{~}{X}_k^o`$, where $`_{0,k}`$ is the moduli space of genus-zero curve with $`k`$-marked points. It has a natural partial compatification $`\overline{}_k`$, which is an almost complex orbifold and can be identified with $`\overline{}_{0,k}\times \stackrel{~}{X}_k^o`$, where $`\overline{}_{0,k}`$ is the Deligne-Mumford compatification of $`_{0,k}`$. Remarks 4.2.6: The natural partial compatification $`\overline{}_k`$ of $`_k`$ ($`k4`$) can be interpreted geometrically as adding nodal ghost curves into $`_k`$. The space $`\stackrel{~}{X}_2^o`$ is naturally identified with the graph of the map $`I:\stackrel{~}{\mathrm{\Sigma }X}\stackrel{~}{\mathrm{\Sigma }X}`$ in $`\stackrel{~}{\mathrm{\Sigma }X}\times \stackrel{~}{\mathrm{\Sigma }X}`$, where $`I`$ is defined by $`(p,(g))(p,(g^1))`$. Next, we construct a complex orbifold bundle $`E_k`$, a kind of obstruction bundle in nature, over the moduli space $`_k`$ of ghost curves with k-marked points. The rank of $`E_k`$ may vary over different connected components of $`_k`$. When $`k=3`$, the restriction of $`E_3`$ to each component gives a geometric construction of obstruction bundle $`E_{(𝐠)}`$ in the last section under identification $`_3=\stackrel{~}{X}_3^o`$. Let us consider the space $`𝒞_k`$ of all triples $`((\mathrm{\Sigma },𝐳),E_{p,𝐠},\mathrm{\Phi }_{p,𝐠})`$ where $`(\mathrm{\Sigma },𝐳)`$ is a genus-zero curve with k-marked points $`𝐳=(z_1,\mathrm{},z_k)`$, $`E_{p,𝐠}`$ is a rank-n holomorphic orbifold bundle over $`\mathrm{\Sigma }`$, and $`\mathrm{\Phi }_{p,𝐠}:E_{p,𝐠}TX`$ a morphism covering the constant map sending $`\mathrm{\Sigma }`$ to the point $`p`$ in $`X`$. To each point $`x𝒞_k`$ we assign a complex vector space $`V_x`$, which is the cokernel of the operator $$\overline{}:\mathrm{\Omega }^{0,0}(E_{p,𝐠})\mathrm{\Omega }^{0,1}(E_{p,𝐠}).$$ $`(\mathrm{4.2.11})`$ We introduce an equivalence relation $``$ amongst pairs $`(x,v)`$ where $`x𝒞_k`$ and $`vV_x`$ as follows: Let $`x=((\mathrm{\Sigma },𝐳),E_{p,𝐠},\mathrm{\Phi }_{p,𝐠})`$ and $`x^{}=((\mathrm{\Sigma }^{},𝐳^{}),E_{p^{},𝐠^{}},\mathrm{\Phi }_{p^{},𝐠^{}})`$, then $`(x,v)(x^{},v^{})`$ if there is a morphism $`\stackrel{~}{\psi }:E_{p,𝐠}E_{p^{},𝐠^{}}`$ such that $`\mathrm{\Phi }_{p,𝐠}=\mathrm{\Phi }_{p^{},𝐠^{}}\stackrel{~}{\psi }`$ and $`\stackrel{~}{\psi }`$ covers a biholomorphism $`\psi :\mathrm{\Sigma }\mathrm{\Sigma }^{}`$ satisfying $`\psi (𝐳)=𝐳^{}`$ (as ordered sets), and $`v^{}=\psi _{}(v)`$ where $`\psi _{}:V_xV_x^{}`$ is induced by $`\stackrel{~}{\psi }`$. We define $`E_k`$ to be the quotient space of all $`(x,v)`$ under this equivalence relation. There is obviously a surjective map $`pr:E_k_k`$ induced by the projection $`(x,v)x`$. Lemma 4.2.7: The space $`E_k`$ can be given a topology such that $`pr:E_k_k`$ is a complex orbifold bundle over $`_k`$. Proof: First we show that the dimension of $`V_x`$ is a local constant function of the equivalence class $`[x]`$ in $`_k`$. Recall a neighborhood of $`[x]`$ in $`_k`$ is given by $`𝒪\times V_p^𝐠/C(𝐠)`$ where $`𝒪`$ is a neighborhood of the genus-zero curve with k-marked points $`(\mathrm{\Sigma },𝐳)`$ in the moduli space $`_{0,k}`$. In fact, we will show that the kernel of $`(\mathrm{4.2.11})`$ is identified with $`(TV_p^𝐠)_p`$, whose dimension is a local constant. Then it follows that $`dimV_x`$ is locally constant as the dimension of cokernel of $`(\mathrm{4.2.11})`$, since by Proposition 4.2.2, the index of $`(\mathrm{4.2.11})`$ is locally constant. For the identification of the kernel of $`(\mathrm{4.2.11})`$, recall that the holomorphic orbifold bundle $`E_{p,𝐠}`$ over the genus-zero curve $`\mathrm{\Sigma }`$ is uniformized by the trivial holomorphic vector bundle $`\underset{¯}{𝐂}^n`$ over a Riemann surface $`\stackrel{~}{\mathrm{\Sigma }}`$ with a holomorphic action of a finite group $`G`$. Hence the kernel of $`(\mathrm{4.2.11})`$ is identified with the $`G`$-invariant holomorphic sections of the trivial bundle $`\underset{¯}{𝐂}^n`$, which are constant sections invariant under $`G`$. Through morphism $`\mathrm{\Phi }_{p,𝐠}:E_{p,𝐠}TX`$, the kernel of $`\overline{}`$ is then identified with $`(TV_p^𝐠)_p`$. Recall that the moduli space $`_{0,k}`$ is a smooth complex manifold. Let $`𝒪`$ be a neighborhood of $`(\mathrm{\Sigma }_0,𝐳_0)`$ in $`_{0,k}`$. Then a neighborhood of $`[x_0]=[(\mathrm{\Sigma }_0,𝐳_0),E_{p,𝐠},\mathrm{\Phi }_{p,𝐠}]`$ in $`_k`$ is uniformized by $`(𝒪\times V_p^𝐠,C(𝐠))`$ (cf. Lemma 4.1.1). More precisely, to any $`((\mathrm{\Sigma },𝐳),y)𝒪\times V_p^𝐠`$, we associate a rank-n holomorphic orbifold bundle over $`(\mathrm{\Sigma },𝐳)`$ as follows: Let $`q=\pi _p(y)U_p`$, then the pair $`(y,𝐠)`$ canonically determines a $`𝐡_yG_q\times \mathrm{}\times G_q`$, and there is a canonically constructed holomorphic orbifold bundle $`E_{q,𝐡_y}`$ over $`(\mathrm{\Sigma },𝐳)`$ with morphism $`\mathrm{\Phi }_{q,𝐡_y}:E_{q,𝐡_y}TX`$ covering the constant map to $`q`$. Hence we have a family of holomorphic orbibundles over genus-zero curve with k-marked points, which are parametrized by $`𝒪\times V_p^𝐠`$. Moreover, it depends on the parameter in $`𝒪`$ holomorphically and the action of $`C(𝐠)`$ on $`V_p^𝐠`$ coincides with the equivalence relation between the pairs of holomorphic orbifold bundle and morphism $`(E_{q,𝐡_y},\mathrm{\Phi }_{q,𝐡_y})`$. Now we put a Kahler metric on each genus-zero curve in $`𝒪`$ which is compatible to the complex structure and depends smoothly on the parameter in $`𝒪`$, and we also put a hermitian metric on $`X`$. Then we have a family of first order elliptic operators depending smoothly on the parameters in $`𝒪\times V_p^𝐠`$: $$\overline{}^{}:\mathrm{\Omega }^{0,1}(E_{q,𝐡_y})\mathrm{\Omega }^{0,0}(E_{q,𝐡_y})$$ and whose kernel gives rise to a complex vector bundle $`E_{x_0}`$ over $`𝒪\times V_p^𝐠`$. The finite group $`C(𝐠)`$ naturally acts on the complex vector bundle which coincides with the equivalence relation amongst the pairs $`(x,v)`$ where $`x𝒞_k`$ and $`vV_x`$. Hence $`(E_{x_0},C(𝐠))`$ is a uniformizing system for $`pr^1(𝒪\times V_p^𝐠/C(𝐠))`$, which fits together to give an orbifold bundle structure for $`pr:E_k_k`$. $`\mathrm{}`$ Remark 4.2.8: Recall that each holomorphic orbifold bundle $`E_{p,𝐠}`$ over $`(S^2,𝐳,𝐦)`$ can be uniformized by a trivial holomorphic vector bundle $`\underset{¯}{𝐂}^n`$ over a Riemann surface $`\mathrm{\Sigma }`$ with a holomorphic group action by $`G`$. Hence each element $`\xi `$ in the kernel of $$\overline{}^{}:\mathrm{\Omega }^{0,1}(E_{p,𝐠})\mathrm{\Omega }^{0,0}(E_{p,𝐠})$$ can be identified with a $`G`$-invariant harmonic $`(0,1)`$-form on $`\mathrm{\Sigma }`$ with value in $`(TV_p)_p`$ (here we identify each fiber of $`\underset{¯}{𝐂}^n`$ with $`(TV_p)_p`$ through the morphism $`\mathrm{\Phi }_{p,𝐠}`$), i.e., $`\xi =w\alpha `$ where $`w(TV_p)_p`$, $`\alpha `$ is a harmonic $`(0,1)`$-form on $`\stackrel{~}{\mathrm{\Sigma }}`$, and $`\xi `$ is $`G`$-invariant. Therefore, when $`k=3`$, it agrees with $`E_{(𝐠)}`$. We observe that with respect to the taken hermitian metric on $`X`$, $`w(TV_p)_p`$ must lie in the orthogonal complement of $`(TV_p^𝐠)_p`$ in $`(TV_p)_p`$. This is because: For any $`u(TV_p^𝐠)_p`$ and a harmonic $`(0,1)`$-form $`\beta `$ on $`\mathrm{\Sigma }`$, if $`u\beta `$ is $`G`$-invariant, then $`\beta `$ is $`G`$-invariant too, which means that $`\beta `$ descents to a harmonic $`(0,1)`$-form on $`S^2`$, and $`\beta `$ must be identically zero. $`\mathrm{}`$ Recall the cup product is defined by equation $$<\alpha _1_{orb}\alpha _2,\gamma >_{orb}=(__3^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(\gamma )e(E_3)),$$ where $`e(E_3)`$ is the Euler form of the complex orbifold bundle $`E_3`$ over $`_3`$ and $`\gamma H_{orb,c}^{}(X)`$. We take a basis $`\{e_j\},\{e_k^o\}`$ of the total orbifold cohomology group $`H_{orb}^{}(X),H_{orb,c}^{}(X)`$ such that each $`e_j,e_k^o`$ is of homogeneous degree. Let $`<e_j,e_k^o>_{orb}=a_{jk}`$ be the Poincare pairing matrix and $`(a^{jk})`$ be the inverse. It is easy to check that the Poincare dual of graph of $`I`$ in $`\stackrel{~}{\mathrm{\Sigma }}^2`$ can be written as $`_{j,k}a^{jk}e_je_k^o`$. Then, $$\alpha _1_{orb}\alpha _2=\underset{j,k}{}e_ja^{kj}(__3^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(e_k^o)e(E_3)).$$ $`(\mathrm{4.2.12})`$ Proof of Theorem 4.1.5: We postpone the proof of associativity of $`_{orb}`$ to the next subsection. We first show that if $`\alpha _1H_{orb}^p(X)`$ and $`\alpha _2H_{orb}^q(X)`$, then $`\alpha _1_{orb}\alpha _2`$ is in $`H_{orb}^{p+q}(X)`$. For the integral in (4.2.12) to be nonzero, $$\mathrm{deg}(e_1^{}(\alpha _1))+\mathrm{deg}(e_2^{}(\alpha _2))+\mathrm{deg}(e_3^{}(e_k^o))+\mathrm{deg}(e(E_3))=2dim_𝐂_3.$$ $`(\mathrm{4.2.13})`$ Here $`\mathrm{deg}`$ stands for the degree of a cohomology class without degree shifting. The degree of Euler class $`e(E_3)`$ is equal to the dimension of cokernel of $`(\mathrm{4.2.11})`$, which by index formula (cf. Proposition 4.2.2) equals $`2dim_𝐂_3^{(i)}(2n2_{j=1}^3\iota (p,g_j))`$ on a connected component $`_3^{(i)}`$ containing point $`(p,(𝐠))`$ where $`𝐠=(g_1,g_2,g_3)`$. Hence $`(\mathrm{4.2.13})`$ becomes $$\mathrm{deg}(\alpha _1)+deg(\alpha _2)+deg(e_k^o)+2\underset{j=1}{\overset{3}{}}\iota (p,g_j)=2n,$$ $`(\mathrm{4.2.14})`$ from which it is easily seen that $`\alpha _1_{orb}\alpha _2`$ is in $`H_{orb}^{p+q}(X)`$. Next we show that $`e_X^0`$ is a unit with respect to $`_{orb}`$, i.e., $`\alpha _{orb}e_X^0=e_X^0_{orb}\alpha =\alpha `$. First observe that there are connected components of $`_3`$ consisting of points $`(p,(𝐠))`$ for which $`𝐠=(g_1,g_2,g_3)`$ satisfies the condition that one of the $`g_i`$ is $`1_{G_p}`$. Over these components the Euler class $`e(E_3)=1`$ in the $`0^{th}`$ cohomology group since $`(\mathrm{4.2.11})`$ has zero cokernel. Let $`\alpha H^{}(X_{(g)})`$. Then $`e_1^{}(\alpha )e_2^{}(e_X^0)e_3^{}(e_k^o)`$ is non-zero only on the connected component of $`_3`$ which is the image of the embedding $`X_{(g)}_3`$ given by $`(p,(g)_{G_p})(p,((g,1_{G_p},g^1)))`$ and $`e_k^o`$ must be in $`H_c^{}(X_{(g^1)})`$. Moreover, we have $`\alpha _{orb}e_X^0`$ $`:=`$ $`{\displaystyle \underset{j,k}{}}({\displaystyle __3^{orb}}e_1^{}(\alpha )e_2^{}(e_X^0)e_3^{}(e_k^o)e(E_3))a^{kj}e_j`$ $`=`$ $`{\displaystyle \underset{j,k}{}}({\displaystyle _{X_{(g)}}^{orb}}\alpha I^{}(e_k^o))a^{kj}e_j`$ $`=`$ $`\alpha `$ Similarly, we can prove that $`e_X^0_{orb}\alpha =\alpha `$. Now we consider the case $`_{orb}:H_{orb}^d(X)\times H_{orb}^{2nd}(X)H_{orb}^{2n}(X)=H^{2n}(X)`$. Let $`\alpha H_{orb}^d(X)`$ and $`\beta H_{orb}^{2nd}(X)`$, then $`e_1^{}(\alpha )e_2^{}(\beta )e_3^{}(e_X^0)`$ is non-zero only on those connected components of $`_3`$ which are images under embedding $`\stackrel{~}{X}_3`$ given by $`(p,(g))(p,((g,g^1,1_{G_p})))`$, and if $`\alpha `$ is in $`H^{}(X_{(g)})`$, $`\beta `$ must be in $`H^{}(X_{(g^1)})`$. Moreover, let $`e_X^{2n}`$ be the generator in $`H^{2n}(X)`$ such that $`e_X^{2n}[X]=1`$, then we have $`\alpha _{orb}\beta `$ $`:=`$ $`{\displaystyle \underset{j,k}{}}({\displaystyle __3^{orb}}e_1^{}(\alpha )e_2^{}(\beta )e_3^{}(e_k^o)e(E_3))a^{kj}e_j`$ $`=`$ $`({\displaystyle __3^{orb}}e_1^{}(\alpha )e_2^{}(\beta )e_3^{}(e_X^0)e(E_3))e_X^{2n}`$ $`=`$ $`({\displaystyle _{\stackrel{~}{X}}^{orb}}\alpha I^{}(\beta ))e_X^{2n}`$ $`=`$ $`<\alpha ,\beta >_{orb}e_X^{2n}`$ from which we see that $`_X\alpha _{orb}\beta =<\alpha ,\beta >_{orb}`$. The rest of the assertions are obvious. $`\mathrm{}`$ ### 4.3 Proof of associativity In this subsection, we give a proof of associativity of the orbifold cup products $`_{orb}`$ defined in the last subsection. We will only present the proof for the orbifold cohomology groups $`H_{orb}^{}(X)`$. The proof for orbifold Dolbeault cohomology is the same. We leave it to readers. Recall the moduli space of ghost curves with k-marked points $`_k`$ for $`k4`$ can be identified with $`_{0,k}\times \stackrel{~}{X}_k^o`$ which admits a natural partial compatification $`\overline{}_{0,k}\times \stackrel{~}{X}_k^o`$ by adding nodal ghost curves. We will first give a detailed analysis on this for the case when $`k=4`$. Let $`\mathrm{\Delta }`$ be the graph of map $`I:\stackrel{~}{\mathrm{\Sigma }X}\stackrel{~}{\mathrm{\Sigma }X}`$ in $`\stackrel{~}{\mathrm{\Sigma }X}\times \stackrel{~}{\mathrm{\Sigma }X}`$ given by $`I:(p,(g))(p,(g^1))`$. To obtain the orbifold structure, one can view $`\mathrm{\Delta }`$ as orbifold fiber product of identify map and $`I`$, which has an induced orbifold structure since both identify and $`I`$ are so called ”good map” (see \[CR\]). Consider map $`\mathrm{\Lambda }:\stackrel{~}{X}_3^o\times \stackrel{~}{X}_3^o\stackrel{~}{\mathrm{\Sigma }X}\times \stackrel{~}{\mathrm{\Sigma }X}`$ given by $`((p,(𝐠)),(q,(𝐡)))((p,(g_3)),(q,(h_1)))`$. We wish to consider the preimage of $`\mathrm{\Delta }`$. Remark: Suppose that we have two maps $$f:XZ,g:YZ.$$ In general, ordinary fiber product $`X\times _ZY`$ may not have a natural orbifold structure. The correct formulation is to use ”good map” introduced in \[CR\]. If $`f,g`$ are good maps, there is a canonical orbifold fiber product (still denoted by $`X\times _ZY`$) obtained by taking fiber product on uniformizing system. It has an induced orbifold structure and there are good map projection to both $`X,Y`$ to make appropriate diagram to commute. However, as a set, such an orbifold fiber product is not usual fiber product. Throughout this paper, we will use $`X\times _ZY`$ to denote orbifold fiber product only. It is clear that the pre-image of $`\mathrm{\Delta }`$ can be viewed as fiber product of $$e_3,Ie_1:\stackrel{~}{X}_3^0\stackrel{~}{X}.$$ Then, we define the pre-image $`\mathrm{\Lambda }^1(\mathrm{\Delta })`$ as orbifold fiber product of $`e_3,,Ie_1`$. It is easy to check that $`\mathrm{\Lambda }^1(\mathrm{\Delta })=\stackrel{~}{X}_4^o`$. Next, we describe explicitly the compatification $`\overline{}_4`$ of $`_4`$. Recall the moduli space of genus-zero curves with 4-marked points $`_{0,4}`$ can be identified with $`𝐏^1\{0,1,\mathrm{}\}`$ by fixing the first three marked points to be $`\{0,1,\mathrm{}\}`$. The Deligne-Mumford compactification $`\overline{}_{0,4}`$ is then identified with $`𝐏^1`$ where each point of $`\{0,1,\mathrm{}\}`$ corresponds to a nodal curve obtained as the last marked point is running into this point. It is easy to see that part of the compatification $`\overline{}_4`$ by adding a copy of $`\stackrel{~}{X}_4^o`$ at $`\mathrm{}\overline{}_{0,4}=𝐏^1`$ where intuitively we associate $`(g_1g_2)^1,g_1g_2`$ at nodal point. In the same way, the compatification at $`0`$ is by adding a copy of $`\stackrel{~}{X}_4^o`$ where we associate $`(g_1g_4)^1,g_1g_4`$ at nodal point, and at $`1`$ by associating $`(g_1g_3)^1,g_1g_3`$ at nodal point. Next, we define an orbifold bundle to measure the failure of transversality of $`\mathrm{\Lambda }`$ to $`\mathrm{\Delta }`$. Definition 4.3.1: We define a complex orbifold bundle $`\nu `$ over $`\mathrm{\Lambda }^1(\mathrm{\Delta })_{(g_1,g_2,g_3,g_4)}`$ as follows: over each uniformizing system $`(V_p^𝐠,C(𝐠))`$ of $`\mathrm{\Lambda }^1(\mathrm{\Delta }_{(𝐠)})`$, where $`𝐠=(g_1,g_2,g_3,g_4)`$, we regard $`V_p^𝐠`$ as the intersection of $`V_p^{g_1}V_p^{g_2}`$ with $`V_p^{g_3}V_p^{g_4}`$ in $`V_p^g`$ where $`g=(g_1g_2)^1`$. We define $`\nu `$ to be the complex orbifold bundle over $`\mathrm{\Lambda }^1(\mathrm{\Delta })`$ whose fiber is the orthogonal complement of $`V_p^{g_1}V_p^{g_2}+V_p^{g_3}V_p^{g_4}`$ in $`V_p^g`$. The associativity is based on the following Lemma 4.3.2: The complex orbifold bundle $`pr:E_4_4`$ can be extended over the compatification $`\overline{}_4`$, denoted by $`\overline{pr}:\overline{E}_4\overline{}_4`$, such that $`\overline{E}_4\|_{\{\}\times \stackrel{~}{X}_4^o}=(E_3E_3)\|_{\mathrm{\Lambda }^1(\mathrm{\Delta })}\nu `$ under the above identification, where $`\{\}`$ represents a point in $`\{0,1,\mathrm{}\}\overline{}_{0,4}`$. Proof: We fix an identification of infinite cylinder $`𝐑\times S^1`$ with $`𝐂^{}\{0\}`$ via the biholomorphism defined by $`t+ise^{(t+is)}`$ where $`t𝐑`$ and $`sS^1=𝐑/2\pi 𝐙`$. Through this identification, we regard a punctured Riemann surface as a Riemann surface with cylindrical ends. A neighborhood of a point $`\{0,1,\mathrm{}\}\overline{}_{0,4}`$, as a family of isomorphism classes of genus-zero curves with 4-marked points, can be described by a family of curves $`(\mathrm{\Sigma }_{r,\theta },𝐳)`$ obtained by gluing of two genus-zero curves with a cylindrical end and two marked points on each, parametrized by $`(r,\theta )`$ where $`0rr_0`$ and $`\theta S^1`$, as we glue the two curves by self-biholomorphisms of $`(\mathrm{ln}r,3\mathrm{ln}r)\times S^1`$ defined by $`(t,s)(4\mathrm{ln}rt,(s+\theta ))`$ ($`r=0`$ represents the nodal curve $``$). Likewise, thinking of points in $`_4`$ as equivalence classes of triples $`((\mathrm{\Sigma },𝐳),E_{p,𝐠},\mathrm{\Phi }_{p,𝐠})`$ where $`(\mathrm{\Sigma },𝐳)`$ is a genus-zero curve of 4-marked points $`𝐳`$, a neighborhood of $`\{\}\times (X\stackrel{~}{X}_4^o)`$ in $`\overline{}_4`$ are described by a family of holomorphic orbifold bundles on $`(\mathrm{\Sigma }_{r,\theta },𝐳)`$ with morphisms obtained by gluing two holomorphic orbifold bundles on genus-zero curves with two marked points and one cylindrical end on each. We denote them by $`(E_{r,\theta },\mathrm{\Phi }_{r,\theta })`$. The key is to construct a family of isomorphisms of complex orbifold bundle $$\mathrm{\Psi }_{r,\theta }:E_3E_3\nu \|_{\mathrm{\Lambda }^1(\mathrm{\Delta })}E_4$$ for $`(r,\theta )(0,r_0)\times S^1`$. Recall the fiber of $`E_3`$ and $`E_4`$ is given by kernels of the $`\overline{}^{}`$ operators. In fact, $`\mathrm{\Psi }_{r,\theta }`$ are given by gluing maps of kernels of $`\overline{}^{}`$ operators. More precisely, suppose $`((\mathrm{\Sigma }_{r,\theta },z),E_{r,\theta },\mathrm{\Phi }_{r,\theta })`$ are obtained by gluing $`((\mathrm{\Sigma }_1,𝐳_1),E_{p,𝐠},\mathrm{\Phi }_{p,𝐠})`$ and $`((\mathrm{\Sigma }_2,𝐳_2),E_{p,𝐡},\mathrm{\Phi }_{p,𝐡})`$ where $`𝐠=(g_1,g_2,g)`$ and $`𝐡=(g^1,h_2,h_3)`$. Let $`m=\|g\|`$. Then $`E_{r,\theta }\|_{(\mathrm{ln}r,3\mathrm{ln}r)\times S^1}`$ is uniformized by $`(\frac{\mathrm{ln}r}{m},\frac{3\mathrm{ln}r}{m})\times S^1\times TV_p`$ with an obvious action by $`𝐙_m=g`$. Let $`\xi _1\mathrm{\Omega }^{0,1}(E_{p,𝐠})`$, $`\xi _2\mathrm{\Omega }^{0,1}(E_{p,𝐡})`$ such that $`\overline{}^{}\xi _i=0`$ for $`i=1,2`$. On the cylindrical end, if we fix the local coframe $`d(t+is)`$, then each $`\xi _i`$ is a $`TV_p`$-valued, exponentially decaying holomorphic function on the cylindrical end. We fix a cut-off function $`\rho (t)`$ such that $`\rho (t)1`$ for $`t0`$ and $`\rho (t)0`$ for $`t1`$. We define the gluing of $`\xi _1`$ and $`\xi _2`$, which is a section of $`\mathrm{\Omega }^{0,1}(E_{r,\theta })`$ and denoted by $`\xi _1\mathrm{\#}\xi _2`$, by $$\xi _1\mathrm{\#}\xi _2=\rho (2\mathrm{ln}r+t)\xi _1+(1\rho (2\mathrm{ln}r+t))\xi _2$$ on the cylindrical end. Let $`\mathrm{\Psi }_{r,\theta }(\xi _1,\xi _2)`$ be the $`L^2`$-projection of $`\xi _1\mathrm{\#}\xi _2`$ onto $`\mathrm{ker}\overline{}^{}`$, then the difference $`\eta =\xi _1\mathrm{\#}\xi _2\mathrm{\Psi }_{r,\theta }(\xi _1,\xi _2)`$ satisfies the estimate $`\overline{}^{}\eta _{L^2}Cr^\delta (\xi _1+\xi _2)`$ for some $`\delta =\delta (\xi _1,\xi _2)>0`$. Hence $`\eta _{L^2}C\|\mathrm{ln}r\|r^\delta (\xi _1+\xi _2)`$ (cf. \[Ch\]), from which it follows that for small enough $`r`$, $`\mathrm{\Psi }_{r,\theta }`$ is an injective linear map. Now given any $`\xi V_p^g`$ which is orthogonal to both $`V_p^{g_1}V_p^{g_2}`$ and $`V_p^{g_3}V_p^{g_4}`$, we define $`\mathrm{\Psi }_{r,\theta }(\xi )`$ as follows: fixing a cut-off function, we construct a section $`u_\xi `$ over the cylindrical neck $`(\mathrm{ln}r,3\mathrm{ln}r)\times S^1`$ with support in $`(\mathrm{ln}r+1,3\mathrm{ln}r1)\times S^1`$ and equals $`\xi `$ on $`(\mathrm{ln}r+2,3\mathrm{ln}r2)\times S^1`$. We write $`\overline{}^{}u_\xi =v_{\xi ,1}+v_{\xi ,2}`$ where $`v_{\xi ,1}`$ is supported in $`(\mathrm{ln}r+1,\mathrm{ln}r+2)\times S^1`$ and $`v_{\xi ,2}`$ in $`(3\mathrm{ln}r2,3\mathrm{ln}r1)\times S^1`$. Since $`\xi `$ is orthogonal to both $`V_p^{g_1}V_p^{g_2}`$ and $`V_p^{g_3}V_p^{g_4}`$, we can arrange so that $`v_{\xi ,1}`$ is $`L^2`$-orthogonal to $`V_p^{g_1}V_p^{g_2}V_p^g`$ and $`v_{\xi ,2}`$ is $`L^2`$-orthogonal to $`V_p^{g^1}V_p^{g_3}V_p^{g_4}`$, which are the kernels of the $`\overline{}`$ operators on $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ acting on sections of $`E_{p,𝐠}`$ and $`E_{p,𝐡}`$ respectively. Hence there exist $`\alpha _1\mathrm{\Omega }^{0,1}(E_{p,𝐠})`$ and $`\alpha _2\mathrm{\Omega }^{0,1}(E_{p,𝐡})`$ such that $`\overline{}^{}\alpha _i=v_{\xi ,i}`$ and $`\alpha _i`$ are $`L^2`$-orthogonal to the kernels of the $`\overline{}^{}`$ operators respectively. We define $`\mathrm{\Psi }_{r,\theta }(\xi )`$ to be the $`L^2`$-orthogonal projection of $`u_\xi \alpha _1\mathrm{\#}\alpha _2`$ onto $`\mathrm{ker}\overline{}^{}`$, then $`\mathrm{\Psi }_{r,\theta }(\xi )`$ is linear on $`\xi `$. On the other hand, observe that $`\overline{}^{}(u_\xi \alpha _1\mathrm{\#}\alpha _2)_{L^2}Cr^\delta \xi `$ for some $`\delta >0`$, if let $`\eta `$ be the difference of $`\mathrm{\Psi }_{r,\theta }(\xi )`$ and $`u_\xi \alpha _1\mathrm{\#}\alpha _2`$, then $`\eta _{L^2}C\|\mathrm{ln}r\|r^\delta \xi `$ (cf. \[Ch\]), from which we see that for sufficiently small $`r>0`$, $`\mathrm{\Psi }_{r,\theta }(\xi )0`$ if $`\xi 0`$. Hence we construct a family of injective morphisms $$\mathrm{\Psi }_{r,\theta }:E_3E_3\nu \|_{\mathrm{\Lambda }^1(\mathrm{\Delta })}E_4$$ for $`(r,\theta )(0,r_0)\times S^1`$. We will show next that each $`\mathrm{\Psi }_{r,\theta }`$ is actually an isomorphism. We denote by $`\overline{}_i`$ the $`\overline{}`$ operator on $`\mathrm{\Sigma }_i`$, and $`\overline{}_{r,\theta }`$ the $`\overline{}`$ operator on $`\mathrm{\Sigma }_{r,\theta }`$. Then index formula tells us that (cf. Proposition 4.2.2) $`index\overline{}_1`$ $`=`$ $`n{\displaystyle \underset{j=1}{\overset{3}{}}}\iota (p,g_j),`$ $`index\overline{}_2`$ $`=`$ $`n{\displaystyle \underset{j=1}{\overset{3}{}}}\iota (p,h_j),`$ $`index\overline{}_{r,\theta }`$ $`=`$ $`n(\iota (p,g_1)+\iota (p,g_2)+\iota (p,h_2)+\iota (p,h_3)),`$ from which we see that $`index\overline{}_1+index\overline{}_2=index\overline{}_{r,\theta }+dim_𝐂V_p^g`$. Since $`dim\mathrm{ker}\overline{}_1+dim\mathrm{ker}\overline{}_2=dim\mathrm{ker}\overline{}_{r,\theta }+dim_𝐂V_p^grank\nu `$, we have $$dimcoker\overline{}_1+dimcoker\overline{}_2+rank\nu =dimcoker\overline{}_{r,\theta }.$$ Hence $`\mathrm{\Psi }_{r,\theta }`$ is an isomorphism for each $`(r,\theta )`$. $`\mathrm{}`$ Before we prove the associativity, let’s review some of basic construction of smooth manifold and its orbifold analogue. Recall that if $`ZX`$ is a submanifold, then Poincare dual of $`Z`$ can be constructed by Thom form of normal bundle $`N_Z`$ via the natural identification between normal bundle and tubuler neighborhood of $`Z`$. Here, Thom form $`\mathrm{\Theta }_Z`$ is a close form such that its restriction on each fiber is a compact supported form of top degree with volume one. In orbifold category, the same is true provided that we interpret “suborbifold” correctly. Here, a suborbifold is a good map $`f:ZX`$ such that locally, $`f`$ can be lifted to a $`G`$-invariant embedding to “general” uniformizing system $`\stackrel{~}{f}:(U_Z,G,\pi _Z)(U_X,G,\pi _X)`$. Here, “general” means that $`U_Z,U_X`$ could be disconnected. For example, orbifold fiber product $`\mathrm{\Lambda }^1(\mathrm{\Delta })`$ is a suborbifold of $`\stackrel{~}{X}_3^o\times \stackrel{~}{X}_3^0`$. It is clear that Poincare dual of $`Z`$ can be represented by Thom class of normal bundle $`Z`$. Proposition 4.3.4: Choose a basis $`\{e_j\},\{e_k^o\}`$ of the total orbifold cohomology group $`H_{orb}^{}(X),H_{orb,c}^{}(X)`$ such that each $`e_j,e_k^o`$ is of homogeneous degree. Let $`<e_j,e_k^o>_{orb}=a_{jk}`$ be the Poincare pairing matrix and $`(a^{jk})`$ be the inverse. Then, $$_{(\stackrel{~}{X}_4^o)_{(𝐠)}}^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(\alpha _3)e_4^{}(e_l^o)e(E_4)$$ $$=\underset{j,k}{}(_{\stackrel{~}{X}_3^o}^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(e_k^o)e(E_3))(_{\stackrel{~}{X}_3^o}^{orb}e_1^{}(e_j)e_2^{}(\alpha _3)e_3^{}(e_l^o)e(E_3))a^{kj}$$ Proof: Key observation is $`\mathrm{\Lambda }^{}N_\mathrm{\Delta }=N_{\mathrm{\Lambda }^1(\mathrm{\Delta })}\nu .`$ Hence, $`\mathrm{\Lambda }^{}\mathrm{\Theta }_\mathrm{\Delta }=\mathrm{\Theta }_{\mathrm{\Lambda }^1(\mathrm{\Delta })}\mathrm{\Theta }_\nu `$. $$\begin{array}{cc}& _{\stackrel{~}{X}_4^o}^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(\alpha _3)e_4^{}(e_l^o)e(E_4)\hfill \\ =\hfill & _{\mathrm{\Lambda }^1(\mathrm{\Delta })}^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(\alpha _3)e_4^{}(e_l^o)e(E_3)e(E_3)e(\nu )\hfill \\ =\hfill & _{\stackrel{~}{X}_3^o\times \stackrel{~}{X}_3^o}^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(\alpha _3)e_4^{}(e_l^o)e(E_3)e(E_3)\mathrm{\Theta }_\nu \mathrm{\Theta }_{\mathrm{\Lambda }^1(\mathrm{\Delta })}\hfill \\ =\hfill & _{\stackrel{~}{X}_3^o\times \stackrel{~}{X}_3^o}^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(\alpha _3)e_4^{}(e_l^o)e(E_3)e(E_3)\mathrm{\Lambda }^{}\mathrm{\Theta }_\mathrm{\Delta }\hfill \\ =\hfill & _{j,k}(_{\stackrel{~}{X}_3^o}^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(e_k^o)e(E_3)))(_{\stackrel{~}{X}_3^o}^{orb}e_1^{}(e_j)e_2^{}(\alpha _3)e_3^{}(e_l^o)e(E_3))a^{kj}\hfill \end{array}$$ Now we are ready to prove Proposition 4.3.5: The cup product $`_{orb}`$ is associative, i.e., for any $`\alpha _i`$, $`i=1,2,3`$, we have $$(\alpha _1_{orb}\alpha _2)_{orb}\alpha _3=\alpha _1_{orb}(\alpha _2_{orb}\alpha _3).$$ Proof: By definition of cup product $`_{orb}`$, we have $`(\alpha _1_{orb}\alpha _2)_{orb}\alpha _3`$ equals $$\underset{j,k,l,s}{}(__3^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(e_k^o)e(E_3))(__3^{orb}e_1^{}(e_j)e_2^{}(\alpha _3)e_3^{}(e_l^o)e(E_3))a^{kj}a^{ls}e_s$$ and $`\alpha _1_{orb}(\alpha _2_{orb}\alpha _3)`$ equals $$\underset{j,k,l,s}{}(__3^{orb}e_1^{}(\alpha _1)e_2^{}(e_j)e_3^{}(e_l^o)e(E_3))(__3^{orb}e_1^{}(\alpha _2)e_2^{}(\alpha _3)e_3^{}(e_k^o)e(E_3))a^{kj}a^{ls}e_s.$$ By Proposition 4.3.4, $$\underset{j,k}{}(__3^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(e_k^o)e(E_3))(__3^{orb}e_1^{}(e_j)e_2^{}(\alpha _3)e_3^{}(e_l^o)e(E_3))a^{kj}$$ equals $$_{\{\mathrm{}\}\times \stackrel{~}{X}_4^o}^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(\alpha _3)e_4^{}(e_l^o)e(E_4),$$ and $$\underset{j,k}{}(__3^{orb}e_1^{}(\alpha _1)e_2^{}(e_j)e_3^{}(e_l^o)e(E_3))(__3^{orb}e_1^{}(\alpha _2)e_2^{}(\alpha _3)e_3^{}(e_k^o)e(E_3))a^{kj}$$ equals $$_{\{0\}\times \stackrel{~}{X}_4^o}^{orb}e_1^{}(\alpha _1)e_2^{}(\alpha _2)e_3^{}(\alpha _3)e_4^{}(e_l^o)e(E_4).$$ Hence $`(\alpha _1_{orb}\alpha _2)_{orb}\alpha _3=\alpha _1_{orb}(\alpha _2_{orb}\alpha _3)`$. $`\mathrm{}`$ ## 5 Examples In general, it is easy to compute orbifold cohomology once we know the action of local group. Example 5.1-Kummer surface: Consider Kummer surface $`X=T^4/\tau `$, where $`\tau `$ is the involution $`xx`$. $`\tau `$ has 16 fixed points, which give 16 twisted sectors. It is easily seen that $`\iota _{(\tau )}=1`$. Hence, we should shift the cohomology classes of a twisted sector by $`2`$ to obtain 16 degree two classes in orbifold cohomology. The cohomology classes of nontwisted sector come from invariant cohomology classes of $`T^4`$. It is easy to compute that $`H^0(X,𝐑),H^4(X,𝐑)`$ has dimension one and $`H^2(X,𝐑)`$ has dimension 6. Hence, we obtain $$b_0^{orb}=b_4^{orb}=1,b_1^{orb}=b_3^{orb}=0,b_2^{orb}=22.$$ Note that orbifold cohomology group of $`T^4/\tau `$ is isomorphic to ordinary cohomology of $`K3`$-surface, which is the the crepant resolution of $`T^4/\tau `$. However, it is easy to compute that Poincare pairing of $`H_{orb}^{}(T^4/\tau ,𝐑)`$ is different from Poincare pairing of $`K3`$-surface. We leave it to readers Example 5.2-Borcea-Voisin threefold: An important class of Calabi-Yau 3-folds due to Borcea-Voisin is constructed as follows: Let $`E`$ be an elliptic curve with an involution $`\tau `$ and $`S`$ be a $`K3`$-surface with an involution $`\sigma `$ acting by $`(1)`$ on $`H^{2,0}(S)`$. Then, $`\tau \times \sigma `$ is an involution of $`E\times S`$, and $`X=E\times S/<\tau \times \sigma >`$ is a Calabi-Yau orbifold. The crepant resolution $`\stackrel{~}{X}`$ of $`X`$ is a smooth Calabi-Yau 3-fold. This class of Calabi-Yau 3-folds occupy an important place in mirror symmetry. Now, we want to compute the orbifold Dolbeault cohomology of $`X`$ to compare with Borcea-Voisin’s calculation of Dolbeault cohomology of $`\stackrel{~}{X}`$. Let’s give a brief description of $`X`$. Our reference is \[Bo\]. $`\tau `$ has 4 fixed points. $`(S,\sigma )`$ is classified by Nikulin. Up to deformation, it is decided by three integers $`(r,a,\delta )`$ with following geometric meaning. Let $`L^\sigma `$ be the fixed part of $`K3`$-lattice. Then, $$r=rank(L^\sigma ),(L^\sigma )^{}/L^\sigma =(𝐙/2𝐙)^a.$$ $`(5.1)`$ $`\delta =0`$ if the fixed locus $`S_\sigma `$ of $`\sigma `$ represents a class divisible by $`2`$. Otherwise $`\delta =1`$. There is a detail table for possible value of $`(r,a,\delta )`$ \[Bo\]. The cases we are interested in are $`(r,a,\delta )(10,10,0)`$, where $`S_\sigma \mathrm{}`$. When $`(r,a,\delta )(10,8,0)`$, $$S_\sigma =C_gE_1\mathrm{},E_k$$ $`(5.2)`$ is a disjoint union of a curve $`C_g`$ of genus $$g=\frac{1}{2}(22ra)$$ and $`k`$ rational curves $`E_i`$, with $$k=\frac{1}{2}(ra).$$ For $`(r,a,\delta )=(10,8,0)`$, $$S_\sigma =C_1\stackrel{~}{C}_1$$ the disjoint union of two elliptic curves. Now, let’s compute its orbifold Dolbeault cohomology. We assume that $`(r,a,\delta )(10,8,0)`$. The case that $`(r,a,\delta )=(10,8,0)`$ can be computed easily as well. We leave it as an exercise for the readers. An elementary computation yields $$h^{1,0}(X)=h^{2,0}(X)=0,h^{3,0}(X)=1,h^{1,1}(X)=r+1,h^{2,1}(X)=1+(20r).$$ $`(5.3)`$ Notes that twisted sectors consist of 4 copies of $`S_\sigma `$. $$h^{0,0}(S_\sigma )=k+1,h^{1,0}(S_\sigma )=g.$$ $`(5.4)`$ It is easy to compute that the degree shifting number for twisted sectors is 1. Therefore, we obtain $$h_{orb}^{1,0}=h_{orb}^{2,0}=0,h_{orb}^{3,0}=1,h_{orb}^{1,1}=1+r+4(k+1),h_{orb}^{2,1}=1+(20r)+4g.$$ $`(5.5)`$ Compared with the calculation for $`\stackrel{~}{X}`$, we get a precise agreement. Next, we compute the triple product on $`H_{orb}^{1,1}`$. $`H_{orb}^{1,1}`$ consists of contributions from nontwisted sector with dimension $`1+r`$ and twisted sectors with dimension $`4(k+1)`$. Only nontrivial one is the classes from twisted sector. Recall that we need to consider the moduli space of 3-point ghost maps with weight $`g_1,g_2,g_3`$ at three marked points satisfying the condition $`g_1g_2g_3=1`$. In our case, the only possibility is $`g_1=g_2=g_3=\tau \times \sigma `$. But $`(\tau \times \sigma )^3=\tau \times \sigma 1`$. Therefore, For any class $`\alpha `$ from twisted sectors, $`\alpha ^3=0`$. On the other hand, we know the triple product or exceptional divisor of $`\stackrel{~}{X}`$ is never zero. Hence, $`X,\stackrel{~}{X}`$ have different cohomology ring. Example 5.3-Weighted projective space: The examples we compute so far are global quotient. Weighted projective spaces are the easiest examples of non-global quotient orbifolds. Let’s consider weighted projective space $`CP(d_1,d_2)`$, where $`(d_1,d_2)=1`$. Thurston’s famous tear drop is $`CP(1,d)`$. $`CP(d_1,d_2)`$ can be defined as the quotient of $`S^3`$ by $`S^1`$, where $`S^1`$ acts on the unit sphere of $`𝐂^2`$ by $$e^{i\theta }(z_1,z_2)=(e^{id_1\theta }z_1,e^{id_2\theta }z_2).$$ $`(5.6)`$ $`CP(d_1,d_2)`$ has two singular points $`x=[1,0],y=[0,1]`$. $`x,y`$ gives rise $`d_21,d_11`$ many twisted sectors indexed by the elements of isotropy subgroup. The degree shifting numbers are $`\frac{i}{d_2},\frac{j}{d_1}`$ for $`1id_21,1jd_11`$. Hence, the orbifold cohomology are $$h_{orb}^0=h_{orb}^2=h_{orb}^{\frac{2i}{d_2}}=h_{orb}^{\frac{2j}{d_1}}=1.$$ $`(5.7)`$ Note that orbifold cohomology classes from twisted sectors have rational degree. Let $`\alpha H_{orb}^{\frac{2}{d_1}},\beta H_{orb}^{\frac{2}{d_2}}`$ be the generators corresponding to $`1H^0(pt,𝐂)`$. An easy computation yields that orbifold cohomology is generated by $`\{1,\alpha ^j,\beta ^i\}`$ with relation $$\alpha ^{d_1}=\beta ^{d_2},\alpha ^{d_1+1}=\beta ^{d_2+1}=0.$$ $`(5.8)`$ The Poincare pairing is for $`1i_1,i_2,i<d_21,1j_1,j_2,j<d_11`$ $$<\beta ^i,\alpha ^j>_{orb}=0,<\beta ^{i_1},\beta ^{i_2}>_{orb}=\delta _{i_1,d_2i_2},<\alpha ^{j_1},\alpha ^{j_2}>_{orb}=\delta _{j_1,d_1j_2}.$$ The last two examples are local examples in nature. But they exhibit a strong relation with group theory. Example 5.4: The easiest example is probably a point with a trivial group action of $`G`$. In this case, a sector $`X_{(g)}`$ is a point with the trivial group action of $`C(g)`$. Hence, orbifold cohomology is generated by conjugacy classes of elements of $`G`$. All the degree shifting numbers are zero. Only Poincare pairing and cup products are interesting. Poincare paring is obvious. Let’s consider cup product. First we observe that $`X_{(g_1,g_2,(g_1g_2)^1)}`$ is a point with the trivial group action of $`C(g_1)C(g_2)`$. We choose a basis $`\{x_{(g)}\}`$ of the orbifold cohomology group where $`x_{(g)}`$ is given by the constant function $`1`$ on $`X_{(g)}`$. Then the inverse of the intersection matrix $`(<x_{(g_1)},x_{(g_2)}>_{orb})`$ has $`a^{x_{(g)}x_{(g^1)}}=\|C(g)\|`$. Now by Lemma 4.1.4 and Equation $`(\mathrm{4.2.12})`$, we have $$x_{(g_1)}x_{(g_2)}=\underset{(h_1,h_2),h_1(g_1),h_2(g_2)}{}\frac{\|C(h_1h_2)\|}{\|C(h_1)C(h_2)\|}x_{(h_1h_2)},$$ where $`(h_1,h_2)`$ is the conjugacy class of pair $`h_1,h_2`$. On the other hand, recall that the center $`Z(𝐂[G])`$ of group algebra $`𝐂[G]`$ is generated by $`_{h(g)}h`$. We can define a map from the orbifold cohomology group onto $`Z(𝐂[G])`$ by $$\mathrm{\Psi }:x_{(g)}\underset{h(g)}{}h.$$ $`(5.9)`$ The map $`\mathrm{\Psi }`$ is a ring homomorphism, which can be seen as follows: $$(\underset{h(g_1)}{}h)(\underset{k(g_2)}{}k)=\underset{h(g_1),k(g_2)}{}hk=\underset{(h_1,h_2),h_1(g_1),h_2(g_2)}{}\frac{A}{B}(\underset{h(h_1h_2)}{}h),$$ $`(5.10)`$ where $`A=\frac{\|G\|}{\|C(h_1)C(h_2)\|}`$ is the number of elements in the orbit of $`(h_1,h_2)`$ of the action of $`G`$ given by $`g(h_1,h_2)=(gh_1g^1,gh_2g^1)`$, and $`B=\frac{\|G\|}{\|C(h_1h_2)\|}`$ is the number of elements in the orbit of $`h_1h_2`$ of the action of $`G`$ given by $`gh=ghg^1`$. Therefore, the orbifold cup product is the same as product of $`Z(𝐂[G])`$, and the orbifold cohomology ring can be identified with the center $`Z(𝐂[G])`$ of group algebra $`𝐂[G]`$ via $`(5.9)`$. Example 5.5: Suppose that $`GSL(n,𝐂)`$ is a finite subgroup. Then, $`𝐂^n/G`$ is an orbifold. $`H^{p,q}(X_{(g)},𝐂)=0`$ for $`p>0`$ or $`q>0`$ and $`H^{0,0}(X_{(g)},𝐂)=𝐂`$. Therefore, $`H_{orb}^{p,q}=0`$ for $`pq`$ and $`H_{orb}^{p,p}`$ is a vector space generated by conjugacy class of $`g`$ with $`\iota _{(g)}=p`$. Therefore, we have a natural decomposition $$H_{orb}^{}(𝐂^n/G,𝐂)=Z[𝐂[G])=\underset{p}{}H_p,$$ $`(5.11)`$ where $`H_p`$ is generated by conjugacy classes of $`g`$ with $`\iota _{(g)}=p`$. The ring structure is also easy to describe. Let $`x_{(g)}`$ be generator corresponding to zero cohomology class of twisted sector $`X_{(g)}`$. We would like to get a formula for $`x_{(g_1)}x_{(g_2)}`$. As we showed before, the multiplication of conjugacy classes can be described in terms of center of twisted group algebra $`Z(𝐂[G])`$. But we have further restrictions in this case. Let’s first describe the moduli space $`X_{(h_1,h_2,(h_1h_2)^1)}`$ and its corresponding GW-invariants. It is clear $$X_{(h_1,h_2,(h_1h_2)^1)}=X_{h_1}X_{h_2}/C(h_1,h_2).$$ To have nonzero invariant, we require that $$\iota _{(h_1h_2)}=\iota _{(h_1)}+\iota _{(h_2)}.$$ $`(5.12)`$ Then, we need to compute $$_{X_{h_1}X_{h_2}/C(h_1,h_2)}^{orb}e_3^{}(vol_c(X_{h_1h_2}))e(E),$$ $`(5.13)`$ where $`vol_c(X_{h_1h_2})`$ is the compact supported $`C(h_1h_2)`$-invariant top form with volume one on $`X_{h_1h_2}`$. It is also viewed as a form on $`X_{h_1}X_{h_2}/C(h_1)C(h_2)`$. However, $$X_{h_1}X_{h_2}X_{h_1h_2}$$ is a submanifold. Therefore, (5.13) is zero unless $$X_{h_1}X_{h_2}=X_{h_1h_2}.$$ $`(5.14)`$ In this case, we call $`(h_1,h_2)`$ transverse. In this case, it is clear that obstruction bundle is trivial. Let $$I_{g_1,g_2}=\{(h_1,h_2);h_i(g_i),\iota _{(h_1)}+\iota _{(h_2)}=\iota _{(h_1h_2)},(h_1,h_2)transverse\}.$$ $`(5.15)`$ Then, using decomposition lemma 4.1.4, $$x_{(g_1)}x_{(g_2)}=\underset{(h_1,h_2)I_{g_1,g_2}}{}d_{(h_1,h_2)}x_{(h_1h_2)}.$$ $`(5.16)`$ A similar computation as previous example yields $`d_{(h_1,h_2)}=\frac{\|C(h_1h_2)\|}{\|C(h_1)C(h_2)\|}`$. ## 6 Some General Remarks Physics indicated that orbifold quantum cohomology should be equivalent to ordinary quantum cohomology of crepant resolution. It is a rather difficult problem to find the precise relations between orbifold quantum cohomology with the quantum cohomology of a crepant resolution. At the classical level, there is an indication that equivariant $`K`$-theory is better suited for this purpose. For GW-invariant, orbifold GW-invariant defined in \[CR\] seems to be equivalent to the relative GW-invariant of pairs studied by Li-Ruan \[LR\]. We hope that we will have a better understanding of this relation in the near future. There are many interesting problems in this orbifold cohomology theory. As we mentioned at the beginning, many Calabi-Yau 3-folds are constructed as crepant resolutions of Calabi-Yau orbifolds. The orbifold string theory suggests that there might be a mirror symmetry phenomenon for Calabi-Yau orbifolds. Another interesting question is the relation between quantum cohomology and birational geometry \[R\],\[LR\]. In fact, this was our original motivation. Namely, we want to investigate the change of quantum cohomology under birational transformations. Birational transformation corresponds to wall crossing phenomenon for symplectic quotients. Here, the natural category is symplectic orbifolds instead of smooth manifolds. From our work, it is clear that we should replace quantum cohomology by orbifold quantum cohomology. Then, it is a challenge problem to calculate the change of orbifold quantum cohomology under birational transformation. The first step is to investigate the change of orbifold cohomology under birational transformation. This should be an interesting problem in its own right.
warning/0004/cond-mat0004118.html	ar5iv	text	# Temperature dependence of the nematic anchoring energy: mean field analysis ## Abstract In the mean field approximation, we evaluate the temperature dependence of the anchoring energy strength of a nematic liquid crystal in contact with a solid substrate due to thermal fluctuations. Our study is limited to the weak anchoring case, where the microscopic surface energy is small with respect to the mean field energy due to the nematic phase. We assume furthermore that the physical properties of the substrate can be considered temperature independent in the range of the nematic phase. According to the thermodynamical perturbative approach, the macroscopic surface energy is deduced by averaging the microscopic one, with a density matrix containing only the nematic mean field. We show that the thermal renormalization of the anchoring energy coefficients is proportional to the generalized nematic order parameters. Our analysis shows also that the thermal renormalization of the anchoring energy coefficients predicted by means of Landau-like theories is a first and rather rough approximation in the whole nematic temperature range. One of the least-understood areas of physics and chemistry of liquid crystals concern the anchoring phenomenon and the temperature surface transitions at the interface liquid crystal-solid or soft substrate. These phenomena are important also from a practical point of view, since they play fundamental role in the realization of displays. In this Letter we analyse these phenomena in nematic liquid crystal media. Nematic liquid crystals are anisotropic fluids made by anisometric molecules having quadrupolar symmetry. Their intermolecular interaction, $`V_N`$, is such to orient the molecular axes, $`𝐮`$, along a common direction $`𝐧`$, called director . It coincides with the optical axis of the medium. When a nematic liquid crystal is in contact with a substrate, the orientation of $`𝐧`$ at the surface results from a balance of the anisotropic interactions with the bulk and with the substrate. In the absence of bulk distortions, the surface orientation of the nematic director coincides with the “easy axis”, $`𝐧_0`$ . It is such to minimize the anisotropic part of the surface energy characterizing the interface between the nematic and the substrate. Long ago Bouchiat and Langevin-Cruchon found a strong temperature dependence of the easy axis. The measurements of Ref. have been repeated by other groups with similar results . Several models have been proposed to interpret this phenomenon. According to Parson the easy axis results from the competition between dipolar and quadrupolar interactions, which depend on the temperature in different manner. In special situations a surface nematic orientation temperature dependent can be observed. The idea of Parson has been generalized by Sluckin and Poniewierski and by Sen and Sullivan . In all the models, the temperature surface transitions are due to a temperature dependence of the anisotropic part of the anchoring energy, which depends on the symmetry of the substrate and on the symmetry of the nematic phase. The aim of our paper is to analyze the temperature dependence of the anisotropic part of the surface energy. We assume that the nematic liquid crystal is not polar. From the molecular point of view it has a quadrupolar symmetry, whose principal axis coincides with $`𝐮`$. The elements of the relevant tensor are $`q_{ij}=(3/2)[u_iu_j(1/3)\delta _{ij}]`$. In our analysis, we expand the surface energy in series of spherical harmonic functions. The coefficients of the expansion are the experimentally detectable anchoring coefficients. According to our model, all the anchoring coefficients of the same order depend on the temperature in the same manner. From this result it follows that in nematic liquid crystals the alignment transitions driven by the surface (the so called temperature surface transitions) are due to a surface anchoring energy which contains contributions of different orders. We analyze the temperature dependence of the anchoring energy using an approach based on the mean field theory. In our analysis we neglect all the inhomogeneities. We assume, furthermore that the surface potential is short range. Let us consider a surface molecule of the nematic liquid crystal. It is submitted to the mean field due to the other nematic molecules, whose corresponding energy is $`V_N`$, and to the interaction with the substrate, $`V_S`$. In this framework, the total energy, $`V`$, of a given surface molecule is $`V=V_N+V_S`$. If $`V_NV_S`$ the extrapolation length $`b=K/WaV_N/V_S`$, where $`K`$ is an average elastic constant and $`a`$ a molecular dimension, is of the order of a molecular dimension . In this case, in the continuum limit it is possible to put $`b=0`$, and assume that the surface nematic orientation is fixed by the surface interaction. This situation is known as the strong anchoring case, and it is not interesting for us here. The interesting case is the one in which $`V_NV_S`$, corresponding to a situation where $`ba`$. This case corresponds to the weak anchoring situation, to which we will limit our investigation. In our analysis the small parameter used to expand the surface energy in power series is $`V_S/V_N1`$, in the weak anchoring situation. On the contrary, the surface scalar order parameter $`S`$ is not supposed to be a small quantity. $`V_N`$ describes the tendency of $`𝐮`$, which defines the molecular orientation, to be oriented along the nematic director $`𝐧`$. Usually, it is approximated by means of the Maier-Saupe’s mean field , $`V_N^M`$, according to which $`V_N^Mn_iq_{ij}n_j=P_2(𝐧𝐮)`$, where $`P_2`$ is the second order Legendre Polynomial. In this framework $`V_N^M=vP_2(𝐧𝐮)S`$, where $`v`$ is a molecular constant and $`S=P_2(𝐧𝐮)`$ the nematic scalar parameter. A generalization of the Maier-Saupe theory has been proposed by Humphries et al. . According to this generalized mean field theory the nematic mean field is given by $$V_N(𝐧𝐮)=\underset{l}{}v_{2l}P_{2l}(𝐧𝐮)S_{2l},$$ (1) where $`v_{2l}`$ are molecular parameters, and $`S_{2l}=P_{2l}(𝐧𝐮)`$ the nematic order parameters, given by the self-consistent equations $$S_{2l}=\frac{_0^1P_{2l}(𝐧𝐮)\mathrm{exp}[\beta V_N(𝐧𝐮)]d(𝐧𝐮)}{_0^1\mathrm{exp}[\beta V_N(𝐧𝐮)]d(𝐧𝐮)}.$$ (2) The Maier-Saupe potential, $`V_N^M`$, is obtained from $`V_N`$ putting $`v_{2k}=v\delta _{1,k}`$. The interaction connected to $`V_S`$ has to describe the tendency of the surface to orient the surface nematic molecules along the “easy direction”, $`𝐧_0`$. This direction depends on the symmetry of the surface and on the molecular properties of the mesophase. Since we limit our analysis to non polar media, $`V_S`$ has to be an even function of $`𝐮`$. It follows that $`V_S`$ is, actually, a function of the tensor $`\stackrel{}{q}`$ and can be written, in general, as $`V_S(𝐮)=V_S(\stackrel{}{q})=_kw_k(0)L_k(\stackrel{}{q})`$, where $`L_k(\stackrel{}{q})`$ indicate the scalar quantities we can build with the molecular tensor of elements $`q_{ij}=(3/2)[u_iu_j(1/3)\delta _{ij}]`$ and the elements of symmetry characterizing the surface. Each term of the expansion of $`V_S(\stackrel{}{q})`$ represents a given interaction, like induced dipole-induced dipole or quadrupole-quadrupole and so on ; the “intrinsic”anchoring coefficients $`w_k(0)`$ are physical parameters connected with the type of interaction described by $`L_k(\stackrel{}{q})`$. $`w_k(0)`$ refer to specific fundamental interactions, and are assumed to be temperature independent. In this case thermal effects arise only from the temperature dependence of the degree of alignment of the nematic molecules. This conclusion is valid only if in the temperature range of the nematic phase the physical properties of the substrate can be considered constant. In the opposite case $`w_k(0)`$ depend also on the temperature, via the substrate. Since we assume that $`w_k(0)`$ are temperature independent, our theory works well when the substrate is a solid crystal. Deviations from our prediction are expected for nematic samples oriented by means of surfactants. For our future considerations it is useful to describe the molecular direction and the nematic director in terms of the polar angles with respect to a cartesian reference frame having the $`z`$-axis parallel to the geometrical normal to the flat surface and the $`x`$-axis along the possible surface anisotropy. Let $`\mathrm{\Theta },\mathrm{\Phi }`$ and $`\theta ,\varphi `$ be the polar and azimuthal angles defining $`𝐮`$ and $`𝐧`$, respectively. Consequently $$V_S(\stackrel{}{q})=V_S(\mathrm{\Theta },\mathrm{\Phi })=\underset{k}{}w_k(0)L_k(\mathrm{\Theta },\mathrm{\Phi }).$$ (3) By decomposing the functions $`L_k(\mathrm{\Theta },\mathrm{\Phi })`$ in series of spherical harmonics functions $`Y_k^m(\mathrm{\Theta },\mathrm{\Phi })`$ we obtain $`L_k(\mathrm{\Theta },\mathrm{\Phi })=_ma_k^mY_k^m(\mathrm{\Theta },\mathrm{\Phi })`$. Since $`L_k=L_k(\stackrel{}{q})`$ and hence $`L_k(\mathrm{\Theta },\mathrm{\Phi })=L_k(\pi \mathrm{\Theta },\pi +\mathrm{\Phi })`$ for all $`k`$, we deduce that $`k=2l`$. It follows that for non-polar nematic liquid crystals $`L_{2l}(\mathrm{\Theta },\mathrm{\Phi })=_ma_{2l}^mY_{2l}^m(\mathrm{\Theta },\mathrm{\Phi })`$, and the microscopic surface energy can be written as $$V_S(\mathrm{\Theta },\mathrm{\Phi })=\underset{l}{}w_{2l}(0)\underset{m}{}a_{2l}^mY_{2l}^m(\mathrm{\Theta },\mathrm{\Phi }).$$ (4) The macroscopic anchoring energy $`W(𝐧)=W(\theta ,\varphi )`$ is obtained by averaging $`V_S`$ over the molecular orientations $`𝐮`$, or over $`\mathrm{\Theta }`$ and $`\mathrm{\Phi }`$. Since in the problem under consideration $`V_SV_N`$, $`V_S`$ can be treated as a perturbation. According to the thermodynamic perturbation theory we have $`W(\theta ,\varphi )=V_S(\mathrm{\Theta },\mathrm{\Phi })`$, and hence, as it follows from Eq.(4), $$W(\theta ,\varphi )=\underset{l}{}w_{2l}(0)\underset{m}{}a_{2l}^mY_{2l}^m(\mathrm{\Theta },\mathrm{\Phi }),$$ (5) where $`A=Tr(\rho A)/Tr(\rho )`$, and $`\rho =\mathrm{exp}(\beta V_N)`$ is the density matrix. In order to derive the macroscopic surface energy $`W(\theta ,\varphi )`$ we have first to express $`V_S(\mathrm{\Theta },\mathrm{\Phi })`$ in terms of a polar coordinates system based on the director $`𝐧`$ as polar axis. The cartesian reference frame has to be rotated in such a way that $`𝐳^{}=𝐧`$. We will indicate with $`\vartheta ,\phi `$ the polar angles of $`𝐮`$ with respect to the rotated coordinate system. In this case $$Y_l^m(\mathrm{\Theta },\mathrm{\Phi })=\underset{m^{}}{}D_{m,m^{}}^l(\theta ,\varphi )Y_l^m^{}(\vartheta ,\phi ),$$ (6) where $`D_{m,m^{}}^l(\theta ,\varphi )`$ are the elements of Wigner’s matrix. Since $`Y_l^m^{}(\vartheta ,\phi )=Y_l^0(\vartheta )\delta _{m^{},0}`$ we obtain from Eq.(6) $`Y_l^m(\mathrm{\Theta },\mathrm{\Phi })=D_{m,0}^l(\theta ,\varphi )Y_l^0(\vartheta )`$. By taking into account that $`D_{m,0}^l(\theta ,\varphi )=Y_l^m(\theta ,\varphi )`$, we have finally, as it follows from Eqs.(3,4,5), $$W(\theta ,\varphi )=\underset{l}{}w_{2l}(0)S_{2l}L_{2l}(\theta ,\varphi ),$$ (7) where we have taken into account that $`Y_{2l}^0(\vartheta ,\phi )=P_{2l}(\mathrm{cos}\vartheta )`$. Eq.(7) is a consequence of the fact that we regard all anisotropic effects as perturbation, so that they do not need to be included in the computation of the averages values. There is axial symmetry about the direction of $`𝐧`$ in the imperturbed system and only the member $`m=0`$ of the $`Y_l^m`$ is different from zero. By comparing Eq.(7) with Eq.(3) we deduce that the temperature dependence of the parameters describing the anisotropic part of the surface energy is given by $$w_{2l}(T)=w_{2l}(0)S_{2l}.$$ (8) This means that the temperature dependence of $`w_{2l}(T)/w_{2l}(0)`$ coincides with the temperature dependence of the $`2l`$-th scalar order parameter. According to the analysis presented above, where the macroscopic anchoring energy is given by the series expansion in spherical harmonic functions shown in Eq.(7), the thermal renormalization of the anchoring coefficients is given by Eq.(8). ¿From these results it follows that the anchoring coefficients of the same order in the expansion have the same temperature dependence. Consequently, in the frame of our model, temperature surface transitions are possible only in nematic samples whose anchoring energy contains contributions from different order in the spherical harmonic functions expansion. The ratios $`S_{2l}/Svs.S`$, for $`l=2,3`$ and $`4`$, in the Maier-Saupe approximation, can be easily evaluated in the nematic phase, where $`0.4S0.8`$. A direct calculation shows that $`S_{2l}/S0.2`$, for $`l=3,4`$, as it is shown in Fig.1. This explains why, usually, the anisotropic part of the surface anchoring energy given by Eq.(7) is well approximated by few terms . In the low temperature region, where $`\beta V_N1`$, the fluctuations of $`𝐮`$ with respect to $`𝐧`$ are small. In this region $`𝐧𝐮=\mathrm{cos}\vartheta 1(1/2)\vartheta ^2+𝒪(4)`$, i.e. $`\vartheta 1`$, and $`P_{2l}(\mathrm{cos}\vartheta )=1[l(2l+1)/2]\vartheta ^2+𝒪(4)`$. Consequently, from Eq.(2), the order parameter $`S_{2l}`$ is found to be $$S_{2l}1\frac{l(2l+1)}{B}\mathrm{exp}\left\{\frac{l(2l+1)}{B}\right\},$$ (9) where $`B=\beta _kk(2k+1)v_{2k}S_{2k}`$. The main nematic scalar order parameter $`S=P_2(𝐧𝐮)`$ is given by $`S=\mathrm{exp}(3/B)`$,as it follows from Eq.(9). The other order parameters can be determined in terms of $`S`$ by $`S_{2l}=S^{l(2l+1)/3}`$. In the low temperature region, the thermal renormalization of the anchoring coefficient is then given by $$w_{2l}(T)=w_{2l}(0)S^{\frac{l(2l+1)}{3}}.$$ (10) In particular, in this range of temperature, $`w_2(T)/w_2(0)=S`$ and $`w_4(T)/w_4(0)=S^{10/3}`$. The temperature dependence given by Eq.(10) reminds the Akulov-Zener law for magnetic anisotropy, well known in ferromagnetism theory As an example, we consider now a nematic liquid crystal limited by an isotropic substrate. In this case only the polar angle $`\theta `$ enters in the description. The analysis of the temperature surface transitions in a system of this kind is usually performed by means of a Landau’s expansion of the anisotropic part of the surface energy . According to this approach $`W(𝐧)`$ is expanded in power series of the invariants made with the elements of symmetry characterizing the nematic phase (which is the nematic tensor order parameter of elements $`Q_{ij}=(3/2)S[n_in_j(1/3)\delta _{ij}]`$), and the substrate (which is the geometrical normal $`𝐳`$). In the Landau-like approaches the quantity playing the role of expansion parameter is $`S`$. However, since the nematic-isotropic phase transition is first order, $`S`$ is never very small (at the transition point it is of the order of 0.3 ). At the second order in $`S`$, $`W^L(\theta )=w_2^LP_2(\mathrm{cos}\theta )+w_4^LP_4(\mathrm{cos}\theta )+𝒪(3)`$, where $`w_2^L=a_1S+a_2S^2`$ and $`w_4^L=a_3S^2`$, in which $`a_1,a_2`$ and $`a_3`$ are constant parameters, temperature independent . Now we want to compare the prediction of a Landau’s expansion up to the second order in $`S`$ with the result of our mean field analysis. In the case under consideration the angular functions $`L_{2l}(\theta ,\varphi )`$ reduce to $`L_{2l}(\theta )=P_{2l}(\mathrm{cos}\theta )`$, and $$W(\theta )=\underset{l}{}w_{2l}(T)P_{2l}(\mathrm{cos}\theta ).$$ (11) ¿From Eq.(8) we obtain $`w_2(T)/w_2(0)=S`$. This means that at the first order in $`S`$ the temperature dependence of the anchoring energy deduced by means of symmetry considerations, $`w_2^L`$, and by means of the mean field agree. However, for $`l=2`$ there is a discrepancy between the two approaches. In fact, according to the mean field we have $`w_2(T)/w_2(0)=S`$, and $`w_4(T)/w_4(0)=S_4S^2`$, whereas the Landau’s approach predicts the temperature dependencies $`w_2^L(T)=a_1S+a_2S^2`$ and $`w_4^L(T)=a_3S^2`$. More precisely, it predicts a renormalization of the coefficient of $`P_2(\mathrm{cos}\theta )`$, by means of a $`S^2`$ contribution, and a temperature dependence of the coefficient of $`P_4(\mathrm{cos}\theta )`$ like $`S^2`$. Of course, in the limit of small $`S`$ the two predictions agree. In fact, if $`S1`$ the renormalization of $`P_2(\mathrm{cos}\theta )`$ in $`S^2`$ can be neglected with respect to the linear term in $`S`$. Furthermore, in this approximation, $`S_4S^2`$. However, in the case of large $`S`$ the discrepancy between the two approaches can be large. In the low temperature region, where it is possible to use the approximate expressions given by Eq.(10) for the thermal renormalization of the anchoring coefficients, our mean field approach predicts $`w_2(T)S`$ and $`w_4(T)S^{10/3}`$. In Fig.2 we show $`S^2`$, predicted by Landau-like models, and $`S^{10/3}`$, predicted by our mean field theory in the low temperature region, $`\mathrm{𝑣𝑠}`$. $`S`$. As it is evident from this figure, our theory represents an improvement with respect to the Landau-like approaches in the whole temperature range. To conclude we stress the main results reported in the paper. We have shown that the renormalization due to the thermal fluctuations of the anchoring coefficients $`w_{2l}`$ is of the kind $`w_{2l}(T)/w_{2l}(0)=S_{2l}`$ where $`S_{2l}`$ is the $`2l`$-th scalar order parameter. In the particular case in which the nematic phase is described by the Maier-Saupe theory, $`w_{2l}(T)/w_{2l}(0)`$ coincides with the average value of the $`2l`$-th Legendre polynomial. We have also shown that only at the lowest order in the scalar order parameter the simple approach based just on the symmetry of the problem agrees with our mean field approach. This is a consequence of the hypothesis of small $`S`$, over which is based the validity of the Landau-like expansions of $`W(𝐧)`$ in power of $`S`$. We have proposed also approximate expressions for the thermal renormalization of the anchoring coefficients, valid in the low temperature region, where the fluctuations of the molecular directions with respect to the nematic director are small. Acknowledgments A.K.Z. has been partially supported by CNR-NATO by means of a NATO Guest fellowships program. Many thanks are due to S. Faetti, C. Oldano and S. Zumer for useful discussions.
warning/0004/cond-mat0004456.html	ar5iv	text	# Emission Spectrum of a Dipole in a Semi–infinite Periodic Dielectric Structure: Effect of the Boundary ## I Introduction It is well known that fluorescence lifiteme of an atom can be drastically changed when placed in an inhomogenous medium. In his pioneering work Purcell predicted a strong enhancement of the radiative decay rate of an emitter placed inside a resonant microcavity. In contrast, it was also suggested that spontaneous emission can be totally inhibited if the emitter transition frequency lies below that of the fundamental resonator mode. This effect can be understood in terms of redistribution of photonic density of states (DOS) caused by inhomogeneity and/or nontrivial boundary conditions imposed on the radiative field. During the years the spontaneous emission of a dipole coupled to various optical enviroments such as metallic cavities, Fabry Perot two–mirror cavities, dielectric microspheres and nanobubbles, has been a subject of extensive theoretical and experimental studies. Recently there has been a growing interest in studies of radiative properties of fluorescent molecules inside periodic dielectric structures, so called photonic crystals (PC). The Bragg diffraction of light that occurs in these systems opens up a spectral gap (or a pseudo–gap) in the photonic DOS in analogy with electronic energy gaps in semiconductors. A quantum electro–dynamical approach for the radiative decay inside PC has revealed novel physical phenomena such as strong suppression of spontaneous emission within the gap and formation of photon–atom bound states, localization of superradiant modes near the band edges, $`\mathrm{e}\mathrm{t}\mathrm{c}.`$. Calculations of emission spectra within the framework of classical theory were performed for one–dimensional Kronig-Penny type model as well as for three–dimensional fcc lattice structures. It has been established that inhibition of spontaneous emission within the gap is accompanied with strong enhancement at the band edges. It was also shown that the emission spectrum strongly depends on the position of the emiter within the unit cell, as well as on its orientation. Experimental observations of inhibited spontaneous emission have been reported for different periodic structures. So far, the exsisting theoretical studies have considered infinite periodic structures. As a result, the emission power was identically zero within the gap. In the experiments, however, the photonic gap appears as a drop by a factor of $`2`$ in the emission power within a certain spectral interval. This points at the important role of the boundary between the photonic crystal and the air, which is studied in the present paper. As we will show below, accounting for a nearby interface results in a nontrivial dependence of the emission spectrum on the $`\mathrm{i}\mathrm{n}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{a}\mathrm{l}`$ $`\mathrm{p}\mathrm{h}\mathrm{a}\mathrm{s}\mathrm{e}`$ of the dielectric modulation. If the distance of the dipole from the boundary surface plane and the dielectric modulation period are respectively $`d`$ and $`a`$, then one might expect that the dependence of the emission spectrum on the $`\mathrm{i}\mathrm{n}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{a}\mathrm{l}`$ $`\mathrm{p}\mathrm{h}\mathrm{a}\mathrm{s}\mathrm{e}`$ be small in the parameter $`a/d`$. On the contrary, we found that the strong dependence of the emission spectrum on the $`\mathrm{i}\mathrm{n}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{a}\mathrm{l}`$ $`\mathrm{p}\mathrm{h}\mathrm{a}\mathrm{s}\mathrm{e}`$ persists even in the limit $`d/a\mathrm{}`$ (provided there is no absorption in the system). We illustrate this effect in the frame of the simplest possible model. Namely, we choose the dielectric modulation to be $`\mathrm{i}`$) weak, $`\mathrm{i}\mathrm{i}`$) one–dimensional, $`\mathrm{i}\mathrm{i}\mathrm{i}`$) sinusoidal. To quantitavely study the effect of a plane boundary we generalize the standard calculations of the emission rate in periodic media for the case of semi–infinite geometry (Sec. 2). In Sec. 3 we present numerical results for emission spectra illustrating the role of the $`\mathrm{i}\mathrm{n}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{a}\mathrm{l}`$ $`\mathrm{p}\mathrm{h}\mathrm{a}\mathrm{s}\mathrm{e}`$. Discussion of our results and their relevance to recent measurements is presented in Sec. 4. ## II Derivation of the Power Emission Spectrum We schematically depict the system under consideration in Fig. 1. The dielectric function for the left half–space is constant and equals $`\epsilon _0`$, whereas for $`x>0`$ is given by $$\epsilon (x)=\epsilon _1+\delta \epsilon \mathrm{cos}(\sigma x+\varphi ),$$ (1) where $`\varphi `$ is the initial phase of the dielectric modulation, $`\sigma =2\pi /a`$ is the modulation wave vector, and $`\delta \epsilon `$ is the amplitude of the modulation. Below we assume $`\delta \epsilon \epsilon _1`$. The wave equations for the elctric and magnetic fields are $$^2𝐄(𝐫)(𝐄(𝐫))+\frac{\omega ^2}{c^2}\epsilon (x)𝐄(𝐫)=i\frac{4\pi \omega }{c^2}𝐉(𝐫)$$ (2) $$^2𝐁(𝐫)[\mathrm{ln}\epsilon (x)]\times \times 𝐁(𝐫)+\frac{\omega ^2}{c^2}\epsilon (x)𝐁(𝐫)=\frac{4\pi }{c}\times 𝐉(𝐫)$$ (3) where the radiation source $`𝐉(𝐫)=𝐣\delta (𝐫𝐫_\mathrm{𝟎})`$ is located at point $`𝐫_\mathrm{𝟎}=(d,0,0)`$. We note that the term proportional to $`\delta \epsilon `$ has been neglected in the r.h.s. of Eq. (3). The time averaged radiative power per unit solid angle is given by $$\frac{dP}{d\mathrm{\Omega }}=\frac{c}{8\pi }Re\left[r^2𝐧(𝐄\times 𝐁^{})\right],$$ (4) where $`B^{}`$ is the complex conjugate of $`B`$, $`r=\|𝐫\|`$ and $`𝐧=𝐫/r`$ is the unit radius vector. Without any loss of generality we choose $`n_z=0`$ (see Fig. 1). Then it is very convenient to separately treat two possible orientations of the dipole. Indeed, one can easily check that the current density components $`J_z`$ and $`J_x,J_y`$ give rise to Electro–Magnetic ($`EM`$) radiation with respectively electric ($`TE`$ polarization) and magnetic ($`TM`$ polarization) fields polarized in the $`z`$ direction. Since these two modes do not interfere, their contributions to the radiation power are additive. The corresponding $`EM`$ wave equations for the two polarizations are obtained from the $`z`$–components of Eqs. (2) and (3) by taking the Fourier transform with respect to $`y`$ and $`z`$ coordinates: $$\frac{d^2E_z}{dx^2}+\left(\frac{\omega ^2}{c^2}\epsilon (x)k_y^2\right)E_z=\frac{4\pi i\omega }{c^2}j_z\delta (xd),(TE)$$ (5) $$\frac{d^2B_z}{dx^2}\left(\frac{\mathrm{ln}\epsilon }{x}\right)\frac{dB_z}{dx}+\left(\frac{\omega ^2}{c^2}\epsilon (x)k_y^2\right)B_z=\frac{4\pi }{c}\left(\frac{j_y}{x}ik_yj_x\right)\delta (xd),(TM)$$ (6) where $`k_y`$ is the $`y`$–component of the wave vector ($`E_z(x;k_y),B_z(x;k_y)e^{ik_yy}`$). Since we want to calculate the power emmited in $`xy`$ plane we have set $`k_z=0`$ in Eqs. (5) and ($`\text{6})`$. The solution of the corresponding homogenous equations may be written as a sum of incident, reflected and transmitted $`EM`$ waves, with two linearly independent terms $`E_1(x)`$, $`E_2(x)`$ and $`B_1(x)`$, $`B_2(x)`$ corresponding to the incident $`EM`$ wave from the right and left, respectively (see inset of Fig. 1). To solve Eqs. (56) we employ the variation of a constant method. We seek solution in the form: $$E_z(x)=C_1^E(x)E_1(x)+C_2^E(x)E_2(x),B_z(x)=C_1^B(x)B_1(x)+C_2^B(x)B_2(x)$$ (7) Upon substituting (7) into (56) we find for the variational coefficients $$C_{1,2}^E(x)=i\omega \frac{4\pi }{c^2}W_E^1_{X_{1,2}^E}^x𝑑x^{}E_{2,1}(x^{})j_z\delta (x^{}d)$$ (8) $$C_{1,2}^B(x)=\frac{4\pi }{c}W_B^1_{X_{1,2}^B}^x𝑑x^{}B_{2,1}(x^{})\left(\frac{j_y}{x^{}}ik_yj_x\right)\delta (x^{}d)$$ (9) where $`W_E`$, $`W_B`$ are the Wronskians $$W_E=E_1\frac{dE_2}{dx}E_2\frac{dE_1}{dx},W_B=B_1\frac{dB_2}{dx}B_2\frac{dB_1}{dx}$$ (10) and $`X_{1,2}^E`$, $`X_{1,2}^B`$ in the lower integration limits are constants of integration. They are determined from the boundary conditions that there are no incoming $`EM`$ waves, implying $`X_1^E=X_1^B=\mathrm{}`$, $`X_2^E=X_2^B=\mathrm{}`$. Hence, the solutions of Eqs. (56) for large negative $`x`$ which satisfy the boundary conditions can be written as follows $$E_z(x)=i\omega \frac{4\pi }{c^2}W_E^1j_zE_1(d)E_2(x),$$ (11) $$B_z(x)=\frac{4\pi }{c}W_B^1(j_y\frac{dB_1}{dx}\|_{x=d}+ik_yj_xB_1(d))B_2(x).$$ (12) It can be easily shown that for negative $`x`$ one has $`W_E^1E_2(x)=W_B^1B_2(x)=ie^{ik_xx}/2k_x`$. Thus, the remaining task is to find $`E_1(d)`$, $`B_1(d)`$ and $`(dB_1/dx)_{x=d}`$. The solution of the homogeneous equations for the left half–space can be written as a sum of two plane waves, $$E_1(x)=e^{ik_xx}+R_Ee^{ik_xx},B_1(x)=e^{ik_xx}+R_Be^{ik_xx}$$ (13) where $`R_E`$, $`R_B`$ are the optical reflection coefficients for $`TE`$ and $`TM`$ polarizations, respectively. For the right half–space one may use the Bloch theorem to find the solution $$E_1(x)=e^{iq_Ex}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}A_n^Ee^{i\sigma nx},B_1(x)=e^{iq_Bx}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}A_n^Be^{i\sigma nx}.$$ (14) When substituting Eq.(14) into Eqs.(56) we obtain two infinite systems of linear, homogenous equations for the coefficents $`A_n^E`$, $`A_n^B`$ $$\left(\frac{\omega ^2}{c^2}\epsilon _1k_y^2(q_E+n\sigma )^2\right)A_n^E+\frac{\omega ^2}{c^2}\frac{\delta \epsilon }{2}\left(e^{i\varphi }A_{n1}^E+e^{i\varphi }A_{n+1}^E\right)=0$$ (15) $`\left({\displaystyle \frac{\omega ^2}{c^2}}\epsilon _1k_y^2(q_B+n\sigma )^2\right)A_n^B+{\displaystyle \frac{\delta \epsilon }{2\epsilon _1}}e^{i\varphi }\left({\displaystyle \frac{\omega ^2}{c^2}}\epsilon _1+\sigma (q_B+(n1)\sigma )\right)A_{n1}^B`$ (16) $`+{\displaystyle \frac{\delta \epsilon }{2\epsilon _1}}e^{i\varphi }\left({\displaystyle \frac{\omega ^2}{c^2}}\epsilon _1\sigma (q_B+(n+1)\sigma )\right)A_{n+1}^B=0.`$ (17) We note that the $`\mathrm{i}\mathrm{n}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{a}\mathrm{l}`$ $`\mathrm{p}\mathrm{h}\mathrm{a}\mathrm{s}\mathrm{e}`$ $`\varphi `$ explicitly enters into these equations. Near the Bragg resonance that occurs at $`q_{E,B}\sigma /2`$, the main coefficents which contribute to the sums in Eq.(14) are $`A_0^E`$, $`A_0^B`$ and $`A_1^E`$, $`A_1^B`$ since the rest of $`A_n^{E,B}`$ are small in the parameter $`\delta \epsilon /\epsilon _1`$. In this approximation the equation systems (1516) may be simplified into two $`2\times 2`$ matrix equations. Requiring for the determinants of these matrices to vanish, one finds the dispersion relations for the two $`EM`$ polarizations near resonance $$\delta q_E=\frac{\sigma }{2\kappa }\sqrt{\left(\frac{\delta \omega }{\omega _0}\right)^2\mathrm{\Delta }^2}$$ (18) $$\delta q_B=\frac{\sigma }{2\kappa }\sqrt{\left(\frac{\delta \omega }{\omega _0}\right)^2\mathrm{\Delta }^2(12\kappa )^2},$$ (19) where we have introduced $`\mathrm{\Delta }=\delta \epsilon /4\epsilon _1`$, $`q_{E,B}\sigma /2=\delta q_{E,B}\sigma /2`$, $`\omega \omega _0=\delta \omega \omega _0`$, $`\kappa =\sigma ^2c^2/4\omega _0^2ϵ_1`$ and the $`k_y`$–dependent resonant frequency $`\omega _0`$ is given by $$\omega _0=\frac{c}{\epsilon _1}\sqrt{\sigma ^2/4+k_y^2}.$$ (20) Eqs. (1819) show that there is a spectral gap for $`EM`$ waves propagating in the system centered at $`\omega _0`$. For both $`TE`$ and $`TM`$ polarizations the central gap position shifts to higher frequencies with increasing $`k_y`$ (which also determines the propagation direction of the radiative field). For $`TE`$ polarization the gap broadens with increasing angle $`\theta ^{}`$ (see Fig. 1) whereas for $`TM`$ polarization the gap narrows and disappears at $`2\kappa =1`$, which corressponds to the propagation direction for which $`k_y=\sigma /2`$. This can be defined as a Brewster angle for Bragg diffraction. If one increases $`k_y`$ further then the gap reopens again. Using Eqs. (15)–(19) we obtain the following expressions for the electric and magnetic fields for the right half space: $$E_1(x)=A_1^Ee^{i\delta q_Ex}\left(e^{i\frac{\sigma }{2}x}F_Ee^{i\varphi }e^{i\frac{\sigma }{2}x}\right)$$ (21) $$B_1(x)=A_1^Be^{i\delta q_Bx}\left(e^{i\frac{\sigma }{2}x}F_Be^{i\varphi }e^{i\frac{\sigma }{2}x}\right)$$ (22) Here the functions $`F_E`$ and $`F_B`$ describe the coupling between incident and Bragg reflected waves and are defined as follows: $$F_E(\delta \omega ,\theta )=\frac{1}{\mathrm{\Delta }}\left(\frac{\delta \omega }{\omega _0}\sqrt{\left(\frac{\delta \omega }{\omega _0}\right)^2\mathrm{\Delta }^2}\right)$$ (23) $$F_B(\delta \omega ,\theta )=\frac{1}{\mathrm{\Delta }(12\kappa )}\left(\frac{\delta \omega }{\omega _0}\sqrt{\left(\frac{\delta \omega }{\omega _0}\right)^2\mathrm{\Delta }^2(12\kappa )^2}\right).$$ (24) Finally, we match the fields components and their derivatives at $`x=0`$, take the inverse Fourier transform with respect to $`k_y,k_z`$ and use Eq. (4) to get the following expression for the total radiated power $`dP/d\mathrm{\Omega }`$, normalized to the radiation power when there is no dielectric modulation: $`\left({\displaystyle \frac{dP}{d\mathrm{\Omega }}}\right)_N={\displaystyle \frac{j_z^2T_{0E}^2}{j_z^2T_{0E}^2+\epsilon _1j_1^2T_{0B}^2}}\left\|{\displaystyle \frac{1F_Ee^{i(\sigma d+\varphi )}}{1R_{0E}F_Ee^{i\varphi }}}e^{i\delta q_Ed}\right\|^2`$ (25) $`+{\displaystyle \frac{\epsilon _1j_1^2T_{0B}^2}{j_1^2T_{0E}^2+\epsilon _1j_1^2T_{0B}^2}}\left\|{\displaystyle \frac{1\chi F_Be^{i(\sigma d+\varphi )}}{1R_{0B}F_Be^{i\varphi }}}e^{i\delta q_Bd}\right\|^2.`$ (26) Here $`T_{0E},R_{0E}`$ and $`T_{0B},R_{0B}`$ are the conventional Fresnel’s transmission and reflection coefficients from the dielectric interface without the dielectric modulation for the amplitudes of $`TE`$ and $`TM`$ polarized waves, respectively, $`j_1=j_y\mathrm{cos}\theta ^{}+j_x\mathrm{sin}\theta ^{}`$, where $`\theta ^{}`$ is related to the observation angle $`\theta `$ by Snell’s law. The quantity $`\chi `$, which depends on the dipole orientation in the $`xy`$ plane is defined as $`\chi =(j_yj_x\mathrm{tan}\theta ^{})/(j_y+j_x\mathrm{tan}\theta ^{})`$. Eq. (25) is the main result of this paper. We note that the phase $`\varphi `$ is explicitly present in both the numerator and denominator of Eq. (25), indicating the important role of the $`\mathrm{i}\mathrm{n}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{a}\mathrm{l}`$ $`\mathrm{m}\mathrm{o}\mathrm{d}\mathrm{u}\mathrm{l}\mathrm{a}\mathrm{t}\mathrm{i}\mathrm{o}\mathrm{n}`$ $`\mathrm{p}\mathrm{h}\mathrm{a}\mathrm{s}\mathrm{e}`$ at the boundary interface. In the next section we numerically analyze the role of $`\varphi `$ in the emission spectrum for different situations. ## III Numerical Results Before presenting our numerical results let us concentrate on a particular experimental case of opals and opal replica. These artificial photonic crystals consist of closely packed $`SiO_2`$ spheres forming an fcc structure that contains fully interconnected voids. The opal replica is formed by filling these voids with a precursor polymer solution and then etching the $`SiO_2`$ spheres after polymerization. The opals and opal replica PC are infiltrated with various fluorescent dye solutions to provide a radiation source inside the crystal. Inhibited spontaneous emission of the dye molecules in such PC have been recently studied by several groups. The refractive index $`n`$ of $`SiO_2`$ is $`n1.46`$ and the refractive index contrast $`\mathrm{\Delta }n`$ between the $`SiO_2`$ balls and the dye solution ranges between $`\mathrm{\Delta }n0.10.3`$. This is not sufficient for a formation of a $`\mathrm{c}\mathrm{o}\mathrm{m}\mathrm{p}\mathrm{l}\mathrm{e}\mathrm{t}\mathrm{e}`$ photonic band gap. Instead, the system posseses pseudogaps (or partial gaps) with an angle–dependent central frequency. To compare our results to the experiments we have chosen in our model $`\mathrm{\Delta }=\delta \epsilon /4\epsilon _1=0.1`$. First, we consider the case when the emitter is many periods away from the interface ($`d=5a`$). In Fig. 2 we show the emission power $`\mathrm{a}\mathrm{v}\mathrm{e}\mathrm{r}\mathrm{a}\mathrm{g}\mathrm{e}\mathrm{d}`$ over the orientations of the emitter $`\mathrm{a}\mathrm{s}`$ $`\mathrm{w}\mathrm{e}\mathrm{l}\mathrm{l}`$ $`\mathrm{a}\mathrm{s}`$ $`\mathrm{o}\mathrm{v}\mathrm{e}\mathrm{r}`$ $`\mathrm{i}\mathrm{t}\mathrm{s}`$ $`\mathrm{p}\mathrm{o}\mathrm{s}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{o}\mathrm{n}`$ within the unit cell (the zero on the frequency axis in all plots corresponds to the central gap frequency $`\mathrm{\Omega }_0`$ at $`\theta =0`$, where $`\mathrm{\Omega }_0=\sigma c/(2\sqrt{\epsilon _1})`$). Inside the gap, the emission power is strongly suppressed. It is seen, however, that even in the limit of large $`d/a`$ the features of the $`\mathrm{a}\mathrm{v}\mathrm{e}\mathrm{r}\mathrm{a}\mathrm{g}\mathrm{e}\mathrm{d}`$ emission power still depend on the $`\mathrm{i}\mathrm{n}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{a}\mathrm{l}`$ $`\mathrm{p}\mathrm{h}\mathrm{a}\mathrm{s}\mathrm{e}`$. The evolution of $`(dP/d\mathrm{\Omega })_N`$ with $`\varphi `$ outside the gap is described as follows. At $`\varphi =0`$ there is a well pronounced singularity at the lower edge of the spectral gap. With increasing $`\varphi `$, this singularity diminishes whereas another singularity starts to develop at the upper edge; the spectrum becomes symmetric at $`\varphi =\pi /2`$. Further increase of $`\varphi `$ leads to a gradual transformation of the initial curve into its mirror image with respect to the central gap frequency, the singularity now occuring at the upper band edge. Fig. 2 corresponds to $`\theta =0`$. The sensitivity of the emission power to $`\varphi `$ appears to be even more pronounced for $`\theta 0`$. We illustrate this effect by plotting in Fig. 3 the $`\mathrm{a}\mathrm{v}\mathrm{e}\mathrm{r}\mathrm{a}\mathrm{g}\mathrm{e}\mathrm{d}`$ $`(dP/d\mathrm{\Omega })_N`$ for $`\theta =60^o`$ for the set of $`\mathrm{i}\mathrm{n}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{a}\mathrm{l}`$ $`\mathrm{p}\mathrm{h}\mathrm{a}\mathrm{s}\mathrm{e}\mathrm{s}`$ $`\varphi =0,\pi /4,\pi /2`$. With increasing $`\varphi `$ again, it is seen that there is a tendency for the emission spectrum to become symmetric near $`\varphi =\pi /2`$. At $`\varphi =\pi `$ (not presented here) the emission spectrum is again transformed into the mirror image of the initial curve at $`\varphi =0`$ but now with respect to the shifted central gap frequency $`\mathrm{\Omega }_0/cos\theta ^{}`$. We also note the appearance of additional singularities in the emission spectrum. Their origin lies in the different angle–dependenceies of the gap width for $`TE`$ and $`TM`$ polarizations, as seen before in Eqs (18) and (19). To be more specific, we note that the highest and lowest frequency peaks correspond to the band edges for $`TE`$ polarization. Similarly, two peaks at the intermidiate frequencies determine the band edges for $`TM`$ polarization. Both pairs of the singularities are located symetrically around the shifted Bragg frequency $`\omega _0=\mathrm{\Omega }_0/\mathrm{cos}\theta ^{}`$. Let us now turn to the discussion of the case where the emitter is close to the interface. The main feature of this situation is that $`(dP/d\mathrm{\Omega })_N`$ $`\mathrm{i}\mathrm{n}\mathrm{s}\mathrm{i}\mathrm{d}\mathrm{e}`$ the gap remains finite. One can see from Eq. (25) that moving the emitter $`N`$ periods away from the boundary decreases the emission power at the center of the gap by a factor of $`exp(2\pi N\mathrm{\Delta })`$. Therefore, to study the features of the spontaneous emission inside the gap we choose $`N=1,2`$. This situation is illustrated in Fig. 4 where we plot the averaged $`(dP/d\mathrm{\Omega })_N`$ for the observation angle $`\theta =60^o`$. One can see that the evolution of the emission power with increasing the $`\mathrm{i}\mathrm{n}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{a}\mathrm{l}`$ $`\mathrm{p}\mathrm{h}\mathrm{a}\mathrm{s}\mathrm{e}`$ is very strong. Again, we note that there are four well pronounced singularities in the emission spectrum that corresponds to the spectral gap edges of the $`TE`$ and $`TM`$ polarizations. In particular, for $`\varphi =0`$ there is a noticable enhancement of the emission rate (by a factor of $`2`$) at the frequency which determines the lower band–edge of the $`TE`$ polarization, whereas for $`\varphi =\pi /2`$ a similar enhancement occurs at the edges of the $`TM`$ polarization gap. Finally, in Fig. 5 we plot the emission power integrated over the observation angle $`\theta `$ (i.e., the total power emitter in $`xy`$ plane) for $`\varphi =0`$ and $`\varphi =\pi `$. In this case too an averaging over the dipole position within the unit cell has been performed. To allow for an unpolarized emission we have chosen $`j_x=j_y=j_z`$. Remarkably, even after angular averaging a weak dependence of the emission power on the initial phase still persists. ## IV Discussion We have studied the emission spectrum of a dipole inside a one–dimensional periodic structure in the presence of a nearby plane boundary. As expected, the emission rates are strongly suppressed for frequencies inside the spectral gap, provided that the emitter is many periods away from the interface ($`d/a>5`$). For frequencies near the band edges, the emission spectrum changes drastically with the $`\mathrm{i}\mathrm{n}\mathrm{i}\mathrm{t}\mathrm{i}\mathrm{a}\mathrm{l}`$ $`\mathrm{p}\mathrm{h}\mathrm{a}\mathrm{s}\mathrm{e}`$ of the dielectric modulation. We also observed enhancement of the emission rates near the band edges, however, by a factor much smaller than predicted in the previous studies. This can be attributed to the fact that the modulation in our model is weak. In Ref. , where the total radiation power from a dipole inside a infinite $`3D`$ fcc lattice was calculated numerically, the contrast in the dielectric constant $`\mathrm{\Delta }\epsilon `$ was more than $`10`$ allowing for the formation of a complete photonic band gap. This resulted in an enhancement of the radiated power at the band edges by a factor of $`25`$. Similarly, in Ref. where the authors considered $`1D`$ Kronig–Penney type modulation of the refractive index, the enhancement factor at the band edges was about $`30`$, whereas inside the gap it was identically zero. Apparently, this resulted from consideration of radiative modes polarized parallel to the dipole direction. Allowing for nonpolarized radiation ($`\mathrm{i}.\mathrm{e}.`$ in all directions) should lead to qualitatively different results (see for example Fig. 5). Strictly speaking, the straightforward comparision of our calculation with those of previous studies is not possible due to the different approach developed here. However, it is clear from our consideration that the sensitivity of the power emission to the boundary conditions should persist also for strong and non–sinusoidal modulation and become even stronger. In this paper we also studied the case where the emitter is sufficiently close to the interface, so that the emission power for the frequencies inside the gap is finite. In this case also we showed that the features of the emission spectrum are very sensitive to the initial phase of the periodic modulation, hence emphasizing the effect of the boundary. We note that in our calculations we assumed that no deffects exist in the system. In the presence of weak disorder the emission spectrum would be significantly modified. Let us consider radiation from a dipole that is many periods away from the boundary with frequency inside the spectral gap. If there are no defects then the radiation ($`e.`$ $`\mathrm{g}.`$, in the direction normal to the boundary) is strongly attenuated (see Fig. 2). However, introducing a small concentration of defects opens up a new mechanism for the light emission to come out from the PC in the direction $`\mathrm{n}\mathrm{o}\mathrm{r}\mathrm{m}\mathrm{a}\mathrm{l}`$ to the boundary. Namely, light can propagate without any attenuation in directions for which the Bragg condition is not satisfied and then scatter off defects that are close to the interface. As a result, the emission power in the direction normal to the boundary is not exponentially small but is finite, similar to the case of a dipole close to the boundary. Remarkably, only defects close to the boundary, within the Bragg atenuation lenght $`\xi _B=(2\pi \sigma \mathrm{\Delta })^1`$ from the interface, contribute to the emission power. Owing to this effect, we conclude that the emission spectrum in the presence of a weak disorder is rather universal, since it does not depend on the dipole–interface distance $`d`$. Experimentally this was demonstrated in Ref. , where the authors found a way to excite fluorescent dye molecules at different distances from the boundary. They have demonstrated that the emission power at the center of the spectral gap is attenuated by a factor of $`2`$ and does not change upon changing $`d`$. This was interpretated in terms of light scattering off defects close to the PC interface. ## V Acknowledgements The authors thank N. Eradat, Drs. A. A. Zakhidov and R. H. Baughman for useful discussions. This work was supported in part by NSF under grant DMR–9732820 and the International Research Foundation of NEDO (Japan). One of the authors (M. E. R.) acknowledges the support of Petroleum Research Fund under grant ACS–PRF #34302–AC6.
warning/0004/cond-mat0004452.html	ar5iv	text	# Optical Absorption Spectra of Bipolarons ## I Introduction The optical and kinetic properties of polar and ionic solids are substantially influenced by the polaron coupling. Large polarons have been most clearly manifested by investigations of the transport phenomena in a magnetic field (see the recent review on polarons ). When two electrons (or two holes) interact with each other simultaneously through the Coulomb force and via the electron-phonon interaction, either two independent polarons can occur or a bound state of two polarons — the bipolaron — can arise (see Refs. on large bipolarons and a comprehensive review concerning small bipolarons). Whether bipolarons are stable or not, depends on the competition between the repulsive forces (the Coulomb interaction) and the attractive forces (mediated through the electron-phonon interaction). Verbist et al. analyzed the large bipolaron using the Feynman path-integral formalism . They introduced a “phase diagram” for the polaron—bipolaron system on the basis of a generalization of Feynman’s trial action and showed that the Pekar-Fröhlich coupling constant as high as 6.8 is needed in 3D to allow for bipolaron formation. Furthermore, in Refs. it was shown that the large bipolaron formation is facilitated in 2D as compared to 3D. Experimental evidences for bipolarons, e. g. from the data on magnetization and electric conductivity in Ti<sub>4</sub>O<sub>7</sub>, as well as in Na<sub>0.3</sub>V<sub>2</sub>O<sub>5</sub> and polyacetylene, was discussed by Mott . In the framework of the renewed interest in bipolaron theory , a preliminary analysis of the absorption of large bipolarons without a magnetic field was given . A variational treatment of both spin-singlet and spin-triplet states of large bipolarons (in 2D) and for sufficiently strong coupling in high magnetic fields has been presented in Ref. . In Ref. , the path-integral approach for a bipolaron was generalized to the case of a bipolaron in a magnetic field for all values of the Pekar-Fröhlich coupling constant and for all magnetic field strengths. It was demonstrated, that the magnetic field favors bipolaron formation. The first investigations of the optical absorption spectrum of large bipolarons in a magnetic field were performed in Refs. . ## II Approach We investigate here the optical properties of a bipolaron using the path-integral memory-function technique, developed in Refs. (see also Ref. ). The path-integral variational principle gives excellent results for the free energy of the electron-phonon systems (see, e. g., Refs. ). Hence, these models adequately describe the physical properties (including the energy spectra) of the polaron and of the bipolaron, respectively. The Hamiltonian of two electrons, interacting with longitudinal optical (LO) phonons and between each other in an external electric field $`𝐄\left(t\right)`$, is $`\widehat{H}`$ $`=`$ $`{\displaystyle \underset{j=1,2}{}}\left[{\displaystyle \frac{𝐩_j^2}{2m}}+e𝐫_j\left(t\right)𝐄\left(t\right)\right]+{\displaystyle \frac{e^2}{\epsilon _{\mathrm{}}\left\|𝐫_1𝐫_2\right\|}}+{\displaystyle \underset{𝐪}{}}\mathrm{}\omega _{\mathrm{LO}}\left(\widehat{a}_𝐪^+\widehat{a}_𝐪+{\displaystyle \frac{1}{2}}\right)`$ (2) $`+{\displaystyle \underset{j=1,2}{}}{\displaystyle \underset{𝐪}{}}\left(V_𝐪\widehat{a}_𝐪e^{\mathrm{i}𝐪𝐫_j}+V_𝐪^{}\widehat{a}_𝐪^+e^{\mathrm{i}𝐪𝐫_j}\right)`$ with the interaction amplitudes $$V_𝐪=\frac{\mathrm{}\omega _{\mathrm{LO}}}{\mathrm{i}q}\left(\frac{2\sqrt{2}\pi \alpha }{V}\right)^{1/2}\left(\frac{\mathrm{}}{m\omega _{\mathrm{LO}}}\right)^{1/4},$$ (3) where $`\alpha `$ is the electron-phonon coupling constant, $`m`$ is the electron band mass, $`V`$ is the volume of the crystal, $`\epsilon _{\mathrm{}}`$ is the high-frequency dielectric constant, $`\omega _{\mathrm{LO}}`$ is the LO-phonon frequency. We consider the equation of motion for the vector function $`𝐑\left(t\right)`$ which has the sense of the average coordinate response per one electron, $$𝐑\left(t\right)\frac{1}{2}\underset{j=1,2}{}𝐫_j\left(t\right)_S.$$ (4) Here, the symbol of averaging denotes the path-integral average over trajectories of the electrons (cf. Ref. ) $$F[\overline{𝐫}\left(t\right),\overline{𝐫}^{}\left(t\right)]_S\frac{D\overline{𝐫}\left(t\right)D\overline{𝐫}^{}\left(t\right)F[\overline{𝐫}\left(t\right),\overline{𝐫}^{}\left(t\right)]\mathrm{exp}\left\{\left(\mathrm{i}/\mathrm{}\right)S[\overline{𝐫}\left(t\right),\overline{𝐫}^{}\left(t\right)]\right\}}{D\overline{𝐫}\left(t\right)D\overline{𝐫}^{}\left(t\right)\mathrm{exp}\left\{\left(\mathrm{i}/\mathrm{}\right)S[\overline{𝐫}\left(t\right),\overline{𝐫}^{}\left(t\right)]\right\}}\left[\overline{𝐫}(𝐫_1,𝐫_2)\right]$$ (5) with the action functional $$S[\overline{𝐫}\left(t\right),\overline{𝐫}^{}\left(t\right)]=S_e\left[\overline{𝐫}\left(t\right)\right]S_e\left[\overline{𝐫}^{}\left(t\right)\right]\mathrm{i}\mathrm{}\mathrm{\Phi }[\overline{𝐫}\left(t\right),\overline{𝐫}^{}\left(t\right)],$$ (6) where $`S_e[𝐫_1\left(t\right),𝐫_2\left(t\right)]`$ is the action of two interacting electrons in an external electric field, $$S_e\left[\overline{𝐫}\left(t\right)\right]=_{\mathrm{}}^{\mathrm{}}\left\{\underset{j=1,2}{}\left[\frac{m\dot{𝐫}_j^2\left(t\right)}{2}e𝐫_j\left(t\right)𝐄\left(t\right)\right]\frac{e^2}{\epsilon _{\mathrm{}}\left\|𝐫_1\left(t\right)𝐫_2\left(t\right)\right\|}\right\}𝑑t,$$ (7) and $`\mathrm{\Phi }[\overline{𝐫}\left(t\right),\overline{𝐫}^{}\left(t\right)]`$ is the “influence phase” of the phonon subsystem $`\mathrm{\Phi }[\overline{𝐫}\left(t\right),\overline{𝐫}^{}\left(t\right)]`$ $`=`$ $`{\displaystyle \underset{𝐪}{}}{\displaystyle \frac{\left\|V_𝐪\right\|^2}{\mathrm{}^2}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑t{\displaystyle \underset{\mathrm{}}{\overset{t}{}}}𝑑t^{}\left[\rho _𝐪\left(t\right)\rho _𝐪^{}\left(t\right)\right]`$ (9) $`\times \left[T_{\omega _𝐪}^{}\left(tt^{}\right)\rho _𝐪\left(t^{}\right)T_{\omega _𝐪}\left(tt^{}\right)\rho _𝐪^{}\left(t^{}\right)\right].`$ This influence phase describes both a retarded interaction between different electrons and a retarded self-interaction of each electron due to the elimination of the phonon coordinates. It contains the Fourier components of the electron density operator $$\rho _𝐪=\underset{j=1,2}{}\mathrm{e}^{\mathrm{i}\mathrm{𝐪𝐫}_j}$$ (10) and the phonon Green’s function $$T_\omega \left(t\right)=\frac{e^{\mathrm{i}\omega t}}{1e^{\beta \omega }}+\frac{e^{\mathrm{i}\omega t}}{e^{\beta \omega }1},\beta \frac{\mathrm{}}{k_BT},$$ (11) $`T`$ is the temperature. We have calculated the optical absorption coefficient for a bipolaron in the memory-function approach , $$\mathrm{\Gamma }\left(\omega \right)=\frac{4\pi }{cn}\frac{n_0e^2}{m}\frac{\omega \text{Im}T\left(\omega \right)}{\left[\omega ^2\text{Re}T\left(\omega \right)\right]^2+\left[\text{Im}T\left(\omega \right)\right]^2},$$ (12) where $`c`$ is the velocity of light, $`n`$ is the refractive index of the crystal, $`n_0`$ is the electron density. The memory function $`T\left(\omega \right)`$ has the form : $$T\left(\omega \right)=\underset{𝐪}{}\frac{\left\|V_𝐪\right\|^2q^2}{3\mathrm{}}\underset{0}{\overset{\mathrm{}}{}}\left(\mathrm{e}^{\mathrm{i}\omega t}1\right)\text{Im}\left[T_{\omega _{\mathrm{LO}}}^{}\left(t\right)\rho _𝐪\left(t\right)\rho _𝐪\left(0\right)_{S_0}\right]𝑑t.$$ (13) It is expressed in terms of the two-point correlation function $`\rho _𝐪\left(t\right)\rho _𝐪\left(0\right)_{S_0}`$ of the electron density operators. The correlation function is calculated as the path-integral average \[cf. Eq. (5)\] with the model action functional $`S_0[\overline{𝐫}\left(t\right),\overline{𝐫}_f\left(t\right),\overline{𝐫}^{}\left(t\right),\overline{𝐫}_f^{}\left(t\right)]`$: $$F_{S_0}\frac{D\overline{𝐫}\left(t\right)D\overline{𝐫}^{}\left(t\right)D\overline{𝐫}_f\left(t\right)D\overline{𝐫}_f^{}\left(t\right)F\mathrm{exp}\left\{\left(\mathrm{i}/\mathrm{}\right)S_0[\overline{𝐫}\left(t\right),\overline{𝐫}_f\left(t\right),\overline{𝐫}^{}\left(t\right),\overline{𝐫}_f^{}\left(t\right)]\right\}}{D\overline{𝐫}\left(t\right)D\overline{𝐫}^{}\left(t\right)D\overline{𝐫}_f\left(t\right)D\overline{𝐫}_f^{}\left(t\right)\mathrm{exp}\left\{\left(\mathrm{i}/\mathrm{}\right)S_0[\overline{𝐫}\left(t\right),\overline{𝐫}_f\left(t\right),\overline{𝐫}^{}\left(t\right),\overline{𝐫}_f^{}\left(t\right)]\right\}}.$$ (14) The model system consists of two electrons harmonically interacting with two fictitious particles and between each other. The model action functional $`S_0[\overline{𝐫}\left(t\right),\overline{𝐫}_f\left(t\right),\overline{𝐫}^{}\left(t\right),\overline{𝐫}_f^{}\left(t\right)]`$ has the form $$S_0[\overline{𝐫}\left(t\right),\overline{𝐫}_f\left(t\right),\overline{𝐫}^{}\left(t\right),\overline{𝐫}_f^{}\left(t\right)]=_{\mathrm{}}^{\mathrm{}}\left[L_0(\dot{\overline{𝐫}},\dot{\overline{𝐫}}_f,\overline{𝐫},\overline{𝐫}_f)L_0(\dot{\overline{𝐫}}^{},\dot{\overline{𝐫}}_f^{},\overline{𝐫}^{},\overline{𝐫}_f^{})\right]𝑑t,$$ (15) with the Lagrangian $`L_0(\dot{\overline{𝐫}},\dot{\overline{𝐫}}_f,\overline{𝐫},\overline{𝐫}_f)`$ $`=`$ $`{\displaystyle \underset{j=1,2}{}}\left({\displaystyle \frac{m\dot{𝐫}_j^2\left(t\right)}{2}}+{\displaystyle \frac{M\dot{𝐫}_{fj}^2\left(t\right)}{2}}\right){\displaystyle \frac{k}{2}}{\displaystyle \underset{j=1,2}{}}\left(𝐫_j𝐫_{fj}\right)^2`$ (17) $`{\displaystyle \frac{k^{}}{2}}\left[\left(𝐫_1𝐫_{f2}\right)^2+\left(𝐫_2𝐫_{f1}\right)^2\right]+{\displaystyle \frac{K}{2}}\left(𝐫_1𝐫_2\right)^2,`$ where $`M`$ is the mass of a fictitious particle, $`k,`$ $`k^{}`$ and $`K`$ are the elastic constants. The oscillator potentials in the Lagrangian (17) imitate the electron-phonon interaction and the electron-electron Coulomb repulsion. In Refs. , the aforesaid model has been first introduced in order to calculate the bipolaron ground-state energy by the Jensen-Feynman variational method . The variational functional of Refs. for the bipolaron free energy contains four variational parameters: $`M,`$ $`k,`$ $`k^{}`$ and $`K`$. It is convenient to use instead of them four other independent parameters: (i) three eigenfrequencies of the internal bipolaron motion $`\mathrm{\Omega }_i`$ $`\left(i=1,2,3\right)`$ (see Ref. ), $`\mathrm{\Omega }_1`$ $`=`$ $`{\displaystyle \frac{m+M}{mM}}\left(k+k^{}\right),`$ (18) $`\mathrm{\Omega }_{2,3}^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{{\displaystyle \frac{m+M}{mM}}\left(k+k^{}\right){\displaystyle \frac{2K}{m}}\pm \left[\left({\displaystyle \frac{Mm}{mM}}\left(k+k^{}\right){\displaystyle \frac{2K}{m}}\right)^2+{\displaystyle \frac{4}{mM}}\left(kk^{}\right)^2\right]^{1/2}\right\},`$ (19) and (ii) the frequency $$v=\left(\frac{k+k^{}}{M}\right)^{1/2},$$ (20) which is analogous to the Feynman parameter $`w`$ in the single-polaron problem . It is seen from Eqs. (18) to (20), that the inequality $`\mathrm{\Omega }_1\mathrm{\Omega }_2v\mathrm{\Omega }_3`$ is fulfilled for the variational frequency parameters. By the variational procedure of Ref. , the optimal values of those variational parameters are found for the physical system of two electrons interacting with the phonon field and between each other. In the zero-temperature limit $`T=0,`$ we have derived from Eq. (13) the analytical expressions for the real and imaginary parts of the memory function for a bipolaron in three dimensions: $`\text{Re}T\left(\omega \right)={\displaystyle \frac{4\sqrt{2}\alpha }{3}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left[1+\left(1\right)^{k+l}\right]}{n!k!l!}}s_1^ns_2^ks_3^l`$ (21) $`\times `$ $`[f_{n+k+l}(\omega 1n\mathrm{\Omega }_1k\mathrm{\Omega }_2l\mathrm{\Omega }_3)+f_{n+k+l}(\omega 1n\mathrm{\Omega }_1k\mathrm{\Omega }_2l\mathrm{\Omega }_3)`$ (22) $``$ $`2f_{n+k+l}(1n\mathrm{\Omega }_1k\mathrm{\Omega }_2l\mathrm{\Omega }_3)],`$ (23) $`\text{Im}T\left(\omega \right)={\displaystyle \frac{4\sqrt{2}\alpha }{3}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left[1+\left(1\right)^{k+l}\right]}{n!k!l!}}s_1^ns_2^ks_3^l`$ (24) $`\times `$ $`\left[g_{n+k+l}\left(\omega 1n\mathrm{\Omega }_1k\mathrm{\Omega }_2l\mathrm{\Omega }_3\right)g_{n+k+l}\left(\omega 1n\mathrm{\Omega }_1k\mathrm{\Omega }_2l\mathrm{\Omega }_3\right)\right].`$ (25) From here on, we use the system of units where $`\mathrm{}=1,`$ $`m=1,`$ $`\omega _{\mathrm{LO}}=1.`$ The factors $`s_i`$ ($`i=1,2,3`$) in Eqs. (23), (25) have the sense of oscillator strengths corresponding to the eigenmodes of the internal motion of a bipolaron. The aforesaid factors are expressed through variational frequency parameters of the bipolaron model functional (18) to (20), $$s_1=\frac{\mathrm{\Omega }_1^2\nu ^2}{\mathrm{\Omega }_1^3},s_2=\frac{\mathrm{\Omega }_2^2\nu ^2}{\mathrm{\Omega }_2\left(\mathrm{\Omega }_2^2\mathrm{\Omega }_3^2\right)},s_3=\frac{\nu ^2\mathrm{\Omega }_3^2}{\mathrm{\Omega }_3\left(\mathrm{\Omega }_2^2\mathrm{\Omega }_3^2\right)}.$$ (26) The functions $`f_n\left(\omega \right)`$ and $`g_n\left(\omega \right)`$ are given by the expressions: $$f_n\left(\omega \right)=\left(1\right)^n\frac{\left\|\omega \right\|^{n+\frac{1}{2}}}{B^{n+\frac{3}{2}}}e^{\frac{A\omega }{B}}\mathrm{\Theta }\left(\omega \right)\frac{\mathrm{\Gamma }\left(n+\frac{1}{2}\right)}{\pi A^{n+\frac{1}{2}}B}{}_{1}{}^{}F_{1}^{}(1,\frac{1}{2}n,\frac{A\omega }{B}),$$ (27) $$g_n\left(\omega \right)=\frac{\omega ^{n+\frac{1}{2}}}{B^{n+\frac{3}{2}}}e^{\frac{A\omega }{B}}\mathrm{\Theta }\left(\omega \right),$$ (28) with the parameters $$A=\frac{\mu }{m\mathrm{\Omega }_1}+\frac{a^2}{\mathrm{\Omega }_2}+\frac{b^2}{\mathrm{\Omega }_3},\mu =\frac{\mathrm{\Omega }_1^2\nu ^2}{\mathrm{\Omega }_1^2},a=\sqrt{\frac{\mathrm{\Omega }_2^2\nu ^2}{\mathrm{\Omega }_2^2\mathrm{\Omega }_3^2}},b=\sqrt{\frac{\nu ^2\mathrm{\Omega }_3^2}{\mathrm{\Omega }_2^2\mathrm{\Omega }_3^2}},B=\frac{\nu ^2}{\mathrm{\Omega }_1^2}.$$ (29) The optical conductivity for a 2D bipolaron is related to that for a 3D bipolaron, which is given by the formulas (23) and (25), by the scaling relation (cf. Ref. ): $$Re\sigma _{2\mathrm{D}}(\omega ,\alpha )=Re\sigma _{3\mathrm{D}}(\omega ,\frac{3\pi }{4}\alpha ).$$ (30) It follows from Eqs. (23) and (25) that the eigenfrequencies $`\mathrm{\Omega }_2`$ and $`\mathrm{\Omega }_3`$ appear in the optical absorption spectra only in such combinations $`\left(k\mathrm{\Omega }_2+l\mathrm{\Omega }_3\right)`$ that $`\left(k+l\right)`$ is an even integer. This selection rule is determined by the symmetry of these eigenmodes (a schematic picture of the internal motion of the bipolaron model system see, e. g., in Ref. ). The only normal coordinates of eigenmodes with the frequencies $`\mathrm{\Omega }_2`$ and $`\mathrm{\Omega }_3`$ (let us denote these coordinates as vectors $`Q_2`$ and $`Q_3`$) are antisymmetric with respect to the permutation of electrons $`r_1r_2.`$ Both the exact and model Lagrangians are symmetric with respect to this permutation. As a result of this symmetry, we obtain the selection rule $$\underset{i=1,2,\mathrm{}}{}\underset{k=x,y,z}{}\left[Q_{2k}^{m_{ki}}\left(t_i\right)Q_{3k}^{m_{ki}^{}}\left(t_i^{}\right)\right]_{S_0}=0,\mathrm{when}\underset{i}{}\underset{k=x,y,z}{}\left(m_{ki}+m_{ki}^{}\right)=\mathrm{odd}.$$ (31) Hence, only the combinations $`\left(_{i=1,2,\mathrm{}}_{k=x,y,z}\left[Q_{2k}^{m_{ki}}\left(t_i\right)Q_{3k}^{m_{ki}^{}}\left(t_i^{}\right)\right]\right)`$ with an even number $`n=_i_{k=x,y,z}\left(m_{ki}+m_{ki}^{}\right)`$ can give a nonzero contribution into the average $`\rho _𝐪\left(t\right)\rho _𝐪\left(0\right)_{S_0}`$ of Eq. (13). The operator of an oscillator coordinate $`Q_{jk}`$ describes quantum transitions between energy levels of the corresponding oscillator with the frequency $`\mathrm{\Omega }_j`$ $`\left(j=1,2,3\right)`$. Consequently, the aforesaid combinations of coordinates provide the transitions, which change the energy of the model system by values $`\left(k\mathrm{\Omega }_2+l\mathrm{\Omega }_3\right)`$ with even $`\left(k+l\right).`$ ## III Optical absorption spectra In order to give a physical interpretation of the calculated bipolaron optical absorption spectra, we refer first to a description of those for a polaron in Refs. . In the polaron optical absorption spectra, for intermediate and large $`\alpha ,`$ there is an intense (zero-phonon) absorption peak corresponding to a transition from the ground state to the first relaxed excited state (RES). The relaxed excited state is created if the electron in the polaron is excited while the lattice readapts to a new electronic configuration . A shoulder at the low-frequency side of the RES peak is attributed to one-phonon transitions through the scattering states (ScS): excitations of the polaron system characterized by the presence of a finite number of real phonons along with the polaron . A broad peak, positioned at a higher frequency than the RES peak, is attributed to transitions to Franck-Condon (FC) states: internal excited polaron states for which the lattice polarization is that of the polaron in its ground state . This picture is well-founded physically and is in agreement with the prediction of the aforesaid peaks, which was formulated within the strong-coupling approach . In Fig. 1, the real $`\text{Re}T\left(\omega \right)`$ and imaginary $`\text{Im}T\left(\omega \right)`$ parts of the memory function for a 3D (2D) bipolaron are plotted as a function of frequency $`\omega `$ for $`\alpha =7`$ ($`2.971`$). These values of $`\alpha `$ are close to the minimal possible $`\alpha `$ for the bipolaron formation: $`\alpha _{\mathrm{min}}=6.8`$ ($`2.9`$)(see Ref. ). The peaks of the imaginary part of the memory function are positioned near the points (explicitly indicated in Fig. 1) $$\omega _{nkl}n\mathrm{\Omega }_1+k\mathrm{\Omega }_2+l\mathrm{\Omega }_3+1,$$ (32) where $`n,`$ $`k,`$ $`l`$ are the non-negative integers ($`k+l=`$ even integer). Analogously to the case of polaron optical absorption , the peaks of the bipolaron optical absorption spectra (Fig. 2) correspond to: (i) peaks of $`\text{Im}T\left(\omega \right)`$ and (ii) roots of the equation $`\omega ^2ReT\left(\omega \right)=0`$ (the crossing points of solid and dotted curves in the panel “a” of Fig. 1). Following the physical interpretation developed in Refs. , we suggest that the peaks corresponding to zeros of $`\left[\omega ^2ReT\left(\omega \right)\right]`$ can be attributed to transitions to bipolaron RES. Due to a larger number of the internal degrees of freedom for a bipolaron, than that for a polaron, there are several types of RES for a bipolaron. The positions of peaks of the optical absorption coefficient, corresponding to those of $`\text{Im}T\left(\omega \right),`$ are close to the values $`\omega _{nkl}`$ determined by Eq. (32), so that $`\omega _{nkl}`$ have the sense of the frequencies of transitions from the ground state to certain internal states of a bipolaron. Following the physical interpretation developed in Refs. , we can attribute these peaks to transitions into Franck-Condon (FC) bipolaron states (in the case when at least one of the numbers $`(n,k,l)0,`$ i. e., for excited states of a bipolaron). Within the same picture, we suggest that the low-intensity peak at $`\omega =1`$ is provided by the transition into the lowest scattering state of a bipolaron. It is worth mentioning, that the main resonances in oscillator strength are RES. In order to classify peaks of the bipolaron optical absorption spectra, we have calculated transition frequencies to several bipolaron RES and FC states within the strong-coupling (adiabatic) approach. The results are shown in Table 1. The transitions are classified with respect to the symmetry of states of the internal bipolaron motion and of the motion of a bipolaron as a whole (labelled, respectively, by small and capital letters within standard denotations of states with different orbital moments). “$`S`$” and “$`s`$” denote the states with the orbital momentum $`l=0,`$ while “$`P`$” and “$`p`$” are those with $`l=1.`$ The trial wave functions are of the same type as those used for the treatment of polarons in Ref. . Table 1. Transition frequencies (in units of the LO phonon frequency) to bipolaron RES and FC states calculated within the adiabatic strong-coupling theory. \| $`\alpha `$ \| \| RES \| \| \| FC \| \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| \| $`pS`$ \| $`sP`$ \| $`pP`$ \| $`pS`$ \| $`sP`$ \| $`pP`$ \| \| $`7`$ \| $`2.2171`$ \| $`5.9877`$ \| $`5.7282`$ \| $`5.4633`$ \| $`9.9396`$ \| $`12.5210`$ \| \| $`8`$ \| $`2.8957`$ \| $`7.8207`$ \| $`7.4818`$ \| $`7.1328`$ \| $`13.0008`$ \| $`16.3627`$ \| \| $`9`$ \| $`3.6648`$ \| $`9.8981`$ \| $`9.4686`$ \| $`8.9364`$ \| $`16.2629`$ \| $`20.4673`$ \| \| $`10`$ \| $`4.5243`$ \| $`12.220`$ \| $`11.6895`$ \| $`11.1019`$ \| $`20.1550`$ \| $`25.4065`$ \| Comparing the peak positions of Fig. 2 with the results of Table 1 (which are obtained in the framework of the adiabatic approach, and therefore can be applied only approximately), we have attributed several optical absorption peaks to transitions to RES and FC states of a definite symmetry, as shown at the figure. It should be mentioned, that every RES peak has a lower frequency, than the corresponding FC peak, since the lattice relaxation leads to a lowering of the energy of the electron-phonon system. The linewidth of peaks of bipolaron optical absorption spectra is determined by the bipolaron recoil in the scattering process. Every FC peak is described by a function $`g_{n+k+l}\left(\omega \omega _{nkl}\right)`$ of Eq. (25), where $`g_n\left(\omega \right)`$ is given by Eq. (28). Consequently, the linewidth of a FC peak can be estimated, using the following characteristic of the function (28): $$\sigma \left[g_n\left(\omega \right)\right]=\sqrt{\omega ^2\omega ^2},\omega ^k\frac{_0^{\mathrm{}}\omega ^kg_n\left(\omega \right)𝑑\omega }{_0^{\mathrm{}}g_n\left(\omega \right)𝑑\omega }.$$ (33) Performing the integrations over $`\omega `$ analytically, we find that $`\sigma \left[g_n\left(\omega \right)\right]={\displaystyle \frac{B}{A}}\left(n+{\displaystyle \frac{3}{2}}\right)^{1/2},`$where the parameters $`A`$ and $`B`$ are given by Eq. (29). The ratio $`\frac{A}{B}`$ is of the same order as $`\left[2\left(M+1\right)\right],`$ which is used for estimation of the bipolaron effective mass . Therefore, the characteristic linewidth of a FC peak corresponding to a definite frequency $`\omega =\omega _{nkl}`$ can be qualitatively estimated as $`\mathrm{\Gamma }_{nkl}^{\left(FC\right)}{\displaystyle \frac{1}{2\left(M+1\right)}}\left(n+k+l+3/2\right)^{1/2}.`$Since the bipolaron effective mass is larger than that of a polaron, the bipolaron optical absorption spectrum consists of a series of peaks which are narrower in comparison with those of the polaron absorption spectrum. Let us define $`N(\omega ,\mathrm{\Delta }\omega )`$ as the number of possible combinations of $`(n,k,l)`$ for which the frequencies $`\left\{\omega _{nkl}\right\}`$ occupy the interval $`\omega \omega _{nkl}\omega +\mathrm{\Delta }\omega .`$ The higher is $`\omega `$ (at a fixed interval length $`\mathrm{\Delta }\omega `$), the larger is $`N(\omega ,\mathrm{\Delta }\omega ).`$ Hence, the density of the FC peaks increases with increasing frequency. Starting from a definite frequency range (where $`N(\omega ,\mathrm{\Delta }\omega )\mathrm{\Gamma }_{FC}\left(\omega \right)\mathrm{\Delta }\omega `$), FC peaks merge into a continuous band. In the limiting case $`\omega 1,`$ the imaginary part of the memory function (25) asymptotically behaves as $$ImT\left(\omega \right)\|_{\omega 1}\frac{2\alpha }{3}\omega ^{1/2}.$$ (34) Hence, the bipolaron absorption coefficient in the limiting case of high frequencies has the same asymptotic behavior as the polaron absorption coefficient (this asymptotic behavior was also derived in Ref. for the weak-coupling polaron optical absorption): $$\mathrm{\Gamma }\left(\omega \right)\|_{\omega 1}\omega ^{5/2}.$$ (35) This behavior has a clear physical explanation: in the high-frequency limit, only one-phonon scattering processes give a contribution into the optical absorption by the electron-phonon system. In this limit, the optical absorption coefficients for bipolarons and polarons at arbitrary coupling strength are described by one and the same asymptotic dependence. It was emphasized in Refs. , that the frequency dependence (35) of the polaron optical absorption coefficient at $`\omega 1`$ differs from the Drude-like dependence $`\mathrm{\Gamma }^{\left(\mathrm{Drude}\right)}\left(\omega \right)\|_{\omega 1}\omega ^2`$. For the bipolaron, the optical absorption spectrum turns into a continuous band at substantially higher frequencies, than for the polaron, due to a comparatively large effective mass of a bipolaron. This frequency region is not presented in Fig. 2. The bipolaron absorption coefficient in this region is very small when compared with that shown in Fig. 2. ## IV Sum rules and conclusions From the general analytical properties of the memory function, for the real part of the optical conductivity per one electron $$\text{Re}\sigma \left(\omega \right)=\frac{e^2}{m}\frac{\omega \text{Im}T\left(\omega \right)}{\left[\omega ^2\text{Re}T\left(\omega \right)\right]^2+\left[\text{Im}T\left(\omega \right)\right]^2}$$ (36) the following sum rule was derived in Ref. : $$\frac{\pi e^2}{m^{}}+_{\epsilon +0}^{\mathrm{}}\text{Re}\sigma \left(\omega \right)𝑑\omega =\frac{\pi e^2}{2m},$$ (37) where $`m^{}`$ is the polaron effective mass. For the bipolaron optical absorption, the same equation (37) is valid with $`m^{}`$ the bipolaron effective mass. Our numerical test of the sum rule (37) has confirmed, that Eq. (37) is fulfilled within the chosen relative accuracy (up to 10<sup>-3</sup>). In Ref. , the ground state theorem for a polaron has been derived which relates the polaron ground state energy to the first moment of the optical absorption spectra. The extension of this theorem to the bipolaron can be performed in the same way as in Ref. . The analog of Eq. (25) from Ref. for a bipolaron is $$E_0(\alpha _2,\eta )E_0(\alpha _1,\eta )=\frac{3m\mathrm{}}{\pi e^2}_{\alpha _1}^{\alpha _2}\frac{d\alpha }{\alpha }_0^{\mathrm{}}\text{Im}\chi _{jj}(\omega ,\alpha ,\eta )𝑑\omega ,$$ (38) where $`\alpha _1`$ and $`\alpha _2`$ are two arbitrary values of $`\alpha ,`$ $`\eta =\epsilon _{\mathrm{}}/\epsilon _0`$ is the ratio of the high-frequency and static dielectric constants, $`\chi _{jj}(\omega ,\alpha ,\eta )`$ is the current-current correlation function . For a bipolaron, we choose $`\eta `$ and both $`\alpha _1`$ and $`\alpha _2`$ within the region of the bipolaron stability. The function $`\text{Im}\chi _{jj}\left(\omega \right)`$ for a bipolaron is expressed through the memory function $`T\left(\omega \right)`$ as $$\text{Im}\chi _{jj}\left(\omega \right)=\frac{2\omega ^2\text{Im}T\left(\omega \right)}{\omega ^42\omega ^2\text{Re}T\left(\omega \right)+\left\|T\left(\omega \right)\right\|^2}.$$ (39) In Fig. 3, the ground-state theorem (38) is illustrated numerically. The solid curve shows the difference between the ground state energies $`E_0(\alpha _2,\eta )E_0(\alpha _1,\eta )`$ (calculated by the variational method using the variational functional of Ref. ) for two values of $`\alpha `$ at $`\eta =0.`$ The value $`\alpha _1`$ is taken $`\alpha _1=7.`$ The dashed curve shows the values of the right-hand side of Eq. (38). As seen from Fig. 3, there exists an excellent agreement between the quantities entering the independently calculated left-hand and the right-hand sides of the equation (38). In the present paper, we have treated the bipolaron optical absorption within the memory-function path-integral approach. The derived optical absorption spectra demonstrate a rich structure of relatively narrow peaks, which is clearly related to the internal states of a bipolaron. We have attributed those peaks of the bipolaron optical absorption spectra to transitions from the ground state to (i) scattering states, (ii) relaxed excited states, (iii) Franck-Condon states. Every RES peak is shifted to a lower frequency with respect to the corresponding FC peak. This interpretation is performed within a unified approach to the polaron and bipolaron problems, which takes its origin in Refs. . The sum rules, developed for the polaron optical conductivity , are extended to the case of a bipolaron. ###### Acknowledgements. This work has been supported by the BOF NOI (UA-UIA), GOA BOF UA 2000, IUAP, FWO-V. projects G.0287.95, 9.0193.97, and the W.O.G. WO.025.99N (Belgium). Figure captions Fig. 1. Real \[panel “a”\] and imaginary \[panel “b”\] parts of the memory function $`T\left(\omega \right)`$ for the optical absorption coefficient of a 3D (2D) bipolaron for $`\alpha =7`$ (2.971), $`\eta =0.0037`$. The optimal values of the variational frequencies are indicated in the figure. Fig. 2. Optical absorption spectra of a 3D (2D) bipolaron for $`\alpha =7`$ (2.971), $`\eta =0.0037`$ \[panel “a”\], for $`\alpha =8`$ (3.395), $`\eta =0.023`$ \[panel “b”\], and for $`\alpha =9`$ (3.820), $`\eta =0.023`$ \[panel “c”\]. Fig. 3. Illustration of the ground-state theorem (38) for the bipolaron optical absorption. The index “1” denotes the ground state energies calculated by the variational method \[the left-hand side of Eq. (38)\], while the index “2” denotes the ground state energies calculated by integration of the bipolaron optical absorption spectra \[the right-hand side of Eq. (38)\].
warning/0004/astro-ph0004052.html	ar5iv	text	# 1 Introduction ## 1 Introduction During its giant outburst in April 1997, the nearby BL Lac object Mkn 501 (at redshift z=0.034) emitted photons up to $`24`$ TeV and $`0.5`$ MeV in the $`\gamma `$-ray and X-ray bands, respectively, and has proved to be the most extreme TeV-blazar observed so far (e.g. Catanese et al 1997, Pian et al 1998, Protheroe et al 1998, Quinn et al 1999, Aharonian et al 1999). This energy is the highest so-far observed for any BL Lac object, and the flux is approximately 2 orders of magnitude higher than the synchrotron peak at its quiescent level. BeppoSAX and OSSE observations (Maraschi 1999) suggest that the X-ray spectrum is curved at all epochs, and the spectrum during flaring has been fitted by a multiply-broken power-law (Bednarek & Protheroe 1999). COMPTEL has not seen any significant signal from Mkn 501 at any time (Collmar 1999), while a $`3\sigma `$ upper limit of $`F(>100\mathrm{MeV})<3.6\times 10^7`$cm<sup>-2</sup> s<sup>-1</sup> has been derived for the April 1997 EGRET viewing period (Catanese et al 1997). A flux increase at TeV-energies was also observed with the Whipple, HEGRA and CAT telescopes (Catanese et al 1997), with the most intense flare peaking on April 16 at a level $`100`$ times higher than during its quiescent flux. The non-detection of Mkn 501 by EGRET indicates that most of the power output of the high energy component is in the GeV-TeV range. The TeV-observations revealed a power-law spectrum with photon index $`2`$ up to $`10`$ TeV and a gradual steepening up to 24 TeV. The extragalactic diffuse infrared background leads to significant extinction of $`\gamma `$-rays through $`\gamma \gamma `$-pair production above 10 TeV. The extinction-corrected TeV-spectrum (e.g. Bednarek & Protheroe 1999), shows the spectral energy distribution (SED) peaking at $`2`$ TeV. Optical observations did not show any significant variations (Buckley & McEnery 1997), indicating that the change in the low energy part of the SED was mainly confined to the X-ray band above $`0.1`$ keV. Various models have been proposed to explain the observed $`\gamma `$-ray emission from TeV-blazars, all of which are identified as high-frequency peaked BL Lac objects. Leptonic models, in which electrons inverse-Compton scatter a population of low energy photons to high energies, currently dominate the thinking of the scientific community. Because of the low luminosity of accretion disks in BL Lacs, the main target photons for the relativistic electrons would be the synchrotron photons produced by the same relativistic electron population, as in the synchrotron self-Compton (SSC) model. An alternative scenario for the production of the observed $`\gamma `$-ray flux has been proposed involving pion photoproduction by energetic protons with subsequent synchrotron-pair cascades initiated by decay products (photons and $`e^\pm `$) of the mesons (e.g. Mannheim et al 1991, Mannheim 1993). These proton-initiated cascade (PIC) models could, in principle, be distinguished by the observation of high energy neutrinos produced as a result of photoproduction. In this paper, we consider the April 1997 flare of Mkn 501 in the light of a modified Synchrotron Proton Blazar (SPB) model. We assume that electrons ($`e^{}`$) and protons ($`p`$) are accelerated by 1st order Fermi acceleration at the same shock. The relativistic $`e^{}`$ radiate synchrotron photons which serve as the target radiation field for proton-photon interactions, and for the subsequent pair-synchrotron cascade which develops as a result of photon-photon pair production. This cascade redistributes the photon power to lower energies where the photons escape from the emission region, or “blob,” which moves relativistically in a direction closely aligned with our line-of-sight. Until recently, this model was not able to reproduce the general features of the double-humped blazar spectral energy distribution (SED), but produced a rather featureless spectrum (see e.g. Mannheim 1993), nor could it explain correlated X-ray/TeV-variability. Here, we present a comprehensive description of our Monte-Carlo simulations of a stationary SPB model, including all relevant emission processes, and show that this model is indeed capable of reproducing a double-humped SED as observed. Here, the origin of the TeV-photons are proton synchrotron radiation, as first proposed by Mücke & Protheroe (1999); a similar model has also been proposed by Aharonian (2000), and Rachen (1999) presented speculations about $`\mu ^\pm `$\- and proton-synchrotron radiation leading to narrow cascade spectra during flares, which might explain correlated X-ray/TeV-variability. Jet energetics and limits from particle shock acceleration, however, put severe constraints on this scenario. The goal of this paper is to discuss the physical processes included in our SPB model Monte-Carlo code, and give the results of applying this code, as an example, to reproduce the SED of the giant flare from Mkn 501 which occurred in April 1997. A comprehensive study of the whole parameter-space (magnetic field, Doppler factor, etc.) for this model will be the subject of a subsequent paper. In Section 2, we discuss constraints on the maximum particle energies imposed by the co-acceleration scenario, and by the pion production threshold. Section 3 is devoted to the emission processes in the present model. Energy losses and particle production are treated in Sect. 3.1, while the cascade calculations, including a brief description of our code, are outlined in Sect. 3.2. In Sect. 4 we apply our model to the April 1997 flare of Mkn 501. The multifrequency photon spectrum is shown in Sect. 4.1, while in Sect. 4.2 the predicted neutrino spectrum is discussed. We conclude with a discussion and summary in Section 5. ## 2 The Co-acceleration Scenario In the present model, shock accelerated protons ($`p`$) interact in the synchrotron photon field generated by the electrons ($`e^{}`$) co-accelerated at the same shock. This scenario may put constraints on the maximum achievable particle energies. The usual process considered for accelerating charged particles in high energy astrophysics is diffusive shock acceleration (see e.g. Bell 1978, Drury 1983, Blandford & Eichler 1987, Biermann & Strittmatter 1987, Jokipii 1987, Jones & Ellison 1991), in which particles undergo collisionless scattering, e.g. by Alfvén waves, in the upstream and downstream plasma. Charged particles with gyroradii larger than the thickness of the shock front propagate with diffusion coefficients $`\kappa _1`$ and $`\kappa _2`$ in the upstream and downstream plasma, respectively, for propagation parallel to the shock normal. In the shock frame the plasma flow velocity changes from $`u_1=\beta _1c`$ in the upstream region to $`u_2`$ in the downstream region. In this paper, for simplicity we restrict ourselves to non-relativistic shocks, and postpone discussion of relativistic shock acceleration to a later paper. The acceleration time scale for non-relativistic shocks is given by $$t_{\mathrm{acc}}=\frac{3r_c\beta }{(r_c1)u_1^2}(\kappa _1+r_c\kappa _2)$$ (1) where $`r_c4`$ is the compression ratio for the case of strong shocks in a non-relativistic monoatomic ideal gas. If the magnetic field is governed by an ordered component, the orientation of the shock normal to the main magnetic field direction becomes important. In general, the diffusion coefficient can be written as $$\kappa _i=\kappa _{i,\|\|}\mathrm{cos}^2\theta _i+\kappa _{i,}\mathrm{sin}^2\theta _i,i=1,2$$ (2) where $`\theta _i`$ is the angle between the magnetic field and the axis connecting the upstream ($`i=1`$) and downstream ($`i=2`$) regions. In the diffusion limit, kinetic theory relates the parallel and perpendicular diffusion coefficients through $$\kappa _{\|\|}=[1+\eta ^2]\kappa _{}=\frac{1}{3}\lambda _{\|\|}\beta c$$ (3) where $`\eta =\lambda _{\|\|}/r_g`$ with $`\lambda _{\|\|}`$ being the mean free path parallel to the magnetic field, $`r_g=\beta \gamma mc^2/eB`$ is the particle’s gyroradius, $`m`$ and $`\gamma =(1\beta ^2)^{1/2}`$ are the particle’s mass and Lorentz factor, respectively, and $`B`$ is the magnetic field strength in the upstream region. The mean free path is, in general, a function of the particle energy through its gyroradius, and is dependent on the spectrum of the magnetic turbulence. In the small angle scattering approximation (i.e., if Alfvén waves dominate the particle deflection with wavelength equal to the particle gyroradius; see Drury 1983) we have $$\lambda =\frac{B^2r_g}{8\pi I(k)k}.$$ (4) This spectrum $`I(k)`$ is usually expressed as a power law of the wave number $`k`$ in the turbulent magnetic field: $$I(k)k^\delta .$$ (5) $`\delta =5/3`$ corresponds to Kolmogorov turbulence which may be common in astrophysical environments (Biermann & Strittmatter 1987), while $`\delta =1`$ corresponds to Bohm diffusion, and is often considered for simplicity. For strong magnetic fields, Kraichnan turbulence $`\delta =3/2`$ (Kraichnan 1965) may be present. In the following, we consider $`\delta `$ as a free parameter. The mean free path may then be expressed as (see Biermann & Strittmatter 1987) $$\lambda _{\|\|}=\frac{r_g}{b(\delta 1)}\left(\frac{r_{g,\mathrm{max}}}{r_g}\right)^{\delta 1}\text{for}\delta 1$$ (6) $$\lambda _{\|\|}=\frac{r_g}{b}\left[\mathrm{ln}\left(\frac{r_{g,\mathrm{max}}}{r_{g,\mathrm{min}}}\right)\right]^1\text{for}\delta =1$$ where $`b`$ is the ratio of the turbulent to ambient magnetic energy density, $`r_{g,\mathrm{max}}`$ is the gyroradius of the most energetic protons, and has the same order of magnitude as the system size, and $`r_{g,\mathrm{min}}`$ corresponds to the smallest turbulence scale. The mean free path, and consequently the acceleration time at maximum energy is only slightly dependent on the turbulence spectrum in the case of protons, whereas for electrons the acceleration time at maximum energy shows a strong dependence on the magnetic turbulence spectrum adopted. We expect $`b1`$, since otherwise the energy density in particles would not be able to be confined by the ambient field (Biermann & Strittmatter 1987). With these relations, the acceleration time scale may be re-written as $$t_{\mathrm{acc}}=\frac{r_g\beta c}{u_1^2}F(\theta _1,\eta )$$ (7) where $$F(\theta _1,\eta )=\frac{\eta r_c}{r_c1}\left[\mathrm{cos}^2\theta _1+\frac{\mathrm{sin}^2\theta _1}{(1+\eta ^2)}+\frac{r_c\mathrm{cos}^2\theta _1+r_c^3\mathrm{sin}^2\theta _1/(1+\eta ^2)}{(\mathrm{cos}^2\theta _1+r_c^2\mathrm{sin}^2\theta _1)^{3/2}}\right]$$ (8) The diffusion approximation used here limits the maximum mean free path to $`\eta <\beta /\beta _1`$ (Jokipii 1987). If the particle spectra are cut off due to synchrotron losses, balancing the acceleration time scale with the loss time scale determines the maximum Lorentz factors of protons, $$\gamma _{p,\mathrm{max}}=2.1\times 10^{11}\beta _1\left[\beta BF(\theta _1,\eta _{p,\mathrm{max}})\right]^{1/2}$$ (9) and electrons $$\gamma _{e,\mathrm{max}}=1.2\times 10^8\beta _1\left[\beta BF(\theta _1,\eta _{e,\mathrm{max}})\right]^{1/2}$$ (10) where $`B`$ is in Gauss. The ratio of the maximum proton Lorentz factor to the maximum electron Lorentz factor is then $$\frac{\gamma _{p,\mathrm{max}}}{\gamma _{e,\mathrm{max}}}\frac{m_p}{m_e}\left[\frac{F(\theta _1,\eta _{e,\mathrm{max}})}{F(\theta _1,\eta _{p,\mathrm{max}})}\right]^{1/2}$$ (11) where the equality corresponds to the maximum proton energy being determined by synchrotron losses, and the inequality to the maximum proton energy being determined instead by adiabatic losses. For parallel shocks this relation is consistent with the results found by Biermann & Strittmatter (1987). The corresponding acceleration time scales at the maximum particle energies, if determined by synchrotron losses, are $$t_{\mathrm{acc},\mathrm{p},\mathrm{max}}=2.2\times 10^7B^{3/2}\beta _1^1[\beta F(\theta _1,\eta _{p,\mathrm{max}}]^{1/2}\mathrm{s}$$ (12) for protons, and $$t_{\mathrm{acc},\mathrm{e},\mathrm{max}}=6.6B^{3/2}\beta _1^1[\beta F(\theta _1,\eta _{e,\mathrm{max}}]^{1/2}\mathrm{s}$$ (13) for electrons. In the present paper, we shall adopt $`\beta _1=0.5`$, $`\beta =1`$ and $`r_c=4`$. The geometry dependent term $`[F(\theta _1,\eta _{p,\mathrm{max}})]^{1/2}`$ is plotted in Fig. 1a as the dashed curves for different $`\eta _{p,\mathrm{max}}`$ values and $`r_c=4`$. It can be seen that highly oblique shocks allow proton acceleration on very short time scales. In the limiting case of a perpendicular shock, drift–shock acceleration drives the energy gain, and the finite size of the shock front restricts the maximum particle energy. The maximum proton Lorentz factor $`\gamma _p`$ is reached for the maximum drift distance, the shock size, and is given by $$\gamma _{p,\mathrm{max}}=3.2\times 10^7\beta _1RB$$ (14) with $`R`$ in cm and $`B`$ in Gauss. At their maximum Lorentz factors, the acceleration process for protons is considerably slower than for electrons, and so the proton acceleration time must be consistent with the observed variability time, $`t_{\mathrm{var}}`$, $$t_{\mathrm{var}}Dt_{\mathrm{acc},\mathrm{p},\mathrm{max}}$$ (15) where $`D`$ the Doppler factor. This can be converted to a constraint on the geometry dependent term using Eq. 12, and is plotted as a function of $`\theta _1`$ in Fig. 1a for a typical set of TeV-blazar parameters ($`B20`$ G, $`D10`$, $`t_{\mathrm{var}}=12`$ hours). The region below the solid line is allowed by the variability constraint (Eq. 15), and gives for each $`\eta _{p,\mathrm{max}}`$ value a minimum shock angle between $`69^{}`$ and $`88^{}`$, depending on $`\eta _{p,\mathrm{max}}`$. Thus, proton shock acceleration on hour time scale in hadronic models can only take place in oblique shocks, and the maximum drift distance may restrict the maximum energy gain rather than the gyroradius. Using the same shock geometry and magnetic turbulence spectra for both protons and electrons, we find that the ratio $`F(\theta _1,\eta _{e,\mathrm{max}})/F(\theta _1,\eta _{p,\mathrm{max}})`$ does not vary by more than a factor of 2 within the allowed shock angle range (see Fig. 1b), allowing us to adopt an average value for this ratio for a given parameter combination. Inserting this ratio into Eq. 11 then restricts the ratio of the allowed maximum particle energies to the range below the solid lines shown in Fig. 2. Points exactly on this line represent models where the particle spectra are limited by synchrotron losses and the acceleration time scale is exactly the variability time scale; points below this line apply if adiabatic losses are dominant for protons or the proton acceleration time scale is shorter than the variability time scale in the jet frame. We note that for a given maximum proton energy, the highest maximum electron energy occurs with Bohm diffusion. In hadronic blazar models pion photoproduction is essential for neutrino production. The threshold for this process is given by $`ϵ_{\mathrm{max}}\gamma _{p,\mathrm{max}}=0.0745`$ GeV where $`ϵ_{\mathrm{max}}`$ is the maximum photon energy of the target field, which in the Synchrotron Proton Blazar models is in turn produced by the co-accelerated $`e^{}`$. In the $`\delta `$-function approximation for the synchrotron emission, $`ϵ_{\mathrm{max}}=(3/8)\gamma _{e,\mathrm{max}}^2(B/B_{\mathrm{cr}})m_ec^2`$ with $`B_{\mathrm{cr}}=4.414\times 10^{13}`$ G. Inserting $`ϵ_{\mathrm{max}}`$ into the threshold condition, we find $$\gamma _{p,\mathrm{max}}1.72\times 10^{16}\left(\frac{B}{1\mathrm{Gauss}}\right)^1\gamma _{e,\mathrm{max}}^2$$ (16) and this is shown in Fig. 2 as the dashed lines for various magnetic field strengths. Together with Eq. 11 the allowed range of maximum particle energies is then restricted to the area below the solid lines and above the dashed lines in Fig. 2, as shown for example as the shaded area for $`B20`$ G and Kolmogorov turbulence. Inspecting hadronic models as presented in the literature (e.g. Mannheim 1993, Mannheim et al 1996, Rachen 1999), we find that most models which are able to fit the observations lie above the Bohm diffusion line (see Fig. 2), indicating that the turbulence spectrum required in common hadronic blazar jet models is likely to be of Kolmogorov/Kraichnan type. Rachen (1999) speculates that the transition between Kolmogorov and Kraichnan type turbulence could be responsible for the difference between low- and high frequency peaked BL Lacs. ## 3 Emission Processes In the present model, the co-accelerated $`e^{}`$ are assumed to produce most of the observed low energy part of the SED by synchrotron radiation, and this is assumed to be the target radiation field for $`p\gamma `$ interactions and subsequent cascading. The observed hardening of the spectrum with rising flux, has recently been convincingly reproduced by a shock model with escape and synchrotron losses (Kirk et al 1998). The spectral slopes in this model are controlled by synchrotron cooling, and thus naturally explain the temporal behaviour of the spectral index as observed in the X-ray band. The flaring behaviour is explained by the shock front running into plasma whose density is locally enhanced, and which thus increases the number of particles injected into the acceleration process. The increase in the plasma density is accompanied by an increase in the magnetic field, leading to a higher acceleration rate and a shift of the maximum particle energies to higher energies. This picture departs from the standard explanation in which the flat synchrotron spectra appear as a result of a superposition of several local self-absorbed synchrotron spectra with changing self-absorption frequency, adopted in previous PIC models (e.g. Mannheim 1993). For simplicity, and because we do not wish to include additional parameters, we use the same magnetic field for synchrotron radiation as for acceleration. For normal shocks this approximation is justified as the magnetic fields either side of the shock are similar. This approximation might even be justified for oblique shocks as a lower magnetic field in the upstream region, compared to the downstream region, implies a higher diffusion coefficient and time spent upstream, increasing the synchrotron losses there and partially compensating for having a lower field upstream. Also, at oblique shocks, reflection at the shock front itself is thought to be more important than diffusion in the downstream region (Kirk & Heavens 1989), so that accelerating particles spend most of their time upstream. In addition, we assume that pitch-angle scattering maintains quasi-isotropic particle distributions, and all radiating particles are confined to the homogeneous emission region. ### 3.1 Energy Losses There are several energy loss/interaction processes which are important for protons, electrons and photons in a dense radiation field produced by relativistic electrons co-accelerated along with the protons: protons interact with photons, resulting in pion production and (Bethe-Heitler) pair production; electrons, muons, protons and charged pions emit synchrotron radiation; photons interact with photons by pair production. We shall show below that, for the present model, Inverse Compton emission by the electrons can be neglected. For simplicity we represent the observed synchrotron spectrum of Mrk 501 during flaring, the target photon field for the $`p\gamma `$-collisions and photon-photon pair production, as a broken power-law: $$n(ϵ)\{\begin{array}{cc}ϵ^{1.6}\hfill & \text{for}10^7\mathrm{eV}ϵ1.6\mathrm{keV}\hfill \\ ϵ^{1.8}\hfill & \text{for}1.6\mathrm{keV}ϵ42\mathrm{keV}\hfill \end{array}$$ For determining the photon density of the target field, the dimension of the emission region, assumed to be spherical, must be known. This can be estimated by setting the photon crossing time equal to the variability time scale (in the jet frame – in the remainder of this section all quantities are in the jet frame unless noted otherwise), making the implicit assumption that the light crossing time scale determines the flux variations. The observed variability time scale can, in general, depend on: (i) the injection time scale for the energetic particles $`t_{\mathrm{acc}}(E)`$; (ii) the time needed for converting their energy into radiation, i.e. their energy loss time scale $`t_{\mathrm{loss}}(E)`$; (iii) the effective light crossing time $`t_{\mathrm{cross}}(E)(2R_{\mathrm{blob}}/c)\times P_{\mathrm{esc}}(E)`$ where $`R_{\mathrm{blob}}`$ is the geometrical blob radius and $`P_{\mathrm{esc}}(E)`$ is the energy dependent probability for photons escaping from the blob taking account of $`\gamma \gamma `$-pair production and diffusion during cascading. Hence, $$Dt_{\mathrm{var}}\mathrm{max}(t_{\mathrm{acc}},t_{\mathrm{loss}},t_{\mathrm{cross}}).$$ (17) For leptonic models, the time scales for energy losses and acceleration are typically significantly shorter than the crossing time, i.e. $`t_{\mathrm{loss}}`$,$`t_{\mathrm{acc}}t_{\mathrm{cross}}`$, and thus the radius of the emission region can be derived from the observed variability time scale. In addition, the emission region is assumed to be optically thin at 1 TeV and X-ray energies, implying that $`R_{\mathrm{TeV}}R_\mathrm{X}`$, where $`R_{\mathrm{TeV}}`$, $`R_\mathrm{X}`$ are the dimensions of the emitting region at 1 TeV and X-ray energies, respectively. This differs in two points from the hadronic blazar jet models: Firstly, the optical depth of the emission region is strongly energy dependent, leading to an effective, i.e. observed, thickness of the emission region $`R_{\mathrm{eff}}(E)R_{\mathrm{blob}}P_{\mathrm{esc}}(E)`$, which also depends on the energy. Diffusion can be neglected during cascading since the cascade processes are of leptonic origin, and are, in general, more rapid than diffusion. Thus, in SPB models the crossing time scale in the optically thick TeV-band is related to the crossing time scale at X-ray energies (optically thin), through $$t_{\mathrm{cross},\mathrm{TeV}}t_{\mathrm{cross},\mathrm{X}}[1\mathrm{exp}(\tau _{\gamma \gamma ,\mathrm{TeV}})]/\tau _{\gamma \gamma ,\mathrm{TeV}}$$ (18) (averaging over a homogeneous emission volume; see Rachen 1999). Hence, $`t_{\mathrm{cross},\mathrm{X}}`$ is the relevant time scale for estimating the radius $`R_{\mathrm{blob}}R_X`$ of the emission region. Secondly, the acceleration and/or energy loss time scales can be of the same order of magnitude as the crossing time scale. Acceleration, however, is always faster than the energy losses of the accelerating particles up to their maximum energy. Note that because of the leptonic nature of the cascade processes, the energy loss time scale will in general be determined by the (slower) hadronic processes. As a further consequence, flux variations are not significantly washed out by the cascading mechanism, but closely follow the crossing or hadronic loss time scales. If $`t_{\mathrm{loss}}t_{\mathrm{cross}}`$ the crossing time scale determines the flux variations, and in this case we can estimate the radius of the emission region through $`R_{\mathrm{blob}}0.5cDt_{\mathrm{var},\mathrm{X}}`$ with $`t_{\mathrm{var},\mathrm{TeV}}t_{\mathrm{var},\mathrm{X}}/\tau _{\gamma \gamma ,\mathrm{TeV}}`$. In our present model, efficient proton synchrotron emission determines the TeV-bump in the blazar SED, and so at the maximum proton energy $`t_{\mathrm{ad}}>t_{p,\mathrm{syn}}`$. The adiabatic loss time scale due to expansion is (Longair 1994) $$t_{\mathrm{ad}}=R_{\mathrm{ad}}^1\|R_{\mathrm{blob}}/\dot{R}_{\mathrm{blob}}\|=2\|B/\dot{B}\|R/u_1$$ (19) assuming magnetic flux conservation $`BR_{\mathrm{blob}}^2`$, and so $`t_{\mathrm{ad}}`$ is related to the size of the emission region $`R_{\mathrm{blob}}`$. Since $`t_{\mathrm{loss}}t_{p,\mathrm{syn}}`$ for the highest proton energies in our model responsible for the TeV-emission, $`t_{\mathrm{loss}}t_{\mathrm{cross}}`$, and so the size of the emission region can be determined by the variability time scale. The shortest doubling time measured by the Whipple Telescope in the 1997 data of Mkn 501 was approximately 2 hours (Quinn et al 1999), while HEGRA reported a lower doubling time of 15 hours (Aharonian et al 1999) or 12 hours (Krawczynski 1999). As a working hypothesis we adopt here a variability time scale of 12 hours, but will discuss also effects of smaller variability time scales. For $`t_{\mathrm{var}}12`$ hours we find $`R_{\mathrm{blob}}8\times 10^{15}`$cm for $`D10`$. Eq. 12 and 13 imply that the acceleration time scale for electrons is much smaller than for protons at their maximum energies, $`t_{acc,e,\mathrm{max}}t_{acc,p,\mathrm{max}}`$. If a $`\gamma `$-ray outburst thus corresponds to particle acceleration of a single $`p/e^{}`$ population, then an increase of the low energy synchrotron flux will occur before the TeV-flare. The amount of the lag depends on the proton acceleration time, and so on the shock parameters, in particular the diffusion coefficients. E.g. for the present model and the Kolmogorov spectrum $`t_{acc,p,\mathrm{max}}10^5`$ s (see Fig. 3), which for D=10 suggests a time lag of $`5`$ hours. For Bohm diffusion, the time lag would be $`14`$ hours. These time lags are consistent with current observational constraints. The relevant radiation and loss time scales for photomeson production, Bethe-Heitler pair production, $`p`$ synchrotron radiation, and adiabatic losses due to jet expansion, are shown in Fig. 3 together with the acceleration time scale. Synchrotron losses, which turn out to be at least as important as losses due to photopion production in our model, limit the injected $`p`$ spectrum to a Lorentz factor of $`\gamma _p3\times 10^{10}`$ for the assumed model parameters. We adopt a Kolmogorov spectrum of turbulence for the magnetic field structure ($`\delta =5/3`$), and so for any $`\eta _p1`$ value, variability arguments constrain the shock angle to $`\theta _175^{}`$ (see Fig. 1). The maximum proton energy could then be achieved, e.g., with $`\eta _p=10`$, $`\theta _1=85^{}`$ and $`u_1=0.5c`$. This is in agreement with the limit imposed on quasi-perpendicular shocks due to their finite shock size. Note that due to the non-zero shock angle, the acceleration time scale shown in Fig. 3 does not follow a strict power-law, but is curved. This is due to the non-linear dependence of $`F(\theta _1,\eta _p)`$ on the particle’s gyroradius (see Eqs. 6–8). A $`\gamma _p^2`$ proton spectrum, typical of shock accelerated particles, is used for $`2\gamma _p\gamma _{p,\mathrm{max}}`$, where $`\gamma _{p,\mathrm{max}}=3\times 10^{10}`$ is obtained from requiring $`t_{\mathrm{acc},p}=t_{\mathrm{syn},p}`$ at $`\gamma _{p,\mathrm{max}}`$. Rachen & Mészáros (1998) noted the importance of synchrotron losses of $`\mu ^\pm `$ (and $`\pi ^\pm `$) prior to their decay in AGN jets and GRBs. The critical Lorentz factors $`\gamma _\mu 2\times 10^9`$ and $`\gamma _\pi 4\times 10^{10}`$, above which synchrotron losses in the assumed magnetic field dominate over decay, lie below the maximum Lorentz factor for $`\mu ^\pm `$ and above the maximum Lorentz factor for $`\pi ^\pm `$. Thus, while $`\pi ^\pm `$-synchrotron losses can safely be ignored, $`\mu ^\pm `$-synchrotron losses should be included. For the parameters employed in this work we receive a target photon energy density of $`u_{\mathrm{target}}0.06\mathrm{TeV}/\mathrm{cm}^3`$, and a magnetic field energy density of $`u_B11.7\mathrm{TeV}/\mathrm{cm}^3`$ using $`B20`$ G. With $`u_Bu_{\mathrm{target}}`$ significant Inverse Compton radiation from the co-accelerated $`e^{}`$ is not expected. We also neglect secondary production interactions of relativistic protons with the ambient thermal plasma. In the present model the dominant loss process turns out to be proton synchrotron radiation. Thus, our assumption is justified if the thermal proton density does not exceed $`10^{9\mathrm{}10}`$ cm<sup>-3</sup>. In Section 4 we estimate the number density of cold protons to be less than this for reasonable values of the jet width. ### 3.2 Simulation of particle production and cascade development For the first time in the context of the SPB-model, we use the Monte-Carlo technique to simulate particle production and cascade development, and this allows us to use exact cross sections. For photomeson production we use the Monte-Carlo code SOPHIA (Mücke et al 2000), and Bethe-Heitler pair production is simulated using the code of Protheroe & Johnson (1996). We calculate the yields for both processes separately, and the results are then combined according to their relative interaction rates. The mean pion production interaction rate for an isotropic photon field is $$r_\pi (E_p)=\frac{1}{8E_p^2\beta _p}_{ϵ_{\mathrm{th}}}^{\mathrm{}}𝑑ϵ\frac{n(ϵ)}{ϵ^2}_{s_{p,\mathrm{th}}}^{s_{p,\mathrm{max}}}𝑑s_p(s_pm_p^2)\sigma _\pi (s_p),$$ (20) where $`s_p=m_p^2+2E_pϵ(1\beta _p\mathrm{cos}\theta _p)`$ is the center-of-momentum (CM) energy squared, $`\theta _p`$ the angle between the proton and the photon, $`\beta _pc`$ the proton’s velocity, $`\sqrt{s_{p,\mathrm{th}}}1.08`$GeV the threshold CM energy, $`ϵ_{\mathrm{th}}=(s_{p,\mathrm{th}}m_p^2)/2(E_p+p_p)`$, $`s_{p,\mathrm{max}}=m_p^2+2E_pϵ(1+\beta _p)`$ and $`\sigma _\pi `$ the pion production cross section. The Bethe-Heitler pair production interaction rate, $`R_{\mathrm{BH}}`$, is calculated using the formulae given in Chodorowski (1992). In highly magnetized environments, proton-photon interactions compete with synchrotron radiation by the protons. To take proton synchrotron losses into account in our code, we sample a $`p\gamma `$ interaction length from an exponential distribution with its corresponding mean $`\overline{x}_{p\gamma }=c/r_{\pi ,BH}`$ as given below. We make the approximation that the proton energy on interacting is given by $$E_pE_p^{(0)}\left[1+\frac{x_{p\gamma }}{\overline{x}_{p,syn}(E_p^{(0)})}\right]^1$$ (21) where $`E_p^{(0)}`$ is the initial proton energy, $`x_{p\gamma }`$ is the sampled interaction length for pion production or Bethe-Heitler pair production, and $`\overline{x}_{p,\mathrm{syn}}=c/r_{p,\mathrm{syn}}`$ is the proton synchrotron loss distance given by $$r_{M,\mathrm{syn}}=\frac{4}{3}\left(\frac{m_e}{M}\right)^2\frac{\sigma _T\gamma u_B}{Mc}.$$ (22) where $`M=m_p`$, $`\sigma _T`$ is the Thomson cross section, and $`u_B`$ is the magnetic energy density. The relatively long mean life time of highly energetic charged pions, muons and kaons might be of the same order of magnitude as their synchrotron loss time scale in highly magnetized environments like AGN jets and GRBs (Rachen & Mészáros 1998). We simulate their synchrotron energy losses by sampling the decay length $`x_{\mathrm{dec}}`$ from an exponential distribution with corresponding mean decay length $`\overline{x}_{\mathrm{dec}}=c\gamma \tau _{\mathrm{dec}}`$. For $`\mu ^\pm `$ we have $`\tau _{\mathrm{dec}}=2.20\times 10^6`$ s, while $`\tau _{\mathrm{dec}}=2.60\times 10^8`$ s for $`\pi ^\pm `$ and $`\tau _{\mathrm{dec}}=1.24\times 10^8`$ s for $`K^\pm `$. The particle’s energy on decaying is then $$E_{K,\pi ,\mu }=E_{K,\pi ,\mu }^{(0)}\left[1+\frac{2x_{dec}}{\overline{x}_{K,\pi ,\mu ,\mathrm{syn}}(E_{K,\pi ,\mu }^{(0)})}\right]^{1/2}.$$ (23) where $`E_{K,\pi ,\mu }^{(0)}`$ is the initial $`K`$,$`\pi `$, or $`\mu `$ energy, $`\overline{x}_{K,\pi ,\mu ,\mathrm{syn}}=c/r_{K,\pi ,\mu ,\mathrm{syn}}`$ is the synchrotron loss distance with $`r_{K,\pi ,\mu ,\mathrm{syn}}`$ given by Eq. 22 for $`M=m_K`$, $`m_\pi `$ or $`m_\mu `$. We next outline the simulations of the cascade development in more detail. Energetic photons will pair produce on the target photon field, and if the magnetic energy density exceeds the injected target field density ($`u_\mathrm{B}>u_{\mathrm{target}}`$) they will initiate a pair-synchrotron cascade. Energetic photons arise directly from $`\pi ^0`$-decay, or indirectly as synchrotron photons from protons, charged mesons and electrons resulting from $`\pi ^\pm \mu ^\pm e^\pm `$ decay and Bethe-Heitler pair production. The optical depth $`\tau _{\gamma \gamma }(E_\gamma )`$ for $`\gamma `$-ray photons with energy $`E_\gamma `$ for $`e^\pm `$ pair production inside the blob is given by $$\tau _{\gamma \gamma }(E_\gamma )=\frac{R_{\mathrm{blob}}}{8E_\gamma ^2}_{ϵ_{\mathrm{min}}}^{\mathrm{}}𝑑ϵ\frac{n(ϵ)}{ϵ^2}_{s_{\gamma ,\mathrm{min}}}^{s_{\gamma ,\mathrm{max}}(ϵ,E_\gamma )}𝑑s_\gamma s_\gamma \sigma _{\gamma \gamma }(s_\gamma )$$ (24) where $`n(ϵ)`$ is the differential photon number density and $`\sigma _{\gamma \gamma }(s)`$ is the total cross section for photon-photon pair production (Jauch & Rohrlich 1955) for a centre of momentum frame energy squared given by $$s_\gamma =2ϵE_\gamma (1\mathrm{cos}\theta _\gamma )$$ (25) where $`\theta _\gamma `$ is the angle between directions of the energetic photon and soft photon, and $`s_{\gamma ,\mathrm{min}}=(2m_ec^2)^2`$, $`s_{\gamma ,\mathrm{max}}=4ϵE_\gamma `$ and $`ϵ_{\mathrm{min}}=(m_ec^2)^2/E_\gamma `$. For simulating photon-photon pair production we approximate $`E_{e^+}=E_e^{}=E_\gamma /2`$. The radiating $`e^\pm `$ are assumed to be continuously isotropized in the blob frame by deflection in the uniform magnetic field, resulting in the synchrotron radiation being isotropic in the blob frame. The spectrum of synchrotron photons, averaged over pitch-angle because of the isotropy of the particle distribution, is calculated using functions given by Protheroe (1990). In general, the cascade can be initiated by photons from $`\pi ^0`$-decay (“$`\pi ^0`$ cascade”), electrons from the $`\pi ^\pm \mu ^\pm e^\pm `$ decay (“$`\pi ^\pm `$ cascade”), $`e^\pm `$ from the proton-photon Bethe-Heitler pair production (“Bethe-Heitler cascade”) and $`p`$ and $`\mu `$-synchrotron photons (“$`p`$-synchrotron cascade” and “$`\mu ^\pm `$-synchrotron cascade”). Here, we assume the cascades develop linearly, and this requires the photon field produced by the cascade, to be negligible as a target field in comparison with the injected synchrotron radiation due to co-accelerated $`e^{}`$. This requirement can be expressed as $`\tau _{\gamma \gamma ,\mathrm{cas}}\tau _{\gamma \gamma ,\mathrm{target}}`$ with $`\tau _{\gamma \gamma ,\mathrm{target}}`$,$`\tau _{\gamma \gamma ,\mathrm{cas}}`$ being the pair production optical depths of photons on the target and cascade photon field, respectively. Our Monte-Carlo results show this condition to be met for the present input. To simplify the calculation, the electrons are completely cooled instantly by synchrotron radiation before pair production by the synchrotron photons takes place. This approximation is equivalent to assuming that $`t_{\mathrm{syn}}t_{\mathrm{pair}}`$ which is justified because of the very short synchrotron life time of electrons in the assumed magnetic field. A matrix method (e.g. Johnson et al 1996) is then used to follow the pair-synchrotron cascade in the ambient synchrotron radiation field and magnetic field. The cascades are considered in the jet frame. Here, electron and photon fluxes are represented by vectors $`G_j^k`$ and $`F_i^k`$, which give the total number of electrons in the energy bin at energy $`E_j`$ and number of photons at energy $`E_i`$, respectively, in the $`k`$th cascade generation. We use a logarithmic stepsize of 0.1 ranging from $`\mathrm{log}(E_e/1\mathrm{GeV})=3`$ to $`12`$ in electron energy, and from $`\mathrm{log}(E_\gamma /1\mathrm{GeV})=13`$ to $`12`$ in photon energy. Averaged over a homogeneous emission region, the probability of gamma-ray interaction by photon-photon pair production at energy $`E_i`$ is given by the vector $`P_{\gamma \gamma ,\mathrm{i}}`$ $$P_{\gamma \gamma ,\mathrm{i}}=[1P_{\mathrm{esc}}(E_i)]=\left[1\frac{1\mathrm{exp}(\tau _{\gamma \gamma }(E_i))}{\tau _{\gamma \gamma }(E_i)}\right].$$ (26) The transfer matrix $`T_{ij}^{(\mathrm{syn})}`$ gives the number of synchrotron photons of energy $`E_i`$ produced by electrons of energy $`E_j`$, $`T_{ij}^{(\mathrm{syn})}`$, and the transfer matrix $`T_{ij}^{(\mathrm{pair})}`$ gives the number of $`e^\pm `$ of energy $`E_j`$ produced through pair production of energetic photons of energy $`E_i`$ on the target field. The vectors and matrices are calculated taking care of energy conservation. The photon and electron fluxes due to synchrotron radiation and photon-photon pair production are then calculated through matrix multiplication. We start by calculating the number of $`\pi ^0`$-decay photons which pair produce in the blob $$F_i^0=I_i^0P_{\gamma \gamma ,\mathrm{i}}$$ where the vector $`I_i^0`$ gives the number of $`\pi ^0`$-decay gamma-rays at an energy $`E_i`$ (the emerging photons, i.e. those which do not pair produce in the blob, are stored in an array). The electron spectrum due to photon-photon pair production is then given by $$F_j^1=\underset{i}{}F_i^0T_{ij}^{(\mathrm{pair})}.$$ These electrons radiate synchrotron photons, and the resulting photon yield is determined by $$I_i^1=\underset{j}{}F_j^1T_{ij}^{(\mathrm{syn})}.$$ This is the 1st cascade generation photon spectrum, which in turn again suffers photon-photon pair production on the target photon field in the blob, etc. In our calculation, we iterate until $`F_i^k0.01`$ at any energy $`E_i`$, or stop at 10 generations. We have injected $`10^4`$ protons at each proton energy equally spaced in $`\mathrm{log}E_p`$ at 0.1 decade intervals from $`10^3`$ GeV to the maximum injection energy. These protons are assumed to be continuously isotropized, and interact with the synchrotron radiation field of the co-accelerated $`e^{}`$ through pion production and Bethe-Heitler pair production. The resulting neutrinos (see Sect. 4.1) escape without further interaction, but the $`e^\pm `$ and high energy $`\gamma `$-rays initiate cascades. The resulting particle spectra were weighted with $`n_{p,0}E_p^2dE_p`$, appropriate to the assumed proton injection spectrum, and divided by the number of injected protons, and by $`4\pi `$ steradians. Figs. 4 and 5 show examples of cascade spectra in the observer’s frame initiated by photons with different origins for model parameters (given above), which satisfactorily reproduce the flare spectrum of Mkn 501. Our Monte-Carlo program only treats the first interactions of protons, and so we must take account of this when making the final flux predictions. In the case of pion production, a fraction $`\kappa 0.25`$ of the initial proton energy goes into particle production, and $`1/3`$ of the time the emerging nucleon is a neutron, and is assumed to escape from the blob without further interaction. Combining these two factors, the resulting spectra must be multiplied by 2.01 to take account of subsequent pion production interactions. Similarly, the multiplication factor which takes account of subsequent Bethe-Heitler pair production interactions is $`10`$. Particle spectra due to pion production and Bethe-Heitler pair production are then weighted according to their mean interaction rates, i.e. $`R_\pi /R_{\mathrm{tot}}`$ for pion production, and $`R_{\mathrm{BH}}/R_{\mathrm{tot}}`$ for Bethe-Heitler pair production where $`R_{\mathrm{tot}}=R_\pi +R_{\mathrm{BH}}+R_{\mathrm{ad}}`$ is the total interaction rate, taking into account adiabatic losses of the protons due to jet expansion approximated by Eq. 18. Adiabatic jet expansion also affects the radius of the emission region, its photon density and magnetic field, and thereby the probability for $`p\gamma `$-interactions and the subsequent cascade development. With the expansion of the emission region, the probability for $`p\gamma `$-interactions decreases because of the decrease of the photon energy density and magnetic field. It follows that $`p\gamma `$-interactions, and their resulting cascades, most likely take place during the initial phase of the jet expansion, and so we neglect any effects on the emerging cascade spectra due to changes in magnetic field or dimension of the emission region. In particular, the flare rise-time can be much shorter than the expansion and light crossing time scales if one associates it with the onset of the $`p\gamma `$-interactions and their resulting cascades. Short rise-times and longer flare decay time scales are typically observed in blazar light curves. Finally, we transform the cascade spectrum to the observer’s frame. For calculating the distance, $`q_0=0.5`$ and $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup> are used. Cascades initiated by photons from $`\pi ^0`$ decay, “$`\pi ^0`$ cascades”, or by electrons from $`\pi \mu e`$ decay, “$`\pi ^\pm `$ cascades”, (Fig. 4) produce rather featureless spectra (see also Mannheim et al 1993). However, cascades initiated by protons “$`p`$-synchrotron cascades”, and by muons from $`\pi \mu e`$ decay, “$`\mu ^\pm `$-synchrotron cascades” (Fig. 5b) produce a double-humped SED as observed for $`\gamma `$-ray blazars (see also Rachen 1999). The contribution from Bethe-Heitler pair production turns out to be negligible (Fig. 5a). In the model for Mkn 501 presented here, the synchrotron flux from directly accelerated $`e^{}`$ significantly exceeds the reprocessed $`\mu `$ and $`p`$ synchrotron flux. Thus, in the present Mkn 501 model direct proton and muon synchrotron radiation is mainly responsible for the high energy hump whereas the low energy hump is mainly synchrotron radiation by the directly accelerated $`e^{}`$. In the present model photons up to 1 TeV in the jet frame are optically thin to $`\gamma \gamma `$-pair production. Thus, only a small fraction of the power emitted as proton synchrotron radiation is redistributed to lower energies via cascading. For models where the optical depth of TeV-photons is significantly higher, a correspondingly larger fraction of the TeV-bump may be redistributed to X-ray energies. In these models, the pairs produced by photons of the “high energy hump” may contribute considerably to the observed X-ray bump, whereas in the present model they are negligible in comparison to the directly accelerated electrons. We note that for the case of strong magnetic fields where $`\mu ^\pm `$-synchrotron cascades cause detectable humps, the $`p`$ cascade spectrum dominates in general over the $`\mu ^\pm `$ cascade spectrum. ## 4 Application to the April 1997 flare of Mkn 501 Adding the four components of the cascade spectra in Fig. 4-5 we obtain the SED shown in Fig. 6 where it is compared with the multifrequency observations of the 16 April 1997 flare of Mkn 501. We use the parametrization of Bednarek & Protheroe (1999) to represent the BeppoSAX+OSSE data (thick straight line at low energies). The broken power-law simplification (Eq. 3.1) used as the target radiation field for $`p\gamma `$-interactions and $`\gamma \gamma `$ interactions is shown by the chain line. The 100 MeV upper limit from Catanese et al (1997) is nearly simultaneous to the TeV-flare observations. The TeV-flux, corrected for pair production on the cosmic background radiation field (Bednarek & Protheroe 1999) for two different background models, is shown as the thick curves. The parameters used for modeling the April 1997 flare are: $`D=12`$, $`B20`$ G, radius of the emission region $`R_{\mathrm{blob}}=8\times 10^{15}`$ cm. For a target photon field for the $`p\gamma `$-interactions, and the cascades as given in Eq. 16, we find a photon energy density of this radiation field of $`u_{\mathrm{target}}=60`$ GeV/cm<sup>-3</sup>. The accelerated protons are assumed to follow a power law $`\gamma _p^2`$ between $`2\gamma _p\gamma _{p,\mathrm{max}}=3\times 10^{10}`$, and in order to fit the emerging cascade spectra to the data a proton number density of $`n_p250`$ cm<sup>-3</sup>, corresponding to an energy density of accelerated protons of $`u_p11.6`$ TeV/cm<sup>-3</sup> is required. With a magnetic field energy density of $`u_B11.7`$ TeV/cm<sup>-3</sup> our model satisfies $`u_{\mathrm{target}}u_pu_B`$ (all parameters are in the co-moving frame of the jet), confirming that a significant contribution from inverse-Compton scattering is not expected. To calculate the total jet luminosity $`L_{\mathrm{jet}}`$, measured in the rest frame of the galaxy, we adapt the formulae of Bicknell (1994) and Bicknell & Dopita (1997) given for the synchrotron self-Compton model to apply for the case of our proton blazar model. We receive $`L_{\mathrm{jet}}`$ $`=`$ $`{\displaystyle \frac{L_{\mathrm{obs}}^{\mathrm{high}}}{D^2\zeta _p}}\left[\chi _p{\displaystyle \frac{(\mathrm{\Gamma }1)}{\mathrm{\Gamma }}}+1+{\displaystyle \frac{p_B}{p_p}}+{\displaystyle \frac{\zeta _pS_{\mathrm{obs}}^{\mathrm{low}}}{\zeta _eS_{\mathrm{obs}}^{\mathrm{high}}}}\right]`$ (27) where $`\mathrm{\Gamma }=(1\beta ^2)^1D/2`$ is a good approximation to the Lorentz factor of jets closely aligned to the line of sight, and the four terms inside the bracket give the relative contributions to the total jet power of cold protons, accelerated protons, magnetic field, and accelerated electrons, respectively. The contribution from cold protons is estimated on the basis of charge neutrality. $`S_{\mathrm{obs}}^{\mathrm{low}}`$ and $`S_{\mathrm{obs}}^{\mathrm{high}}`$ are the observed bolometric fluxes of the low and high energy component, respectively, and $`\zeta _e1`$, $`\zeta _p`$ are the radiative efficiencies for electrons and protons. $`L_{\mathrm{obs}}^{\mathrm{high}}=4\pi d^2S_{\mathrm{obs}}^{\mathrm{high}}`$ with $`d`$ the luminosity distance, $`p_B=[B^2/(2\mu _0)]/3`$ and $$p_p=\frac{L_{\mathrm{obs}}^{\mathrm{high}}}{4D^2\zeta _p\mathrm{\Gamma }^2\beta c\pi R_{blob}^2}$$ (28) gives the jet-frame pressure of relativistic protons that would apply in the absense of energy loss mechanisms, and $$\chi _p=\frac{3}{4}\left(\frac{m_p}{m_e}\frac{\zeta _pS_{\mathrm{obs}}^{\mathrm{low}}}{\zeta _eS_{\mathrm{obs}}^{\mathrm{high}}}\frac{1}{\gamma _{e}^{}{}_{1}{}^{}\mathrm{ln}(\gamma _{e}^{}{}_{2}{}^{}/\gamma _{e}^{}{}_{1}{}^{})}\frac{1}{\gamma _{p}^{}{}_{1}{}^{}\mathrm{ln}(\gamma _{p}^{}{}_{2}{}^{}/\gamma _{p}^{}{}_{1}{}^{})}\right).$$ (29) We use this formula to estimate the total jet power to be $`10^{46}`$ erg/s, and we find the contributions to the total jet power of cold protons, magnetic field, and accelerated electrons, relative to that of accelerated protons, to be 0.8, 1, and 0.01, respectively, with a number density of cold protons $`10^4`$ cm<sup>-3</sup>. Thus, for the shock acceleration mechanism we require more power per particle going into relativistic protons than into electrons. Fig. 7 shows the dependence of the total jet luminosity on the Doppler factor $`D`$ with fixed parameters $`t_{\mathrm{var}}=12`$ hours, $`\beta _1=0.5`$, and with $`B`$, $`R`$, $`n_p`$ and $`u_{\mathrm{target}}`$ being of Mkn 501 during flaring. Clearly visible is the fact that in hadronic models the proton kinetic energy and the poynting flux dominate the total jet luminosity, while the electron kinetic energy is only of minor importance. At high Doppler factors the emission region becomes so large that one needs only relatively small magnetic fields and proton densities to fit the observations. In addition, adiabatic losses become small resulting in a decrease of the required kinetic proton luminosity. For example, $`B5`$ G and $`n_p10^2`$ cm<sup>-3</sup> are sufficient to fit the Mkn 501 flare for $`D=50`$, while for $`D=8`$ magnetic fields of over 30 G and proton densities of $`n_p10^4`$ cm<sup>-3</sup> are needed. The total jet luminosity exhibits a minimum of $`10^{46}`$erg/s at around $`D12`$, which corresponds to the model we have chosen to present here. Smaller variability time scales in the shock acceleration model can only be explained by quasi-perpendicular shocks with a large mean free path (see Fig. 1a), which in turn are only consistent with the diffusion approximation for low shock speeds. Due to the size limit of the shock, the maximum energy gain is then limited to significantly lower values, and one may not reach TeV-photon energies in the proton synchrotron model, unless the magnetic field is increased accordingly, which in turn leads to larger jet luminosities. For example, for $`D=10`$ and a variability time scale of 12 hours, at least B=20 G is needed to comply with all acceleration constraints. For $`t_{\mathrm{var}}=3`$ hours, one needs at least 50 G. In comparison, leptonic models would give a minimum jet luminosity of $`10^{44}`$ erg/s for this flare (Ghisellini 1998), about two orders of magnitude lower. This is due to the much lower magnetic fields invoked there, and the less massive particles which drive the kinetic flow. In the framework of the jet–disk symbiosis (e.g. Falcke & Biermann 1995), the jet luminosity should not exceed the total accretion power $`Q_{\mathrm{accr}}`$ for the equilibrium state. Accretion theory relates the disk luminosity to the accretion power. Page & Thorne (1974) give $`L_{\mathrm{disk}}(0.050.3)Q_{\mathrm{accr}}`$. Disk luminosities for ’typical’ radio-loud AGN lie in the range $`L_{\mathrm{disk}}10^{44}10^{48}`$ erg/s with BL Lac objects tending to the lower end on average. Specifically, for Mkn 501 there are no emission line measurements available, and this complicates the evaluation of the disk luminosity for this object. However, any observed UV-emission in the flaring stage may put an upper limit on it. Again only historical data are available here, leaving room for speculation. We estimate $`L_{\mathrm{disk}}10^{43}10^{44}`$ erg/s (Mufson et al 1984, Pian et al 1998), and obtain for the accretion power, $`Q_{\mathrm{accr}}(3200)10^{43}`$ erg/s, at least a factor 5 below the necessary value to comply with the constraint of the disk–jet symbiosis. Note, however, that the estimate of $`Q_{\mathrm{accr}}`$ is based on archival non-flaring data from Mkn 501, and we could speculate that either the disk has pushed more energy into the jet during TeV-flaring, or that the flaring stage can not be considered as a steady state. Also, accretion theory might predict larger conversion efficiencies of the accretion power into disk radiation than actually might occur in BL Lac objects. It is also instructive to consider the radiative efficiency of the proposed model in comparison to alternative models. Using Eq. 29 we estimate the radiation efficiency of the protons we find $`\zeta _p0.01`$. This is similar to the value of $`\zeta _p`$ in the Proton-Blazar model proposed by Mannheim (1993), and is fully in agreement with the results presented by Celotti & Fabian (1993) on a basis of a sample of 105 sources. In comparison, leptonic models would give two orders of magnitude higher radiation efficiencies. ### 4.1 Neutrino spectra Unlike leptonic models, hadronic blazar models may result in neutrino emission through the production and decay of charged mesons, e.g. $`\pi ^+\mu ^++\nu _\mu `$ followed by $`\mu ^+e^++\nu _e+\overline{\nu }_\mu `$. Neutrinos escape without further interaction, and the predicted neutrino spectrum from Mkn 501 during flaring is shown in Fig. 8. We calculate the $`\nu `$-emission from the source itself, and do not include here any additional contribution from escaping cosmic rays interacting while propagating through the cosmic microwave background radiation (see e.g. Protheroe and Johnson 1996). The proton injection spectrum is modified by the photo-hadronic interaction rate which approximately follows a $`\gamma _p^{+0.6}`$ power-law for proton energies above $`10^5`$ GeV (see Fig. 3), where the nucleons interact preferably in the flatter part of the target photon spectrum ($`ϵ1.6`$ keV) to produce mesons. This causes the resulting neutrino spectrum above $`E_\nu 10^5`$ GeV to be $`dN/dE_\nu E_\nu ^{1.4}`$, whereas below this energy, the target photon field for pion production is the steeper part (1.6 keV $`ϵ42`$ keV), and the corresponding neutrino spectrum is $`dN/dE_\nu E_\nu ^{1.2}`$. At even lower energies a further flattening is due to pion production by the lowest energy protons at threshold. There is a steepening in the spectrum above $`10^9`$ GeV which needs some comment here. Neutrinos with energy $`<10^9`$ GeV are mostly produced near threshold in the $`\mathrm{\Delta }`$-resonance region, while the higher energy neutrinos are mainly produced in the secondary resonant region of the pion production cross section. $`\pi ^{}`$ production is suppressed in the $`\mathrm{\Delta }`$-resonance, and hence we find a suppression of $`\overline{\nu }_e`$-emission in comparison to $`\nu _e`$ emission at low energies, whereas $`\pi ^+`$ and $`\pi ^{}`$ multiplicities are comparable at higher CM energies, leading to roughly equal fluxes of $`\nu _e`$ and $`\overline{\nu }_e`$. The effects of $`\mu ^\pm `$-synchrotron emission show up as a break at $`10^9`$ GeV in the observer’s frame, whereas $`\pi ^\pm `$-synchrotron emission turns out to be unimportant in our model. Another important source of high energy neutrinos is the production and decay of charged kaons. In the case of $`p\gamma `$-interactions positively charged kaons are produced, and are important in the secondary resonance region of the cross section (Mücke et al 1999, Mücke et al 2000). They decay in $`64\%`$ of all cases into muons and direct high energy muon-neutrinos. In contrast to the neutrinos originating from $`\pi ^\pm `$ and $`\mu ^\pm `$-decay, these muon-neutrinos will not have suffered energy losses through $`\pi ^\pm `$\- and $`\mu ^\pm `$-synchrotron radiation, and therefore appear as an excess in comparison to the remaining neutrino flavors at the high energy end ($`E_\nu 10^9`$$`10^{10}`$ GeV) of the emerging neutrino spectrum, and in addition cause the total neutrino spectrum to extend to $`10^{10}`$ GeV. In contrast to previous SPB-jet models in which one expected equal photon and neutrino energy fluxes (e.g. Mannheim 1993, 1995), our model predicts a peak neutrino energy flux approximately two orders of magnitude lower than high energy gamma-rays. This is due to the synchrotron losses of protons being dominant, leading to gamma-ray emission at the expense of neutrino emission. ## 5 Conclusions This paper describes an application of our newly-developed Monte Carlo program which simulates a modified version of the stationary SPB-model. The Monte Carlo technique allows us to use exact cross sections, and all important emission processes are considered here. As an example, we have used our code to model the giant April 1997 TeV-flare from Mkn 501. Here, the TeV-photons are due to synchrotron radiation of the relativistic protons in the highly magnetized emission region. This proton synchrotron model was first proposed by Mücke & Protheroe (1999); a similar model for TeV emission in Mrk 501 has just been proposed by Aharonian (2000), and his conclusions regarding the required Doppler factor and magnetic field are very similar to ours. Our model departs from the standard SPB model as introduced first by Mannheim and co-workers mainly in two areas: (i) we use the observed synchrotron radiation as the target photon field for $`p\gamma `$-interactions and pair-synchrotron cascades assuming it to be produced by co-accelerated electrons, and (ii) our model takes into account synchrotron radiation from muons and protons. The model parameters derived assuming diffusive shock acceleration of $`e^{}`$ and $`p`$ in a Kolmogorov turbulence spectrum are consistent with the X-ray to TeV-data in the flare state. However, the total jet power we obtain is too large to comply with the steady-state jet–disk symbiosis scenario, but then we are not dealing with a stready-state phenomenon. While the emerging cascade SED initiated by $`\pi ^0`$ decay and $`\pi ^\pm `$ synchrotron photons turns out to be relatively featureless, as was also found by, e.g., Mannheim (1993), the $`\mu ^\pm `$ (see also Rachen 1999, and Rachen & Mannheim 2000) and, more importantly, the proton synchrotron radiation and its cascade produces a double-humped SED as is commonly observed in flaring blazars. For the present model, we find that proton synchrotron radiation dominates the TeV emission, while the contribution of the synchrotron radiation from the pairs, produced by photon-photon interactions of gamma-rays from the high energy hump is only minor. Our model considers the emission region to be homogeneous. Inhomogeneities in particle density, magnetic field, etc. within the source would result in a broader X-ray and TeV-peak in the SED. This indicates that for Mkn 501 a homogeneous model of the emission region seems to be appropriate. Being a hadronic model, our model predicts neutrino emission and we give the expected neutrino flux of Mrk 501 during flaring. Comparing our predicted neutrino flux with that for previous proton blazar models (e.g. Mannheim 1993), we find that the neutrino output in our model is significantly less than was previously estimated due to the synchrotron losses dominating the energy losses of protons, producing synchrotron $`\gamma `$-rays at the expense of $`\pi ^0`$ $`\gamma `$-rays and neutrinos. ## Acknowledgements We thank Jörg Rachen, Alina Donea and Geoff Bicknell for helpful discussions. This work was supported by a grant from the Australian Research Council.
warning/0004/hep-th0004165.html	ar5iv	text	# LPT-ORSAY 00/43LPTM-00/25hep-th/0004165 Brane solutions in strings with broken supersymmetry and dilaton tadpoles ## 1 Introduction Nonsupersymmetric string models are generically plagued by divergences which raise the question of their quantum consistency. In particular, ten dimensional nonsupersymmetric models have all dilaton tadpoles and some of them have tachyons in the spectrum. It is however believed that some tadpoles do not signal an internal inconsistency of the theory but merely a bakground redefinition . In particular, in orientifolds of Type II theories there should be a difference between tadpoles of Ramond-Ramond (RR) closed fields and tadpoles of Neveu-Schwarz Neveu-Schwarz (NS-NS) closed fields. While the first ones cannot be cured by a background redefinition and signal an internal inconsistency of the theory asking therefore always to be cancelled , the latter ones could in principle be cured by a background redefinition. The NS-NS tadpoles remove flat directions, generate potentials for the corresponding fields (for example dilaton in ten dimensions) and break supersymmetry. The difference between RR and NS-NS tadpoles play an important role in (some) orientifold models with broken supersymmetry recently constructed . The purpose of this letter is to explicitly find the background of the nonsupersymmetric tachyon-free strings in 10D: the type I model in ten dimensions , containing 32 D$`\overline{9}`$ antibranes and 32 O$`9_{}`$ planes, and the heterotic $`SO(16)\times SO(16)`$. For the two theories, there is no background with maximal $`SO(10)`$ Lorentz symmetry, a result to be expected by various considerations. We find explicitly the classical backgrounds with a 9D Poincaré symmetry. We find an unique solution for the Type I model and two independent solutions for the heterotic one. A remarkable feature (in the Type I and one of the two heterotic backgrounds) of the static solutions is that the tenth coordinate is dynamically compactified in the classical background. Furthemore, the effective nine-dimensional Planck and Yang-Mills constants are finite indicating that the low energy physics is nine-dimensional. Another classical background with maximal symmetry is a cosmological-type solution which we explicitly exihibit for the two theories. They both have big-bang type curvature singularities. In section 2 we briefly review the construction of the type I $`USp(32)`$ string. In section 3, we determine its classical background with maximal symmetry. In section 4, we consider the $`SO(16)\times SO(16)`$ heterotic string and finally we end in section 5 with a discussion of the solutions. ## 2 The Type I nonsupersymmetric $`USp(32)`$ string The unique supersymmetric Type I model in ten dimensions is based on the gauge group $`SO(32)`$ and contains, by using a modern language (see, for example, ) 32 D9 branes and 32 O$`9_+`$ planes. There is however another, nonsupersymmetric tachyon-free model, with the same closed string spectrum at tree-level, containing 32 $`\overline{D}`$9 branes (i.e. branes of positive tension and negative RR charge) and 32 O$`9_{}`$ planes (i.e. nondynamical objects with positive tension and positive RR charge). The open string partition functions are $`𝒦={\displaystyle \frac{1}{2(4\pi ^2\alpha ^{})^5}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\tau _2}{\tau _2^6}}(V_8S_8){\displaystyle \frac{1}{\eta ^8}},`$ $`𝒜={\displaystyle \frac{N^2}{2}}{\displaystyle \frac{1}{(8\pi ^2\alpha ^{})^5}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^6}}(V_8S_8){\displaystyle \frac{1}{\eta ^8}},`$ $`={\displaystyle \frac{N}{2}}{\displaystyle \frac{1}{(8\pi ^2\alpha ^{})^5}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^6}}(V_8+S_8){\displaystyle \frac{1}{\eta ^8}},`$ (2.1) where $`\alpha ^{}M_s^2`$ is the string tension and $$V_8=\frac{\theta _3^4\theta _4^4}{2\eta ^4},S_8=\frac{\theta _2^4\theta _1^4}{2\eta ^4},$$ (2.2) where $`\theta _i`$ are Jacobi functions and $`\eta `$ the Dedekind function (see, for example, ). In (2.1) the various modular functions are defined on the double covering torus of the corresponding (Klein, annulus, Möbius) surface of modular parameter $$\tau =2i\tau _2(\mathrm{Klein}),\tau =\frac{it}{2}(\mathrm{annulus}),\tau =\frac{it}{2}+\frac{1}{2}(\mathrm{Mobius}),$$ (2.3) where $`\tau =\tau _1+i\tau _2`$ is the modular parameter of the torus amplitude and $`t`$ is the (one-loop) open string modulus. As usual, the annulus and Möbius amplitudes have the dual interpretation of one-loop open string amplitudes and tree-level closed string propagation with the modulus $`l`$, related to the open string channel moduli by $$𝒦:l=\frac{1}{2\tau _2},𝒜:l=\frac{2}{t},:l=\frac{1}{2t}.$$ (2.4) ¿From the closed string propagation viewpoint, $`V_8`$ describe the NS-NS sector (more precisely, the dilaton) and $`S_8`$ the RR sector, corresponding to an unphysical 10-form. The tadpole conditions can be derived from the $`t0`$ ($`l\mathrm{}`$) limit of the amplitudes above and read $$𝒦+𝒜+=\frac{1}{2}\frac{1}{(8\pi ^2\alpha ^{})^5}_0^{\mathrm{}}𝑑l\{(N+32)^2\times 1(N32)^2\times 1\}+\mathrm{}.$$ (2.5) It is therefore clear that we can set to zero the RR tappole by choosing $`N=32`$, but we are forced to live with a dilaton tadpole. The resulting open spectrum is nonsupersymmetric (the closed spectrum is supersymmetric and given by the Type I supergravity) and contains the vectors of the gauge group $`USp(32)`$ and a fermion in the antisymmetric (reducible) representation. However, the spectrum is free of gauge and gravitational anomalies and therefore the model should be consistent. It is easy to realize from (2.1) that the model contains 32 $`\overline{D}`$9 branes and 32 O$`9_{}`$ planes , such that the total RR charge is zero but NS-NS tadpoles are present, signaling breaking of supersymmetry in the open sector. The effective action, identified by writing the amplitudes (2.1) in the tree-level closed channel, contains here the bosonic terms $$S=\frac{M_s^8}{2}d^{10}x\sqrt{G}e^{2\mathrm{\Phi }}[R+4(\mathrm{\Phi })^2]T_9d^{10}x[(N+32)\sqrt{G}e^\mathrm{\Phi }(N32)A_{10}]+\mathrm{},$$ (2.6) where $`T_9`$ is the D9 brane tension and we set to zero the RR two-form and the gauge fields, which will play no role in our paper. Notice in (2.6) the peculiar couplings of the dilaton and the 10-form to antibranes and O$`9_{}`$ planes, in agreement with the general properties displayed earlier. The RR tadpole $`N=32`$ is found in (2.6) simply as the classical field equation for the unphysical 10-form $`A_{10}`$. The difference between the supersymmetric $`SO(32)`$ and nonsupersymmetric $`USp(32)`$ model described previously is in the Möbius amplitude describing propagation between (anti)branes and orientifold planes. Indeed, in the nonsupersymmetric case there is a sign change in the vector (or NS-NS in the closed channel) character $`V_8`$. Both supersymmetric and nonsupersymmetric possibilites are however consistent with the particle interpretation and factorization of the amplitudes. As noticed before, the NS-NS tadpoles generate scalar potentials for the corresponding (closed-string) fields, in our case the (10d) dilaton. The dilaton potential read $$V(N+32)e^\mathrm{\Phi },$$ (2.7) and in the Einstein basis is proportional to $`(N+32)\mathrm{exp}(3\mathrm{\Phi }/2)`$. It has therefore the (usual) runaway behaviour towards zero string coupling, a feature which is of course true in any perturbative construction. The dilaton tadpole means that the classical background, around which we must consistently quantize the string, cannot be the ten dimensional Minkowski vacuum and solutions with lower symmetry must be searched for. Once the background is corectly identified, there is no NS-NS tadpole anymore, of course. ## 3 The classical background of the nonsupersymmetric $`USp(32)`$ Type I string We are searching for classical solutions of the effective lagrangian (2.6) of the model (2.1) in the Einstein frame, which reads $$S_E=\frac{1}{2k^2}d^{10}x\sqrt{G}[R\frac{1}{2}(\mathrm{\Phi })^2]T_9^Ed^{10}x[(N+32)\sqrt{G}e^{\frac{3\mathrm{\Phi }}{2}}(N32)A_{10}]+\mathrm{},$$ (3.1) where $`T_9^E`$ indicate that the tension here is in the Einstein frame. The maximal possible symmetry of the background of the model described in the previous paragraph has a nine dimensional Poincaré isometry and is of the following form $$ds^2=e^{2A(y)}\eta _{\mu \nu }dx^\mu dx^\nu +e^{2B(y)}dy^2,\mathrm{\Phi }=\mathrm{\Phi }(y),$$ (3.2) where $`\mu ,\nu =0\mathrm{}8`$ and the antisymmetric tensor field from the RR sector, the gauge fields and all fermion fields are set to zero. The Einstein and the dilaton field equations with this ansatz are given by $`36(A^{})^2+8A^{\prime \prime }8A^{}B^{}+{\displaystyle \frac{1}{4}}(\mathrm{\Phi }^{})^2=\alpha _Ee^{2B+3\mathrm{\Phi }/2},`$ $`36(A^{})^2{\displaystyle \frac{1}{4}}(\mathrm{\Phi }^{})^2=\alpha _Ee^{2B+3\mathrm{\Phi }/2},`$ $`\mathrm{\Phi }^{\prime \prime }+(9A^{}B^{})\mathrm{\Phi }^{}=3\alpha _Ee^{2B+3\mathrm{\Phi }/2},`$ (3.3) where we defined $`\alpha _E=(N+32)k^2T_9^E=64k^2T_9^E`$ and $`A^{}dA/dy`$, etc. The function $`B`$ can be gauge-fixed by using the reparametrisation invariance of the above equations. It is convenient to choose the coordinate $`y`$ where $`B=3\mathrm{\Phi }/4`$ so that the exponential factors in the equations (3.3) disappear. In this coordinate system, the second equation in (3.3) is solved in terms of one function<sup>1</sup><sup>1</sup>1It can be checked that the other possible sign choices in (3.4) lead to the same solution. $`f`$ $$A^{}=\frac{1}{6}\sqrt{\alpha _E}shf,\mathrm{\Phi }^{}=2\sqrt{\alpha _E}chf.$$ (3.4) The two other field equations become then $`{\displaystyle \frac{4}{3}}\sqrt{\alpha _E}f^{}chf+\alpha _Ee^{2f}`$ $`=`$ $`\alpha _E,`$ $`2\sqrt{\alpha _E}f^{}shf+{\displaystyle \frac{3}{2}}\alpha _Ee^{2f}`$ $`=`$ $`{\displaystyle \frac{3}{2}}\alpha _E,`$ (3.5) and the solution is then $$e^f=\frac{3}{2}\sqrt{\alpha _E}y+c,$$ (3.6) where $`c`$ is a constant. By a choice of the $`y`$ origin and rescaling of the $`x^\mu `$ coordinates, the final solution in the Einstein frame reads $`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{3}{4}}\alpha _Ey^2+{\displaystyle \frac{2}{3}}\mathrm{ln}\|\sqrt{\alpha _E}y\|+\mathrm{\Phi }_0,`$ $`ds_E^2`$ $`=`$ $`\|\sqrt{\alpha _E}y\|^{1/9}e^{\alpha _Ey^2/8}\eta _{\mu \nu }dx^\mu dx^\nu +\|\sqrt{\alpha _E}y\|^1e^{3\mathrm{\Phi }_0/2}e^{9\alpha _Ey^2/8}dy^2.`$ (3.7) For physical purposes it is also useful to display the solution in the string frame, related as usual by a Weyl rescaling $`Ge^{\frac{\mathrm{\Phi }}{2}}G`$ to the Einstein frame $$A_s=A+\frac{1}{4}\mathrm{\Phi },B_s=B+\frac{1}{4}\mathrm{\Phi }.$$ (3.8) In the string frame the solution reads $`g_s`$ $``$ $`e^\mathrm{\Phi }=e^{\mathrm{\Phi }_0}\|\sqrt{\alpha }y\|^{2/3}e^{3\alpha y^2/4},`$ $`ds^2`$ $`=`$ $`\|\sqrt{\alpha }y\|^{4/9}e^{\mathrm{\Phi }_0/2}e^{\alpha y^2/4}\eta _{\mu \nu }dx^\mu dx^\nu +\|\sqrt{\alpha }y\|^{2/3}e^{\mathrm{\Phi }_0}e^{3\alpha y^2/4}dy^2,`$ (3.9) where $`\alpha =64M_s^8T_9`$. The solution (3.9) displays two timelike singularities, one at the origin $`y=0`$ and one at infinity $`y=+\mathrm{}`$, so that the range of the $`y`$ coordinate is $`0<y<+\mathrm{}`$. The dilaton, on the other hand, vanishes at $`y=0`$ and diverges at $`y=+\mathrm{}`$. The brane solution found above (3.9) has a striking feature. Suppose the $`y`$ coordinate is noncompact $`0<y<+\mathrm{}`$. In curved space however, the real radius $`R_c`$ is given by the integral $$2\pi R_c=_0^{\mathrm{}}𝑑ye^B=e^{\mathrm{\Phi }_0/2}\alpha ^{\frac{1}{2}}_0^{\mathrm{}}\frac{du}{u^{1/3}}e^{3u^2/8},$$ (3.10) where $`u=\sqrt{\alpha }y`$. The result is finite, meaning that despite apparencies the tenth coordinate is actually compact. The topology of the solution is thus a ten dimensional manifold with two boundaries at $`y=0`$ and $`y=+\mathrm{}`$ that is $`R^9\times S^1/Z_2`$. Moreover, it can be argued that gravity and gauge fields of the D$`\overline{9}`$ branes are confined to the nine-dimensional noncompact subspace, by computing the nine-dimensional Planck mass and gauge couplings, respectively $`M_P^7`$ $`=`$ $`M_s^8{\displaystyle _0^{\mathrm{}}}𝑑ye^{7A_s+B_s2\mathrm{\Phi }}=M_s^8\alpha ^{\frac{1}{2}}e^{3\mathrm{\Phi }_0/4}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{du}{u^{1/9}}}e^{3u^2/4},`$ $`{\displaystyle \frac{1}{g_{YM}^2}}`$ $`=`$ $`M_s^6{\displaystyle _0^{\mathrm{}}}𝑑ye^{5A_s+B_s\mathrm{\Phi }}=M_s^6\alpha ^{\frac{1}{2}}e^{\mathrm{\Phi }_0/4}{\displaystyle _0^{\mathrm{}}}𝑑uu^{1/9}e^{u^2/2}.`$ (3.11) Both of them are finite, which indeed suggest that gravity and gauge fields of the D$`\overline{9}`$ branes are confined to the nine-dimensional subspace. The relations (3.10) and (3.11) are in sharp contrast with the usual flat space relations obtained by compactifying the ten-dimensional theory down to nine-dimensions on a circle of radius $`R_c`$ $$M_P^7e^{2\mathrm{\Phi }_0}R_cM_s^8,\frac{1}{g_{YM}^2}e^{\mathrm{\Phi }_0}R_cM_s^6.$$ (3.12) By using exactly the same method we find a cosmological solution for the $`USp(32)`$ nonsupersymmetric Type I model by searching a homogeneous metric of the form $$ds^2=e^{2B(t)}dt^2+e^{2A(t)}\delta _{\mu \nu }dx^\mu dx^\nu ,\mathrm{\Phi }=\mathrm{\Phi }(t),$$ (3.13) where $`t`$ is a time coordinate. The solution is easily found by following the steps which led to (3.7) and (3.9). The result in the string frame is $`g_s`$ $`=`$ $`e^\mathrm{\Phi }=e^{\mathrm{\Phi }_0}\|\sqrt{\alpha }t\|^{2/3}e^{3\alpha t^2/4},`$ $`ds^2`$ $`=`$ $`\|\sqrt{\alpha }t\|^{2/3}e^{\mathrm{\Phi }_0}e^{3\alpha t^2/4}dt^2+\|\sqrt{\alpha }t\|^{4/9}e^{\mathrm{\Phi }_0/2}e^{\alpha t^2/4}\delta _{\mu \nu }dx^\mu dx^\nu .`$ (3.14) The metric has a spacelike curvature singularities at $`t=0`$ and $`t=+\mathrm{}`$. The laps separating these two singularities, to be interpreted as the real time parameter, $$\tau =𝑑t\|\sqrt{\alpha }t\|^{1/3}e^{\mathrm{\Phi }_0/2}e^{3\alpha t^2/8},$$ (3.15) is infinite. ## 4 Classical background of the $`SO(16)\times SO(16)`$ heterotic string In ten dimensions there is a unique tachyon-free non-supersymmetric heterotic string model . It can be obtained from the two supersymmetric heterotic strings as a $`Z_2`$ orbifold. The resulting bosonic spectrum comprises the gravity multiplet: graviton, dilaton, antisymmetric tensor and gauge bosons of the gauge group $`SO(16)\times SO(16)`$. Compactifications of this theory were considered in , and its strong coupling behavior in nine dimensions was examined in . It has been shown that the cosmological constant (the partition function on the torus) at one loop is finite and positive, furthemore its approximate value is given by $$\mathrm{\Lambda }M_s^{10}\frac{2^6}{(2\pi )^{10}}\times 5.67.$$ (4.1) The effective low energy action for the gravity multiplet is the same as before except for the cosmological constant term, which now reads $$\mathrm{\Lambda }\sqrt{G}$$ (4.2) in the string metric. The absence of the dilaton in this term reflects the one-loop nature of the cosmological constant<sup>2</sup><sup>2</sup>2A direct computation shows also that the one-loop dilaton tadpole is non-zero and proportional to the cosmological constant $`\mathrm{\Lambda }`$.. The same ansatz of the Einstein metric as in the previous paragraph leads, in the Einstein frame, to the equations $`36(A^{})^2+8A^{\prime \prime }8A^{}B^{}+{\displaystyle \frac{1}{4}}(\mathrm{\Phi }^{})^2=\beta _Ee^{2B+5\mathrm{\Phi }/2},`$ $`36(A^{})^2{\displaystyle \frac{1}{4}}(\mathrm{\Phi }^{})^2=\beta _Ee^{2B+5\mathrm{\Phi }/2},`$ $`\mathrm{\Phi }^{\prime \prime }+(9A^{}B^{})\mathrm{\Phi }^{}=5\beta _Ee^{2B+5\mathrm{\Phi }/2},`$ (4.3) where we defined $`\beta _E=\mathrm{\Lambda }_Ek^2`$. The gauge which eliminates the exponential factors is now $`B=5\mathrm{\Phi }/4`$. After solving the second equation in (4.3) $`A^{}=\sqrt{\beta _E}sh(h)/6`$, $`\mathrm{\Phi }^{}=2\sqrt{\beta _E}ch(h)`$, the remaining equations give $$e^h=\frac{1}{2}\frac{e^{\sqrt{\beta _E}y}+ϵe^{\sqrt{\beta _E}y}}{e^{\sqrt{\beta _E}y}ϵe^{\sqrt{\beta _E}y}},$$ (4.4) where $`ϵ=\pm 1`$. An important difference with respect to the type I solution is that here we have two non-equivalent (that is, not related by coordinate transformations) solutions corresponding to $`ϵ=1`$ or $`1`$. Let us first consider the $`ϵ=1`$ case. The solution in the Einstein frame reads $`\mathrm{\Phi }`$ $`=`$ $`\mathrm{\Phi }_0+{\displaystyle \frac{1}{2}}\mathrm{ln}\|sh\sqrt{\beta _E}y\|+2\mathrm{ln}(ch\sqrt{\beta _E}y),`$ $`ds_E^2`$ $`=`$ $`\|sh(\sqrt{\beta _E}y)\|^{\frac{1}{12}}\left(ch(\sqrt{\beta _E}y)\right)^{\frac{1}{3}}dx^2+e^{5\mathrm{\Phi }_0/2}\|sh(\sqrt{\beta _E}y)\|^{\frac{5}{4}}\left(ch(\sqrt{\beta _E}y)\right)^5dy^2`$ (4.5) and in the string frame $`e^{2\mathrm{\Phi }}`$ $`=`$ $`e^{2\mathrm{\Phi }_0}\|sh(\sqrt{\beta }y)\|\left(ch(\sqrt{\beta }y)\right)^4,`$ $`ds^2`$ $`=`$ $`e^{\mathrm{\Phi }_0/2}\|sh(\sqrt{\beta }y)\|^{\frac{1}{3}}\left(ch(\sqrt{\beta }y)\right)^{\frac{2}{3}}dx^2+e^{2\mathrm{\Phi }_0}\|sh(\sqrt{\beta }y)\|^1\left(ch(\sqrt{\beta }y)\right)^4dy^2,`$ (4.6) where here $`\beta =\mathrm{\Lambda }M_s^8`$. In both metrics, the solution has two timelike singularities at $`y=0`$ and $`y=\mathrm{}`$. These singularities are separated by a finite distance which in the string frame reads $$2\pi R_c=(\beta )^{1/2}e^{\mathrm{\Phi }_0}_0^+\mathrm{}𝑑u(shu)^{1/2}(chu)^2.$$ (4.7) The spacetime has therefore the topology of a nine dimensional Minkowski space times an interval. Notice that the nine dimensional Planck mass $$M_p^7=M_s^8(\beta )^{1/2}e^{5\mathrm{\Phi }_0/4}_0^{\mathrm{}}𝑑u\left(shu\right)^{1/3}\left(chu\right)^{11/3},$$ (4.8) as well as the 9D Yang-Mills coupling $$\frac{1}{g_{YM}^2}=M_s^6e^{7\mathrm{\Phi }_0/4}(\beta )^{1/2}_0^{\mathrm{}}𝑑u\left(shu\right)^{2/3}\left(chu\right)^{13/3},$$ (4.9) are finite. This fact together with the finitude of the length of the tenth coordinate indicate that the low energy processes are described by a 9D theory. The coupling constant vanishes at $`0`$ and becomes infinite at large $`y`$. We shall comment more on this fact in the conclusion. The second solution corresponding to $`ϵ=1`$ can be obtained by exchanging $`(ch(\sqrt{\beta }y)`$ with $`\|sh(\sqrt{\beta }y\|`$ in the solutions (4.5) and (4.6). This solution has two singularities at $`0`$ and $`\mathrm{}`$, however the nine-dimensional Planck and Yang-Mills constants as well as the length of the tenth coordinate are infinite. The cosmological solution invariant with respect to the nine dimensional Euclidian group can be also readily found and it reads in the string frame $`ds^2`$ $`=`$ $`e^{2\mathrm{\Phi }_0}(\mathrm{sin}\sqrt{\beta }t)^1(\mathrm{cos}\sqrt{\beta }t)^4dt^2+e^{\mathrm{\Phi }_0/2}(\mathrm{sin}\sqrt{\beta }t)^{1/3}(\mathrm{cos}\sqrt{\beta }t)^{2/3}dx^2,`$ $`e^{2\mathrm{\Phi }}`$ $`=`$ $`e^{2\mathrm{\Phi }_0}(\mathrm{sin}\sqrt{\beta }t)(\mathrm{cos}\sqrt{\beta }t)^4.`$ (4.10) The variable $`t`$ in these equations belongs to the interval $`[0,\pi /(2\sqrt{\beta })]`$. At the boundaries in $`t=0`$ and $`t=\pi /(2\sqrt{\beta })`$ the metric develops curvature singularities. The time separating these two singularities is infinite: $$\tau =_0^{\pi /(2\sqrt{\beta })}𝑑t(\mathrm{sin}\sqrt{\beta }t)^{1/2}(\mathrm{cos}\sqrt{\beta }t)^2=\mathrm{}.$$ (4.11) Notice that the solution obtained by exchanging the sine and cosine in the above equations is not a new one since it can be obtained by shiting the time coordinate $`t\pi /(2\sqrt{\beta })t`$. ## 5 Discussion Before discussing the solutions we have found we should mention that there exists an another interesting non supersymmetric and tachyon-free model in ten dimensions : it is an orientifold of the type 0B string with a projection that removes the tachyon and introduces an open sector with the gauge group $`U(32)`$. This model has both a one loop (positive) cosmological constant and a disk dilaton tadpole so we have a sum of two terms in the low energy effective action $$\mathrm{\Lambda }_1\sqrt{G}e^\mathrm{\Phi }\mathrm{\Lambda }_2\sqrt{G},$$ (5.1) in the string metric. The first term is of the type we encountered in the $`USp(32)`$ case and the second is analogous to the one we met in the $`SO(16)\times SO(16)`$ case. There is no simple gauge choice for $`B`$ that renders the equations as simple as before. However qualitatively the solution should behave as the $`USp(32)`$ case for small $`y`$ where the coupling constant is small and the behavior for large $`y`$ should resemble that of the $`SO(16)\times SO(16)`$ case. In partcular we expect the radius of the tenth dimension to be compactified (since the divergence of the radius in the second solution of $`SO(16)\times SO(16)`$ is due to the behavior at the origin). We have determined the maximally symmetric solutions to the low energy equations of two tachyon-free non-supersymmetric strings. To what extent can we consider these solutions as representing the vacuum of these string theories ? Perturbatively, there are two kind of string corrections to the low energy effective action. The first ones are $`\alpha ^{}`$ corrections involving the string oscillators and the second ones are $`g_s`$ string loop corrections. A common feature of the solutions we found is that the conformal factor $`e^{2A}`$ (in nine dimensions) as well as the string coupling vanish at the origin and diverge at $`y=\mathrm{}`$. The effective string scale at coordinate $`y`$ being given by $`M_s^2(y)=M_s^2e^{2A_s}`$, we expect $`\alpha ^{}`$ corrections to be important at the origin and the loop corrections to be dominant at infinity. So strictly speaking we cannot trust the classical solution near the two singularities where interesting string physics would occur. Another less ambitious question, is the classical stability of our solutions. That is, do small perturbations around the background we found destroy the solution ? The answer to this question is closely related to the determination of the Kaluza-Klein excitations . Another common feature of the static solution of the Type I model (3.9) and of the first heterotic solution (4.6) is that the effective Yang-Mills and Planck constants are finite which means that the gravitational and gauge physics is effectively nine-dimensional. A remarkable feature of the static solutions in the low energy approximation is the spontaneous compactification of one coordinate. Whether this feature will survive the string and loop corrections is an interesting and open question. The fact that the nine dimensional metric is flat in spite of the ten-dimensional cosmological constant is in the spirit of the higher dimensional mechanisms which try to explain the vanishing of the (effective) cosmological constant . As in the new approaches to this problem discussed in , however, a better understanding of the naked singularities present in our solutions is needed in order to substantiate this claim. A notable difference between the type I and heterotic solutions is that the type I background is unique, whereas there are two classical heterotic backgrounds with physically very different properties. It is possible that this is due to the low energy approximation and that this degeneracy will be lifted by string corrections. All of the nontrivial features of the solutions are due to the presence of dilaton tadpoles. The latter are generic for non-supersymmetric string models. According to the Fischler-Susskind mechanism, the quantization around the classical solution should lead to finite string amplitudes . It would be intersting to confirm this explicitly for the present models. Acknowledgments We grateful to C. Angelantonj and A. Sagnotti for illuminating discussions on nonsupersymmetric strings and to C. Grojean and S. Lavignac for discussions concerning naked singularities in relation with the proposals . E.D. would like to thank the Theory Group at LBNL–Berkeley for warm hospitality during the final stage of this work.
warning/0004/hep-th0004147.html	ar5iv	text	# Contents ## 1 Introduction and summary The foundations of noncommutative geometry are inspired in large part by the fundamental principles of quantum mechanics. Just as in quantum mechanics physical observables are replaced with operators, and thereby satisfy the Heisenberg uncertainty principle, in noncommutative geometry spacetime coordinates are replaced with noncommuting operators. The corresponding smearing of spacetime coordinates in this way fits nicely into the ideas behind spacetime uncertainty relations and the concept of minimum length in string theory . This heuristic correspondence has been used to suggest that noncommutative geometry provides a natural framework to describe nonperturbative aspects of string theory . This belief is further supported by the fact that Matrix Theory and the IIB matrix model , which are conjectured to provide nonperturbative definitions of string theories, give rise to noncommutative Yang-Mills theory on toroidal compactifications . The particular noncommutative toroidal compactification is interpreted as being the result of the presence of a background Neveu-Schwarz two-form field, and it can also be understood in the context of open string quantization in D-brane backgrounds . Furthermore, in Ref. it has been shown that the IIB matrix model with D-brane backgrounds is described by noncommutative Yang-Mills theory. The early motivation for studying quantum field theory on noncommutative spacetimes was that, because of the spacetime uncertainty relation, the introduction of noncommutativity would provide a natural ultraviolet regularization. However, more recent perturbative calculations have shown that planar noncommutative Feynman diagrams contain exactly the same ultraviolet divergences that their commutative counterparts do, which implies that the noncommutativity does not serve as an ultraviolet regulator. One therefore needs to introduce some other form of regularization to study the dynamics of noncommutative field theories. On the other hand, it has been found that the ultraviolet divergences in non-planar Feynman diagrams exhibit an intriguing mixing of ultraviolet and infrared scales, which can also be described using string-theoretical approaches . Heuristically, this UV/IR mixing can be understood in terms of the induced uncertainty relations among the spacetime coordinates. If one measures a given spacetime coordinate with some high precision, then the remaining spacetime directions will generally extend because of the smearing. Furthermore, noncommutative solitons which do not have counterparts in ordinary field theory have been discovered for sufficiently large values of the noncommutativity parameters, and it has also been shown that noncommutative Yang-Mills theory in four dimensions naturally includes gravity. In order to investigate further the non-trivial dynamics of noncommutative field theories, it is important therefore to develop a nonperturbative regularization of these theories. Such a program has been put forward in Refs. , and it is similar to earlier works based on the mapping between large $`N`$ matrices and spacetime fields. In particular, in Ref. a unified framework was presented which naturally interpolates between the two ways that noncommutative Yang-Mills theory has appeared in the context of matrix model formulations of string theory, namely the compactification of Matrix theory and the twisted large $`N`$ reduced model. The model proposed was a finite $`N`$ matrix model defined by the twisted Eguchi-Kawai model with a quotient condition analogous to the ones considered in Refs. . It was interpreted as a lattice formulation of noncommutative gauge theory with manifest star-gauge invariance. The formulation naturally includes Wilson lattice gauge theory for a particular choice of parameters even at finite $`N`$. In Ref. , the lattice formulation was reconsidered from a general point of view without specifying any representation of the noncommutative algebra generated by the spacetime coordinates. This enabled an extension of the formalism to arbitrary even dimensions and also to allow for the minimal coupling to matter fields in the fundamental representation of the gauge group. Furthermore, in this approach one produces finite dimensional representations of the noncommutative geometry in much the same spirit as Weyl’s finite version of quantum mechanics . Discrete versions of noncommutative geometry using random lattices and graphs, as well as their applications to finite dimensional versions of noncommutative gauge theories, have also been studied in Ref. . In this paper we will show further that noncommutative Yang-Mills theory in an arbitrary even number of spacetime dimensions, with either rational or irrational valued (dimensionless) noncommutativity parameters and with gauge fields of arbitrary topological charge, can be regularized nonperturbatively by using commutative lattice gauge theories with twisted boundary conditions on the gauge fields . We will do so by presenting an extensive description of the lattice formulation of noncommutative gauge theories which was reported in Ref. . The first striking fact we will find in the general lattice formulation of noncommutative field theories is that the discretization of spacetime inevitably requires that it be compactified, in order to satisfy certain consistency conditions on the noncommutative algebra generated by the spacetime coordinates. It then follows that the lattice spacing must be much smaller than the length scale determined by the noncommutativity parameters, and so the continuum limit does not commute with the commutative limit. This demonstates the UV/IR mixing mentioned above at a completely nonperturbative level. We will also show that the dimensionless noncommutativity parameters are necessarily rational-valued on the lattice, which leads immediately to the possibility of obtaining finite dimensional representations of the noncommutative algebra as alluded to above. Irrational noncommutativity parameters can then be obtained by taking the continuum limit appropriately. Using the lattice formulation, we can regularize noncommutative Yang-Mills theories with periodic boundary conditions on the gauge fields, arbitrary noncommutativity parameters and arbitrary spacetime periodicities. We will then demonstrate on the lattice the Morita equivalence between ordinary, commutative Yang-Mills theory with twisted boundary conditions on the gauge fields and noncommutative Yang-Mills theory with periodic gauge fields. Using this equivalence, we will find that noncommutative Yang-Mills theory with single-valued gauge fields can actually be regularized by means of commutative lattice gauge theory with background magnetic flux. In order to obtain a continuum field theory with irrational-valued dimensionless noncommutativity parameters, the rank of the gauge group of the regulating commutative theory must be taken to infinity as one takes the continuum limit. As a special case, the construction includes the recent proposal of using twisted large $`N`$ reduced models for a concrete definition of noncommutative gauge theory. We shall find, however, that the class of noncommutative Yang-Mills theories that one can obtain using twisted large $`N`$ reduced models is not the most general one. The formalism also allows for the coupling of gauge fields to matter fields in the fundamental representation of the gauge group. We will carry out a large mass expansion of the matter-coupled theory and clarify the physical meaning of Wilson loops in noncommutative gauge theories . A remarkable property of these Wilson loops is that star-gauge invariance does not require that the loop be closed, but it does require that the relative distance vector between the two ends of the loop be proportional to its total momentum. Namely, as one increases the total momentum, the loop extends in spacetime. This is another manifestation of UV/IR mixing due to the noncommutativity. We will demonstrate how noncommutative Wilson loops naturally arise from matter field averages in noncommutative gauge theory. By calculating the effective action induced by matter and star-gauge invariant observables constructed out of the matter fields, we shall see that noncommutative Wilson loops play a very fundamental role, just like ordinary Wilson loops do in commutative gauge theories . We will also show explicitly how star-gauge invariant observables reduce smoothly to ordinary gauge invariant observables in the commutative limit for fixed gauge backgrounds. Morita equivalence in the presence of fundamental matter fields is also proven. As a special case, we obtain the twisted Eguchi-Kawai model with fundamental matter which was introduced in Ref. as a model which reproduces ordinary large $`N`$ gauge theory in the Veneziano limit. These Morita equivalences clarify the interpretations of various quantities in noncommutative gauge theory. For instance, it is known that only planar noncommutative Feynman diagrams contribute to the beta-function at one-loop order , and in the case of noncommutative quantum electrodynamics it is given by $$\beta (g^2)=\frac{g^4}{4\pi ^2}\left(\frac{11}{3}\frac{2}{3}n_f\right),$$ (1.1) where $`n_f`$ is the number of fermion flavours. The beta-function (1.1) coincides with that of ordinary $`\mathrm{U}(N)`$ Yang-Mills theory coupled to $`N_f=n_fN`$ flavours of fermion fields in the large $`N`$ (planar) limit (after an appropriate rescaling of the Yang-Mills coupling constant $`g`$). This coincidence can be understood as a consequence of the Morita equivalence that we shall derive. In fact, the results of such calculations indicate that the phenomenon of Morita equivalence in noncommutative gauge theories holds beyond the classical level in regularized perturbation theory. We will further show that the star-gauge invariant observables of two Morita equivalent noncommutative Yang-Mills theories are in a one-to-one correspondence. Owing to this fact, we can define regularized correlation functions of star-gauge invariant observables in noncommutative Yang-Mills theory with multi-valued gauge fields by using the lattice regularization of the corresponding Morita equivalent noncommutative gauge theory with periodic gauge fields. This Morita equivalence property will be the key feature which permits the discretization of generic noncommutative Yang-Mills theories in the continuum. The organization of the remainder of this paper is as follows. In Section 2 we present a pedagogical introduction to noncommutative field theory, and in particular noncommutative Yang-Mills theory, for the sake of completeness and in order to set up the notations to be used throughout the paper. In Section 3, we present a detailed and very general field theoretical derivation of the Morita equivalence relation between noncommutative gauge theories and demonstrate the one-to-one correspondence between star-gauge invariant observables in the Morita equivalent theories. In Section 4, we construct the discrete version of noncommutative field theory. We show that the lattice formulation exhausts the noncommutative Yang-Mills theories in arbitrary even dimensions, given the Morita equivalence properties described in Section 3. In Section 5 we establish Morita equivalence on the lattice, which allows us to regularize noncommutative Yang-Mills theories by means of commutative lattice gauge theories with twisted boundary conditions on the gauge fields. We also describe the relationship to finite dimensional matrix model representations of noncommutative gauge theories. Finally, in Section 6 we introduce matter fields in the fundamental representation of the gauge group and study the properties of star-gauge invariant observables as well as Morita equivalence in the presence of matter fields. ## 2 Quantum field theory on noncommutative spaces In this Section we will briefly review some aspects of constructing noncommutative gauge theories in the continuum, in a way that will enable us to later construct a lattice version of the field theory. We will start with a simple scalar field theory to introduce the ideas, and then move on to a description of noncommutative Yang-Mills theory and its observables. ### 2.1 Scalar field theory Consider the scalar quantum field theory which is described by the partition function $`Z`$ $`=`$ $`{\displaystyle 𝒟\varphi \text{e}^{S[\varphi ]}}`$ $`S`$ $`=`$ $`{\displaystyle \text{d}^Dx\left(\frac{1}{2}\left(_\mu \varphi \right)^2+\frac{1}{2}\varphi ^2+\frac{1}{4\text{!}}\varphi ^4\right)},`$ (2.1) where $`\varphi (x)`$ is a real-valued scalar field on $`D`$-dimensional Euclidean spacetime $`\text{}^D`$. For fields in a Schwartz space of functions of sufficiently rapid decrease at infinity, we may use the Fourier transformation $$\varphi (x)=\frac{\text{d}^Dk}{(2\pi )^D}\stackrel{~}{\varphi }(k)\text{e}^{ik_\mu x_\mu },$$ (2.2) with $`\stackrel{~}{\varphi }(k)=\stackrel{~}{\varphi }(k)^{}`$. The first step in generalizing a quantum field theory on an ordinary spacetime to one on a noncommutative spacetime is to replace the local coordinates $`x_\mu `$ by hermitian operators $`\widehat{x}_\mu `$ obeying the commutation relations $$[\widehat{x}_\mu ,\widehat{x}_\nu ]=i\theta _{\mu \nu },$$ (2.3) where $`\theta _{\mu \nu }=\theta _{\nu \mu }`$ are real-valued c-numbers with dimensions of length squared. Consequently, fields on spacetime are replaced by operators. Replacing $`x_\mu `$ in (2.2) by $`\widehat{x}_\mu `$, we obtain $$\widehat{\varphi }=\frac{\text{d}^Dk}{(2\pi )^D}\stackrel{~}{\varphi }(k)\text{e}^{ik_\mu \widehat{x}_\mu }.$$ (2.4) The operator-ordering ambiguity which exists when we replace $`x_\mu `$ by $`\widehat{x}_\mu `$ is fixed in (2.4) by requiring covariance and hermiticity of the operator $`\widehat{\varphi }`$. The operator (2.4) is called a Weyl operator or the Weyl symbol of the field $`\varphi (x)`$, and this method of constructing a noncommutative field theory is sometimes refered to as Weyl quantization. Combining (2.2) and (2.4), one can obtain an explicit map $`\widehat{\mathrm{\Delta }}(x)`$ which transforms a field $`\varphi (x)`$ to an operator $`\widehat{\varphi }`$ as $`\widehat{\varphi }`$ $`=`$ $`{\displaystyle \text{d}^Dx\varphi (x)\widehat{\mathrm{\Delta }}(x)}`$ (2.5) $`\widehat{\mathrm{\Delta }}(x)`$ $`=`$ $`{\displaystyle \frac{\text{d}^Dk}{(2\pi )^D}\text{e}^{ik_\mu \widehat{x}_\mu }\text{e}^{ik_\mu x_\mu }}.`$ (2.6) This means that the field $`\varphi (x)`$ can be thought of as the coordinate space representation of the operator $`\widehat{\varphi }`$. Note that the operator (2.6) is hermitian, $`\widehat{\mathrm{\Delta }}(x)^{}=\widehat{\mathrm{\Delta }}(x)`$, and in the commutative case $`\theta _{\mu \nu }=0`$ it reduces to a delta-function $`\delta ^D(\widehat{x}x)`$. However, for $`\theta _{\mu \nu }0`$, $`\widehat{\mathrm{\Delta }}(x)`$ is a complicated map. We may also introduce an anti-hermitian derivation $`\widehat{}_\mu `$ through the commutation relations $$[\widehat{}_\mu ,\widehat{x}_\nu ]=\delta _{\mu \nu },[\widehat{}_\mu ,\widehat{}_\nu ]=ic_{\mu \nu },$$ (2.7) where $`c_{\mu \nu }`$ are real-valued c-numbers.<sup>1</sup><sup>1</sup>1The c-numbers $`c_{\mu \nu }`$ turn out to be irrelevant for our purposes, since the operator $`\widehat{}_\mu `$ will only appear in commutator brackets. For this reason it is conventional to set $`c_{\mu \nu }=0`$. However, this restriction is not necessary, and indeed we will encounter cases with non-vanishing $`c_{\mu \nu }`$ later on. One can show that $`[\widehat{}_\mu ,\widehat{\varphi }]`$ $`=`$ $`{\displaystyle \text{d}^Dx_\mu \varphi (x)\widehat{\mathrm{\Delta }}(x)}`$ $`[\widehat{}_\mu ,\widehat{\mathrm{\Delta }}(x)]`$ $`=`$ $`_\mu \widehat{\mathrm{\Delta }}(x).`$ (2.8) From (2.8) it follows that any translation generator can be represented by the operator $`\text{e}^{v_\mu \widehat{}_\mu }`$, $`v_\mu \text{}`$, which satisfies $$\text{e}^{v_\mu \widehat{}_\mu }\widehat{\mathrm{\Delta }}(x)\text{e}^{v_\mu \widehat{}_\mu }=\widehat{\mathrm{\Delta }}(x+v).$$ (2.9) The existence of such an operator implies that $`\text{Tr }\widehat{\mathrm{\Delta }}(x)`$ is independent of $`x`$ for any trace Tr on the algebra of operators. Therefore, one can represent an integration of fields on the spacetime as $$\text{Tr }\widehat{\varphi }=\text{d}^Dx\varphi (x),$$ (2.10) where the normalization of the operator trace is fixed by requiring $`\text{Tr }\widehat{\mathrm{\Delta }}(x)=1`$. Using (2.5), (2.6) and (2.10) it is straightforward to show that the collection of operators $`\widehat{\mathrm{\Delta }}(x)`$ for $`x\text{}^D`$ form an orthonormal set, $$\text{Tr }\left(\widehat{\mathrm{\Delta }}(x)\widehat{\mathrm{\Delta }}(y)\right)=\delta ^D(xy),$$ (2.11) from which it follows that the inverse of the map (2.6) is given by $$\varphi (x)=\text{Tr }\left(\widehat{\varphi }\widehat{\mathrm{\Delta }}(x)\right).$$ (2.12) Therefore, there is a one-to-one correspondence between fields (of sufficiently rapid decrease at infinity) and operators. This can be thought of as an analog of the operator-state correspondence of local quantum field theory. Using these definitions, one can define a scalar quantum field theory on a noncommutative spacetime as $`Z`$ $`=`$ $`{\displaystyle \text{d}\widehat{\varphi }\text{e}^{S\left[\widehat{\varphi }\right]}}`$ $`S\left[\widehat{\varphi }\right]`$ $`=`$ $`\text{Tr }\left({\displaystyle \frac{1}{2}}[\widehat{}_\mu ,\widehat{\varphi }]^2+{\displaystyle \frac{1}{2}}\widehat{\varphi }^2+{\displaystyle \frac{1}{4\text{!}}}\widehat{\varphi }^4\right),`$ (2.13) where the measure $`\text{d}\widehat{\varphi }`$ is nothing but the ordinary path integral measure $`𝒟\varphi `$. The difference between noncommutative quantum field theory and ordinary field theory comes from the different products of fields that are used in the two cases. Suppose that $`\widehat{\varphi }_3=\widehat{\varphi }_1\widehat{\varphi }_2`$, where $`\widehat{\varphi }_i=\text{d}^Dx\varphi _i(x)\widehat{\mathrm{\Delta }}(x)`$. It is straightforward to compute that the corresponding scalar field $`\varphi _3(x)`$ is given by $`\varphi _3(x)`$ $`=`$ $`\text{Tr }\left(\widehat{\varphi }_1\widehat{\varphi }_2\widehat{\mathrm{\Delta }}(x)\right)`$ (2.14) $`=`$ $`{\displaystyle \frac{1}{\pi ^D\|det\theta \|}}{\displaystyle \text{d}^Dy\text{d}^Dz\varphi _1(y)\varphi _2(z)\text{e}^{2i(\theta ^1)_{\mu \nu }(xy)_\mu (xz)_\nu }}`$ $`=`$ $`\varphi _1(x)\mathrm{exp}\left({\displaystyle \frac{i}{2}}\stackrel{}{_\mu }\theta _{\mu \nu }\stackrel{}{_\nu }\right)\varphi _2(x)`$ $`\stackrel{\mathrm{def}}{=}`$ $`\varphi _1(x)\varphi _2(x),`$ where in the second line we have assumed that $`\theta _{\mu \nu }`$ is an invertible matrix. The multiplication in (2.14) is called the star or Moyal product of the fields $`\varphi _1(x)`$ and $`\varphi _2(x)`$. It is associative but not commutative. For $`\theta _{\mu \nu }=0`$, it reduces to the ordinary product of functions and the theory (2.13) to the usual scalar field theory (2.1). One can easily show that $`\text{d}^Dx\varphi _i(x)\varphi _j(x)\mathrm{}\varphi _k(x)`$ is invariant under cyclic permutation of the fields, justifying the representation of the trace in (2.10). Since $$\text{d}^Dx\varphi _1(x)\varphi _2(x)=\text{d}^Dx\varphi _1(x)\varphi _2(x),$$ (2.15) noncommutative field theory and ordinary field theory are identical at the level of free fields. The differences come when one considers the interaction term $`\text{Tr }\left(\widehat{\varphi }^4\right)`$ $`=`$ $`{\displaystyle \text{d}^Dx\varphi (x)\varphi (x)\varphi (x)\varphi (x)}`$ $`=`$ $`{\displaystyle \left(\underset{i=1}{\overset{4}{}}\frac{\text{d}^Dk_i}{(2\pi )^D}\right)(2\pi )^D\delta ^D\left(\underset{i=1}{\overset{4}{}}k_i\right)\left(\underset{i=1}{\overset{4}{}}\stackrel{~}{\varphi }(k_i)\right)\text{e}^{\frac{i}{2}\theta _{\mu \nu }_{i<j}k_{\mu i}k_{\nu j}}}.`$ We see that the vertex in the noncommutative field theory contains a momentum dependent phase factor. The interaction is therefore non-local. Naively one might expect that the effect of the noncommutativity becomes negligible at energy scales much smaller than $`\|\theta _{\mu \nu }\|^{1/2}`$. This is not quite the case. For example, it has been shown that the noncommutativity drastically alters the infrared dynamics of the theory . The perturbative renormalizability of noncommutative scalar field theory has been studied in Ref. . ### 2.2 Noncommutative Yang-Mills theory Let us now turn to gauge theory on a noncommutative spacetime. A hermitian operator corresponding to a U($`p`$) gauge field $`A_\mu (x)`$ can be introduced as $$\widehat{A}_\mu =\text{d}^Dx\widehat{\mathrm{\Delta }}(x)A_\mu (x),$$ (2.17) where $`\widehat{\mathrm{\Delta }}(x)`$ is given by (2.6). The action can be written as $`S`$ $`=`$ $`\text{Tr }\mathrm{tr}_{(p)}\left([\widehat{}_\mu ,\widehat{A}_\nu ][\widehat{}_\nu ,\widehat{A}_\mu ]+i[\widehat{A}_\mu ,\widehat{A}_\nu ]\right)^2`$ (2.18) $`=`$ $`{\displaystyle \text{d}^Dx\mathrm{tr}_{(p)}\left(F_{\mu \nu }(x)F_{\mu \nu }(x)\right)},`$ (2.19) where $$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu +i(A_\mu A_\nu A_\nu A_\mu ).$$ (2.20) As in (2.10), we use the symbol Tr to denote the operator trace over coordinates, while $`\mathrm{tr}_{(p)}`$ denotes the (finite-dimensional) trace in the fundamental representation of the U$`(p)`$ gauge group. The action (2.19) is invariant under the “star-gauge” transformation $$A_\mu (x)g(x)A_\mu (x)g(x)^{}ig(x)_\mu g(x)^{},$$ (2.21) where the gauge function $`g(x)`$ is a $`p\times p`$ matrix field satisfying $$g(x)g(x)^{}=11_p,$$ (2.22) namely it is star-unitary. Introducing a unitary operator corresponding to the gauge function $`g(x)`$ through $$\widehat{g}=\text{d}^Dx\widehat{\mathrm{\Delta }}(x)g(x),$$ (2.23) the gauge transformation (2.21) can be written in terms of Weyl operators as $$\widehat{A}_\mu \widehat{g}\widehat{A}_\mu \widehat{g}^{}i\widehat{g}[\widehat{}_\mu ,\widehat{g}^{}],$$ (2.24) under which the action (2.18) is invariant. As can be seen from the action (2.19), three-point and four-point gauge interactions exist even for the simplest case of a U(1) gauge group. It is known that noncommutative U(1) gauge theory in four dimensions is asymptotically free, and in fact that its beta-function coincides exactly with that of large $`N`$ Yang-Mills theory (See eq. (1.1)). This can be understood as a consequence of Morita equivalence of noncommutative gauge theories, which we will derive in Section 3. One form of Morita equivalence states that the non-abelian structure of the gauge group can be absorbed into the noncommutativity of spacetime. Note in this regard that Eq. (2.17) already has a suggestive form, as it shows that the spacetime indices and the gauge group indices are treated on the same footing in noncommutative geometry. ### 2.3 Star-gauge invariant observables In Ref. , star-gauge invariant observables in noncommutative Yang-Mills theory have been discovered using a twisted large $`N`$ reduced model. In this subsection, we construct star-gauge invariant observables in the present general formalism following Ref. . For this, we first note that a novel feature of the star-product is that spacetime translations can be represented on fields via star-conjugation by functions on spacetime. Let us consider a star-unitary function $`S_v(x)`$ with the property<sup>2</sup><sup>2</sup>2In principle, $`S_v(x)`$ can be a $`p\times p`$ matrix-valued function, but we shall find below that it is proportional to the unit matrix. $$S_v(x)g(x)S_v(x)^{}=g(x+v)$$ (2.25) for arbitrary functions $`g(x)`$, where $`v`$ is a real $`D`$-dimensional vector. It is straightforward to show that a necessary and sufficient condition for such a function to exist is $$\theta _{\mu \nu }_\nu S_v(x)=iv_\mu S_v(x).$$ (2.26) In particular, when $`\theta _{\mu \nu }`$ is invertible, the solution to (2.26) exists for arbitrary $`v_\mu `$ and is given by the plane wave $$S_v(x)=\text{e}^{ik_\mu x_\mu }\mathrm{\hspace{0.17em}1}1_p\mathrm{with}k_\mu =(\theta ^1)_{\mu \nu }v_\nu .$$ (2.27) The uniqueness of the function (2.27) is immediate. If $`S_v^{}(x)`$ is another function with the property (2.25), then the function $`S_v^{}(x)S_v(x)`$ star-commutes with all functions $`g(x)`$. In particular, by taking $`g(x)=\text{e}^{iw_\mu x_\mu }`$ for arbitrary $`w_\mu \text{}`$ we may conclude that the function $`S_v^{}(x)S_v(x)`$ is independent of $`x`$, and thus the two functions $`S_v(x)`$ and $`S_v^{}(x)`$ coincide up to some irrelevant phase factor. We now construct an analog of the parallel transport operator as $`𝒰(x;C)`$ $`=`$ $`\mathrm{P}\mathrm{exp}_{}\left(i{\displaystyle \underset{C}{}}\text{d}\xi ^\mu A_\mu (x+\xi )\right)`$ (2.28) $`\stackrel{\mathrm{def}}{=}`$ $`1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}i^n{\displaystyle \underset{0}{\overset{1}{}}}\text{d}\sigma _1{\displaystyle \underset{\sigma _1}{\overset{1}{}}}\text{d}\sigma _2\mathrm{}{\displaystyle \underset{\sigma _{n1}}{\overset{1}{}}}\text{d}\sigma _n\xi ^{{}_{}{}^{}\mu _{1}^{}}(\sigma _1)\mathrm{}\xi ^{{}_{}{}^{}\mu _{n}^{}}(\sigma _n)`$ $`\times A_{\mu _1}\left(x+\xi (\sigma _1)\right)\mathrm{}A_{\mu _n}\left(x+\xi (\sigma _n)\right),`$ where $`C`$ is an oriented curve in $`D`$-dimensional spacetime parametrized by the functions $`\xi _\mu (\sigma )`$ with $`0\sigma 1`$. We fix the starting point of the curve to be the origin, $`\xi _\mu (0)=0`$, and denote its endpoint by $`v_\mu =\xi _\mu (1)`$ in what follows. The operator $`𝒰(x;C)`$ can be regarded as a $`p\times p`$ star-unitary matrix field depending on $`C`$. Under the star-gauge transformation (2.21), it transforms as $$𝒰(x;C)g(x)𝒰(x;C)g(x+v)^{}.$$ (2.29) Star-gauge invariant observables can be constructed out of (2.28) by using the function (2.27), with the property (2.25), as $$𝒪(C)=\text{d}^Dx\mathrm{tr}_{(p)}\left(𝒰(x;C)S_v(x)\right).$$ (2.30) In Eq. (2.30), the parameter $`k_\mu `$ in (2.27) can be interpreted as the total momentum of the contour $`C`$. Note that in the commutative case $`\theta _{\mu \nu }=0`$, gauge invariance requires that the loop be closed irrespectively of the total momentum. In the noncommutative case, on the other hand, star-gauge invariance does not require that the loop be closed, but it does require the separation vector $`v_\mu `$ of the loop to be proportional to the total momentum, as is seen from (2.27). The larger the total momentum is, the longer becomes the open loop. We will study further the dynamics of noncommutative Wilson loops in Section 6.1 by introducing matter fields as probes. Let us now rewrite the star-gauge invariant quantity (2.30) in terms of Weyl operators. After a little algebra, the parallel transport operator can be rewritten as $$𝒰(x;C)=\text{Tr }\left(\widehat{U}(C)\widehat{D}(C)^{}\widehat{\mathrm{\Delta }}(x)\right),$$ (2.31) where we have introduced the unitary operators $`\widehat{U}(C)`$ $`=`$ $`\mathrm{P}\mathrm{exp}\left({\displaystyle \underset{C}{}}\text{d}\xi ^\mu \left(\widehat{}_\mu +i\widehat{A}_\mu \right)\right),`$ $`\widehat{D}(C)`$ $`=`$ $`\mathrm{P}\mathrm{exp}\left({\displaystyle \underset{C}{}}\text{d}\xi ^\mu \widehat{}_\mu \right).`$ (2.32) Under a gauge transformation (2.24), $`\widehat{U}(C)`$ transforms as $$\widehat{U}(C)\widehat{g}\widehat{U}(C)\widehat{g}^{}.$$ (2.33) Note that $`\widehat{D}(C)`$ is nothing but a translation operator (up to an irrelevant phase factor), $$\widehat{D}(C)\widehat{\mathrm{\Delta }}(x)\widehat{D}(C)^{}=\widehat{\mathrm{\Delta }}(x+v).$$ (2.34) We also introduce the unitary operator $$\widehat{S}_v=\text{d}^Dx\widehat{\mathrm{\Delta }}(x)S_v(x)=\text{e}^{ik_\mu \widehat{x}_\mu }11_p$$ (2.35) which satisfies $$\widehat{S}_v\widehat{\mathrm{\Delta }}(x)\widehat{S}_v^{}=\widehat{\mathrm{\Delta }}(x+v)11_p$$ (2.36) corresponding to (2.25). The Weyl operator description of star-gauge invariant observables is then given by $$𝒪(C)=\text{Tr }\mathrm{tr}_{(p)}\left(\widehat{U}(C)\widehat{D}(C)^{}\widehat{S}_v\right).$$ (2.37) Its gauge invariance can be checked directly by noting that the operator $`\widehat{D}(C)^{}\widehat{S}_v`$ in (2.37) commutes with $`\widehat{\mathrm{\Delta }}(x)`$ due to (2.34) and (2.36). ### 2.4 The noncommutative torus We now briefly discuss the case where the spacetime is a $`D`$-dimensional torus $`\text{𝕋}^D`$ instead of $`\text{}^D`$. Let us first consider the case discussed in subsection 2.1 and impose periodic boundary conditions on the scalar field $`\varphi (x)`$ (Twisted boundary conditions will be studied in the next Section), $$\varphi (x+\mathrm{\Sigma }_{\mu a}\widehat{\mu })=\varphi (x),a=1,\mathrm{},D,$$ (2.38) where $`\widehat{\mu }`$ is a unit vector in the $`\mu `$-th direction of spacetime, and $`\mathrm{\Sigma }_{\mu a}`$ is the $`D\times D`$ period matrix of the torus which is a vielbein for the metric of $`\text{𝕋}^D`$. Here and in the following we use Greek letters for spacetime indices and Latin letters for frame indices. Due to the periodicity (2.38), the momenta $`k_\mu `$ in the Fourier mode expansion (2.2) are quantized as $$k_\mu =2\pi (\mathrm{\Sigma }^1)_{a\mu }m_a,m_a\text{}.$$ (2.39) We define a mapping of fields into operators as in (2.5) but now with the $`\widehat{\mathrm{\Delta }}(x)`$ given by $$\widehat{\mathrm{\Delta }}(x)=\frac{1}{\|det\mathrm{\Sigma }\|}\underset{\stackrel{}{m}\text{}^D}{}\left(\underset{a=1}{\overset{D}{}}\left(\widehat{Z}_a\right)^{m_a}\right)\text{e}^{\pi i_{a<b}\mathrm{\Theta }_{ab}m_am_b}\text{e}^{2\pi i(\mathrm{\Sigma }^1)_{a\mu }m_ax_\mu },$$ (2.40) where the operators $$\widehat{Z}_a=\text{e}^{2\pi i(\mathrm{\Sigma }^1)_{a\mu }\widehat{x}_\mu }$$ (2.41) satisfy the commutation relations $`\widehat{Z}_a\widehat{Z}_b`$ $`=`$ $`\text{e}^{2\pi i\mathrm{\Theta }_{ab}}\widehat{Z}_b\widehat{Z}_a`$ $`[\widehat{}_\mu ,\widehat{Z}_a]`$ $`=`$ $`2\pi i\left(\mathrm{\Sigma }^1\right)_{a\mu }\widehat{Z}_a,`$ (2.42) with $$\mathrm{\Theta }_{ab}=2\pi \left(\mathrm{\Sigma }^1\right)_{a\mu }\theta _{\mu \nu }\left(\mathrm{\Sigma }^1\right)_{b\nu }$$ (2.43) the dimensionless noncommutativity parameter. The basis (2.40) has the requisite properties $`\widehat{\mathrm{\Delta }}(x+\mathrm{\Sigma }_{\mu a}\widehat{\mu })`$ $`=`$ $`\widehat{\mathrm{\Delta }}(x),a=1,\mathrm{},D,`$ (2.44) $`[\widehat{}_\mu ,\widehat{\mathrm{\Delta }}(x)]`$ $`=`$ $`_\mu \widehat{\mathrm{\Delta }}(x).`$ (2.45) Note that the torus $`\text{𝕋}^D`$ has a discrete geometrical symmetry given by its SL$`(D,\text{})`$ automorphism group, under which the period matrix transforms as $`\mathrm{\Sigma }\mathrm{\Sigma }\mathrm{\Lambda }`$ with $`\mathrm{\Lambda }\mathrm{SL}(D,\text{})`$. The dimensionless noncommutativity parameter (2.43) transforms as $`\mathrm{\Theta }\mathrm{\Lambda }^1\mathrm{\Theta }(\mathrm{\Lambda }^1)^{}`$ and the coordinate operators (2.41) as $`\widehat{Z}_a_{b=1}^D(\widehat{Z}_b)^{(\mathrm{\Lambda }^1)_{ab}}`$. All formulae such as (2.40)–(2.43) are manifestly SL$`(D,\text{})`$ covariant. This symmetry, which persists in the case of generic twisted boundary conditions on the fields, will be exploited in the next Section. Let us now consider star-gauge invariant observables in noncommutative Yang-Mills theory on the torus $`\text{𝕋}^D`$ . We can define the parallel transport operator $`𝒰(x;C)`$ as in (2.28). The star-gauge transformation of $`𝒰(x;C)`$ is given by (2.29), but now the star-unitary function $`g(x)`$ is a single-valued function on the torus. We can define a star-gauge invariant observable as in (2.30), where $`S_v(x)`$ is a function on the $`D`$-dimensional torus which satisfies (2.25) for arbitrary $`g(x)`$ that we are considering. One finds that a necessary and sufficient condition is $$S_v\left(x+2\pi (\mathrm{\Sigma }^1\theta )_{a\mu }\widehat{\mu }\right)=\text{e}^{2\pi i(\mathrm{\Sigma }^1)_{a\mu }v_\mu }S_v(x).$$ (2.46) The solution can again be given by the plane wave $`S_v(x)=\text{e}^{ik_\mu x_\mu }\mathrm{\hspace{0.17em}1}1_p`$, but in the present case $`k_\mu `$ is quantized as in (2.39). The condition that must be met is $$v_\mu =\theta _{\mu \nu }k_\nu +\mathrm{\Sigma }_{\mu a}n_a$$ (2.47) for some integer-valued vector $`n_a`$. Thus we obtain an analog of the Polyakov line in noncommutative Yang-Mills theory represented by the existence of the case with non-zero “winding number” $`n_a`$. ## 3 Morita equivalence In noncommutative geometry, there is a remarkable geometric equivalence relation on certain classes of noncommutative spaces known as Morita equivalence . Roughly speaking, two spaces are Morita equivalent if one space can be regarded as a twisted operator bundle over the other space of a certain topological charge whose fibers are operator algebras. Many topological quantities are preserved by the Morita equivalence relation. In particular, K-theory groups are invariant under it, so that two Morita equivalent spaces should have a canonical mapping between gauge bundles defined over them. In noncommutative Yang-Mills theory, this implies a remarkable duality between gauge theories over different noncommutative tori , for example, which relates a Yang-Mills theory with background magnetic flux to a gauge theory with gauge group of lower rank and no background flux. It allows one to interpolate continuously, through noncommutative Yang-Mills theories, between two ordinary Yang-Mills theories with gauge groups of different rank and appropriate background magnetic fluxes. Furthermore, in certain instances, there is the remarkable fact that the non-abelian nature of a gauge group can be absorbed into the noncommutativity of spacetime by mapping a U$`(p)`$ gauge theory with multi-valued gauge fields to a U(1) gauge theory with single-valued fields on a dual noncommutative torus. In string theory, Morita equivalence coincides with the phenomenon of T-duality of open string backgrounds in the presence of D-branes . The nonperturbative formulation of generic noncommutative Yang-Mills theories that we shall present in Section 4 relies crucially on the Morita equivalence property. We shall therefore present a detailed and rather general derivation of the Morita equivalence relation for noncommutative tori from the point of view adopted in the previous Section to noncommutative Yang-Mills theory. We shall see that, in the present formalism, the equivalence has a natural geometrical interpretation as a change of basis for the map between fields and Weyl operators for the noncommutative geometry. We will also show that this duality holds at the level of star-gauge invariant observables in the two gauge theories, thereby providing a complete equivalence of the noncommutative quantum field theories. ### 3.1 Twisted gauge theory on the noncommutative torus Consider a U($`p`$) gauge theory on a noncommutative torus $`\text{𝕋}_\mathrm{\Theta }^D`$ with gauge fields of non-vanishing topological charge. Such gauge fields are not single-valued functions on the torus. Instead, they must be regarded as functions on the universal covering space $`\text{}^D`$ with the twisted boundary conditions $$A_\mu (x+\mathrm{\Sigma }_{\nu a}\widehat{\nu })=\mathrm{\Omega }_a(x)A_\mu (x)\mathrm{\Omega }_a(x)^{}i\mathrm{\Omega }_a(x)_\mu \mathrm{\Omega }_a(x)^{}.$$ (3.1) The transition functions $`\mathrm{\Omega }_a(x)`$ are star-unitary $`p\times p`$ matrices, which we decompose into an abelian part and an SU$`(p)`$ part via the gauge choice $$\mathrm{\Omega }_a(x)=\text{e}^{i\alpha _{a\mu }x_\mu }\mathrm{\Gamma }_a,$$ (3.2) where $`\alpha `$ is a real-valued constant $`D\times D`$ matrix satisfying $`(\alpha \mathrm{\Sigma })^{}=\alpha \mathrm{\Sigma }`$, and $`\mathrm{\Gamma }_a`$ are SU$`(p)`$ matrices. The matrix $`\alpha `$ will account for the abelian fluxes of the gauge fields. The antisymmetry of the matrix $`\alpha \mathrm{\Sigma }`$ implies that the transition function $`\mathrm{\Omega }_a(x)`$ has periodicity $`\mathrm{\Omega }_a(x+\mathrm{\Sigma }_{\mu a}\widehat{\mu })=\mathrm{\Omega }_a(x)`$. Consistency of the conditions (3.1) requires the $`\mathrm{\Omega }_a`$ to obey the cocycle condition $$\mathrm{\Omega }_a(x+\mathrm{\Sigma }_{\mu b}\widehat{\mu })\mathrm{\Omega }_b(x)=𝒵_{ab}\mathrm{\Omega }_b(x+\mathrm{\Sigma }_{\mu a}\widehat{\mu })\mathrm{\Omega }_a(x),$$ (3.3) where $`𝒵_{ab}=\text{e}^{2\pi i\gamma _{ab}/p}`$ are elements of the center of the SU($`p`$) part of the gauge group, with $`\gamma `$ an antisymmetric integral $`D\times D`$ matrix. In commutative SU($`p`$) gauge theory, a phase factor such as $`𝒵_{ab}`$ would induce a non-abelian ’t Hooft flux for the gauge fields . In that case, the matrix $`\alpha `$ in (3.2) should be set to zero, and therefore having a non-trivial $`𝒵_{ab}`$ is the only way to twist the boundary conditions on the gauge fields. In commutative U($`p`$) gauge theory, keeping both $`\gamma `$ and $`\alpha `$ non-zero is redundant because one can eliminate the U(1) twists $`\alpha `$ by $`\text{}_p`$-valued phase factors according to the global decomposition $`\mathrm{U}(p)=\mathrm{U}(1)\times \mathrm{SU}(p)/\text{}_p`$, and one can set either $`\gamma =0`$ or $`\alpha =0`$ without loss of generality. However, this is not quite the case in noncommutative U($`p`$) gauge theory, as we will see. Eq. (3.3) implies that the matrices $`\mathrm{\Gamma }_a`$ mutually commute up to some phases, $$\mathrm{\Gamma }_a\mathrm{\Gamma }_b=\text{e}^{2\pi iQ_{ab}/p}\mathrm{\Gamma }_b\mathrm{\Gamma }_a.$$ (3.4) Taking the determinant of both sides of (3.4) shows that the antisymmetric matrix $`Q`$ has elements $`Q_{ab}\text{}`$. Furthermore, eq. (3.3) gives the consistency condition $$Q=\frac{p}{2\pi }\left(2\alpha \mathrm{\Sigma }\alpha \theta \alpha ^{}\right)\gamma ,$$ (3.5) where here and in the following we use a matrix multiplication convention. The integer $`Q_{ab}`$ is the magnetic flux of the gauge field through the surface formed by the $`a`$-th and $`b`$-th cycles of the torus. The central extension of the cocycle relation (3.3) shifts $`Q`$ by the matrix $`\gamma `$ of non-abelian ’t Hooft fluxes representing the associated principal curvatures. The fact that the trivial configuration $`A_\mu (x)=0`$ is not a solution of (3.1) motivates the introduction of a fixed, multi-valued background abelian gauge field $`A_\mu ^{(0)}(x)`$ defined by $$A_\mu ^{(0)}(x)=\frac{1}{2}F_{\mu \nu }x_\nu 11_p,$$ (3.6) where $`F`$ is a real-valued constant antisymmetric $`D\times D`$ matrix. The twisted boundary conditions (3.1) for the gauge field $`A_\mu ^{(0)}(x)`$ are equivalent to the relations $`\alpha `$ $`=`$ $`\mathrm{\Sigma }^{}F{\displaystyle \frac{1}{211_D+\theta F}}`$ $`F`$ $`=`$ $`2\alpha ^{}{\displaystyle \frac{1}{\mathrm{\Sigma }\theta \alpha ^{}}}.`$ (3.7) We then decompose $`A_\mu (x)`$ into a part representing a particular solution to the twisted boundary conditions (3.1) and a part representing the fluctuations around the fixed background as $$A_\mu (x)=A_\mu ^{(0)}(x)+𝒜_\mu (x),$$ (3.8) where $`𝒜_\mu (x)`$ satisfies the constraints $$𝒜_\mu (x+\mathrm{\Sigma }_{\nu a}\widehat{\nu })=\mathrm{\Omega }_a(x)𝒜_\mu (x)\mathrm{\Omega }_a(x)^{}.$$ (3.9) This means that the field $`𝒜_\mu (x)`$ is an adjoint section of the corresponding gauge bundle over the noncommutative torus. Let us now consider the noncommutative Yang-Mills action $$S=\frac{1}{g^2}\text{d}^Dx\mathrm{tr}_{(p)}\left(F_{\mu \nu }(x)f_{\mu \nu }\mathrm{\hspace{0.17em}1}1_p\right)_{}^2,$$ (3.10) where the noncommutative field strength tensor $`F_{\mu \nu }(x)`$ is defined in (2.20). The constant background tensor field $`f_{\mu \nu }`$ will be specified later, and the integration is taken over the torus $`\text{𝕋}^D`$. Using the decomposition (3.8), the action (3.10) can be written as $$S=\frac{1}{g^2}\text{d}^Dx\mathrm{tr}_{(p)}\left(_{\mu \nu }(x)+F_{\mu \nu }^{(0)}f_{\mu \nu }\mathrm{\hspace{0.17em}1}1_p\right)_{}^2,$$ (3.11) where $`_{\mu \nu }`$ $`=`$ $`_\mu 𝒜_\nu _\nu 𝒜_\mu +i\left(𝒜_\mu 𝒜_\nu 𝒜_\nu 𝒜_\mu \right)`$ (3.12) $`+i\left(A_\mu ^{(0)}𝒜_\nu 𝒜_\nu A_\mu ^{(0)}\right)i\left(A_\nu ^{(0)}𝒜_\mu 𝒜_\mu A_\nu ^{(0)}\right)`$ $`F_{\mu \nu }^{(0)}`$ $`=`$ $`_\mu A_\nu ^{(0)}_\nu A_\mu ^{(0)}+i\left(A_\mu ^{(0)}A_\nu ^{(0)}A_\nu ^{(0)}A_\mu ^{(0)}\right)=\left(F+{\displaystyle \frac{1}{4}}F\theta F\right)_{\mu \nu }11_p.`$ From (3.7) and the identity $$\left(\frac{1}{11_D\theta \alpha ^{}\mathrm{\Sigma }^1}\right)^211_p=\left(11_D+\frac{1}{2}\theta F\right)^211_p=11_D11_p+\theta F^{(0)}$$ (3.14) it follows that the relation (3.5) is equivalent to $$\mathrm{\Sigma }^{}F^{(0)}\mathrm{\Sigma }=2\pi \frac{1}{p11_D+(Q+\gamma )\mathrm{\Theta }}\left(Q+\gamma \right)11_p.$$ (3.15) Eq. (3.15) gives the relationship between the central curvatures, the topological charges, and the ’t Hooft fluxes of the gauge field configurations. Requiring that $`𝒜_\mu (x)=0`$ be the vacuum field configuration of the theory up to a gauge transformation fixes $`f_{\mu \nu }\mathrm{\hspace{0.17em}1}1_p=F_{\mu \nu }^{(0)}`$, so that the action becomes $$S=\frac{1}{g^2}\text{d}^Dx\mathrm{tr}_{(p)}\left(_{\mu \nu }(x)_{\mu \nu }(x)\right).$$ (3.16) ### 3.2 Transformation of the action In this subsection, we will derive the Morita equivalence relation between noncommutative Yang-Mills theories in arbitrary even dimension $`D=2d`$. Specifically, we will show that the theory (3.10) with the constraint (3.1) is equivalent to another noncommutative gauge theory on a torus with single-valued gauge fields. For this, we first rewrite the gauge theory of the previous subsection in terms of Weyl operators. We introduce the hermitian operator $`\widehat{A}_\mu `$ as in (2.17), and define the constant abelian background $`\widehat{A}_\mu ^{(0)}`$ and the transition function $`\widehat{\mathrm{\Omega }}_a`$ similarly. The explicit gauge choice (3.2) corresponds to the unitary operator $$\widehat{\mathrm{\Omega }}_a=\text{e}^{i\alpha _{a\mu }\widehat{x}_\mu }\mathrm{\Gamma }_a,$$ (3.17) and, corresponding to (3.8), we decompose $`\widehat{A}_\mu `$ as $$\widehat{A}_\mu =\widehat{A}_\mu ^{(0)}+\widehat{𝒜}_\mu .$$ (3.18) The action (3.10) or (3.16) can be written as $`S`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}\text{Tr }\mathrm{tr}_{(p)}\left([\widehat{}_\mu ,\widehat{A}_\nu ][\widehat{}_\nu ,\widehat{A}_\mu ]+i[\widehat{A}_\mu ,\widehat{A}_\nu ]f_{\mu \nu }\mathrm{\hspace{0.17em}1}1_p\right)^2`$ (3.19) $`=`$ $`{\displaystyle \frac{1}{g^2}}\text{Tr }\mathrm{tr}_{(p)}\left([\widehat{}_\mu ^{(0)},\widehat{𝒜}_\nu ][\widehat{}_\nu ^{(0)},\widehat{𝒜}_\mu ]+i[\widehat{𝒜}_\mu ,\widehat{𝒜}_\nu ]\right)^2,`$ (3.20) where $`f_{\mu \nu }`$ may be given in terms of $`\widehat{A}_\mu ^{(0)}`$ as $$f_{\mu \nu }\mathrm{\hspace{0.17em}1}1_p=[\widehat{}_\mu ,\widehat{A}_\nu ^{(0)}][\widehat{}_\nu ,\widehat{A}_\mu ^{(0)}]+i[\widehat{A}_\mu ^{(0)},\widehat{A}_\nu ^{(0)}]$$ (3.21) and we have introduced the fiducial constant curvature connection $$\widehat{}_\mu ^{(0)}=\widehat{}_\mu +i\widehat{A}_\mu ^{(0)}.$$ (3.22) The constraint (3.1) and the equivalent one (3.9) can be written in terms of operators as $`\text{e}^{\mathrm{\Sigma }_{\nu a}\widehat{}_\nu }\widehat{A}_\mu \text{e}^{\mathrm{\Sigma }_{\nu a}\widehat{}_\nu }`$ $`=`$ $`\widehat{\mathrm{\Omega }}_a\widehat{A}_\mu \widehat{\mathrm{\Omega }}_a^{}i\widehat{\mathrm{\Omega }}_a[\widehat{}_\mu ,\widehat{\mathrm{\Omega }}_a^{}]`$ (3.23) $`\text{e}^{\mathrm{\Sigma }_{\nu a}\widehat{}_\nu }\widehat{𝒜}_\mu \text{e}^{\mathrm{\Sigma }_{\nu a}\widehat{}_\nu }`$ $`=`$ $`\widehat{\mathrm{\Omega }}_a\widehat{𝒜}_\mu \widehat{\mathrm{\Omega }}_a^{}.`$ (3.24) Our next task is to solve this constraint for the gauge field configurations. First of all, we see from (3.4) that the $`\mathrm{\Gamma }_a`$ are twist eating solutions for SU$`(p)`$. For generic rank $`p`$ and flux matrix $`Q`$, these matrices may be constructed as follows . For this, it is convenient to use the discrete SL$`(D,\text{})`$ symmetry of $`\text{𝕋}^D`$ to transform the twist eaters into matrices $`\mathrm{\Gamma }_b^{}`$ through $$\mathrm{\Gamma }_a=\underset{b=1}{\overset{D}{}}(\mathrm{\Gamma }_b^{})^{\mathrm{\Lambda }_{ba}},$$ (3.25) where $`\mathrm{\Lambda }\mathrm{SL}(D,\text{})`$. The transformed twist eaters $`\mathrm{\Gamma }_a^{}`$ satisfy the commutation relations $$\mathrm{\Gamma }_a^{}\mathrm{\Gamma }_b^{}=\text{e}^{2\pi iQ_{ab}^{}/p}\mathrm{\Gamma }_b^{}\mathrm{\Gamma }_a^{},$$ (3.26) with $`Q=\mathrm{\Lambda }^{}Q^{}\mathrm{\Lambda }`$. We can now choose $`\mathrm{\Lambda }`$ so that the matrix $`Q^{}`$ takes a canonical skew-diagonal form $$Q^{}=\left(\begin{array}{ccccc}0& q_1& & & \\ q_1& 0& & & \\ & & \mathrm{}& & \\ & & & 0& q_d\\ & & & q_d& 0\end{array}\right).$$ (3.27) Given the $`d`$ independent fluxes $`q_i\text{}`$, we introduce the three integers $$p_i=\mathrm{gcd}(q_i,p),\stackrel{~}{p}_i=\frac{p}{p_i},\stackrel{~}{q}_i=\frac{q_i}{p_i}.$$ (3.28) By construction, $`\stackrel{~}{p}_i`$ and $`\stackrel{~}{q}_i`$ are co-prime. A necessary and sufficient condition for the existence of solutions to (3.26) is that the integer $`\stackrel{~}{p}_1\mathrm{}\stackrel{~}{p}_d`$, which is the dimension of the irreducible representation of the Weyl-’t Hooft algebra (3.4), divides the rank $`p`$. This is a condition which must be met by the geometrical parameters of the given constant curvature bundle. In that case we write $$p=\stackrel{~}{p}_0\underset{i=1}{\overset{d}{}}\stackrel{~}{p}_i$$ (3.29) and the twist eating solutions may then be given on the subgroup $`\mathrm{SU}(\stackrel{~}{p}_1)\mathrm{}\mathrm{SU}(\stackrel{~}{p}_d)\mathrm{SU}(\stackrel{~}{p}_0)`$ of SU$`(p)`$ as $`\mathrm{\Gamma }_{2i1}^{}`$ $`=`$ $`11_{\stackrel{~}{p}_1}\mathrm{}V_{\stackrel{~}{p}_i}\mathrm{}11_{\stackrel{~}{p}_d}11_{\stackrel{~}{p}_0}`$ $`\mathrm{\Gamma }_{2i}^{}`$ $`=`$ $`11_{\stackrel{~}{p}_1}\mathrm{}\left(W_{\stackrel{~}{p}_i}\right)^{\stackrel{~}{q}_i}\mathrm{}11_{\stackrel{~}{p}_d}11_{\stackrel{~}{p}_0}`$ (3.30) for $`i=1,\mathrm{},d`$. Here $`V_p`$ and $`W_p`$ are the SU$`(p)`$ shift and clock matrices $$V_p=\left(\begin{array}{ccccc}0& 1& & & 0\\ & 0& 1& & \\ & & \mathrm{}& \mathrm{}& \\ & & & \mathrm{}& 1\\ 1& & & & 0\end{array}\right),W_p=\left(\begin{array}{ccccc}1& & & & \\ & \text{e}^{2\pi i/p}& & & \\ & & \text{e}^{4\pi i/p}& & \\ & & & \mathrm{}& \\ & & & & \text{e}^{2\pi i(p1)/p}\end{array}\right)$$ (3.31) obeying $`V_pW_p=\text{e}^{2\pi i/p}W_pV_p`$. We can then obtain twist eaters $`\mathrm{\Gamma }_a`$ satisfying (3.4) by using the relation (3.25). Note that the space of matrices which commute with the $`\mathrm{\Gamma }_a`$ is the gl$`(\stackrel{~}{p}_0,\text{})`$ subspace of $`\mathrm{gl}(p,\text{})`$ generated by the matrices $`11_{\stackrel{~}{p}_1}\mathrm{}11_{\stackrel{~}{p}_d}Z_0`$, $`Z_0\mathrm{gl}(\stackrel{~}{p}_0,\text{})`$. Since the set $`\{(V_p)^j(W_p)^j^{}\|j,j^{}\text{}_p\}`$ spans the linear space gl($`p,\text{}`$), the operator $`\widehat{𝒜}_\mu `$ may be expanded as $$\widehat{𝒜}_\mu =\underset{\stackrel{}{j}}{}\text{d}^Dk\text{e}^{ik_\nu \widehat{x}_\nu }\underset{a=1}{\overset{D}{}}(\mathrm{\Gamma }_a)^{j_a}a_\mu (k,\stackrel{}{j}),$$ (3.32) where the $`\stackrel{~}{p}_0\times \stackrel{~}{p}_0`$ matrix-valued coefficients $`a_\mu (k,\stackrel{}{j})`$ are periodic functions of $`\stackrel{}{j}`$ with the periodicity $`j_aj_a+\stackrel{~}{P}_{ab}`$, $`b=1,\mathrm{},D`$. The matrix $`\stackrel{~}{P}`$ will be specified below. The constraint (3.24) then implies that $`a_\mu (k,\stackrel{}{j})`$ vanishes unless $$\zeta _a+\frac{1}{p}Q_{ab}j_b=n_a\text{},$$ (3.33) where we have defined the $`D`$-dimensional real-valued vector $$\zeta _a=\frac{1}{2\pi }k_\nu \left(\mathrm{\Sigma }+\theta \alpha ^{}\right)_{\nu a}.$$ (3.34) We will now find the general solution ($`\zeta _a`$, $`j_b`$) to eq. (3.33). Introducing the $`D\times D`$ integral matrices $$\stackrel{~}{P}^{}=\left(\begin{array}{ccccc}\stackrel{~}{p}_1& 0& & & \\ 0& \stackrel{~}{p}_1& & & \\ & & \mathrm{}& & \\ & & & \stackrel{~}{p}_d& 0\\ & & & 0& \stackrel{~}{p}_d\end{array}\right),\stackrel{~}{Q}^{}=\left(\begin{array}{ccccc}0& \stackrel{~}{q}_1& & & \\ \stackrel{~}{q}_1& 0& & & \\ & & \mathrm{}& & \\ & & & 0& \stackrel{~}{q}_d\\ & & & \stackrel{~}{q}_d& 0\end{array}\right),$$ (3.35) and $`\stackrel{~}{P}=\mathrm{\Lambda }^1\stackrel{~}{P}^{}\mathrm{\Lambda }^{}`$, $`\stackrel{~}{Q}=\mathrm{\Lambda }^{}\stackrel{~}{Q}^{}\mathrm{\Lambda }^{}`$, we can write $$Q=p\stackrel{~}{Q}\stackrel{~}{P}^1,$$ (3.36) where $`\mathrm{\Lambda }^{}\mathrm{SL}(D,\text{})`$ will represent the automorphism symmetry group of the resulting theory that we shall find. Then eq. (3.33) can be written as $`\zeta _a=m_c(\stackrel{~}{P}^1)_{ca}`$, where $`m_c`$ is a $`D`$-dimensional integral vector which satisfies $$m_c=n_b\stackrel{~}{P}_{bc}+j_b\stackrel{~}{Q}_{bc}.$$ (3.37) We will next show that for any given integer $`m_c`$, there exists a set of integers ($`n_b`$, $`j_b`$) which satisfies (3.37). For this, we note that since $`\stackrel{~}{p}_i`$ and $`\stackrel{~}{q}_i`$ are relatively prime, there exists a set of integers $`(a_i,b_i)`$ such that $$a_i\stackrel{~}{p}_i+b_i\stackrel{~}{q}_i=1,i=1,\mathrm{},d.$$ (3.38) Introducing the $`D\times D`$ integral matrices $$A^{}=\left(\begin{array}{ccccc}a_1& 0& & & \\ 0& a_1& & & \\ & & \mathrm{}& & \\ & & & a_d& 0\\ & & & 0& a_d\end{array}\right),B^{}=\left(\begin{array}{ccccc}0& b_1& & & \\ b_1& 0& & & \\ & & \mathrm{}& & \\ & & & 0& b_d\\ & & & b_d& 0\end{array}\right),$$ (3.39) and $`A=(\mathrm{\Lambda }^{})^1A^{}\mathrm{\Lambda }`$, $`B=(\mathrm{\Lambda }^{})^1B^{}(\mathrm{\Lambda }^1)^{}`$, we have $$A\stackrel{~}{P}+B\stackrel{~}{Q}=11_D.$$ (3.40) The integers $`n_b=m_aA_{ab}`$ and $`j_b=m_aB_{ab}`$ then give a solution to (3.37). It is easy to see that for any given $`m_a`$, the $`j_b`$ satisfying (3.37) are unique up to their periodicity. Thus, we find that the general solution to eq. (3.33) can be given by $$\zeta _a=m_b\left(\stackrel{~}{P}^1\right)_{ba},j_a=m_bB_{ba}m_b\text{}.$$ (3.41) Solving (3.34) for the momenta we obtain $`k_\mu =2\pi m_a\beta _{a\mu }`$, where we have defined the $`D\times D`$ matrix $$\beta =\frac{1}{(\mathrm{\Sigma }+\theta \alpha ^{})\stackrel{~}{P}}.$$ (3.42) Therefore, the most general gauge field configuration satisfying the constraint (3.24) is given as $$\widehat{𝒜}_\mu =\underset{\stackrel{}{m}\text{}^D}{}\text{e}^{\pi i_{a<b}\mathrm{\Theta }_{ab}^{}m_am_b}\left(\underset{a=1}{\overset{D}{}}\left(\widehat{Z}_a^{}\right)^{m_a}\right)\stackrel{~}{a}_\mu (\stackrel{}{m}),$$ (3.43) where $`\stackrel{~}{a}_\mu (\stackrel{}{m})`$ are $`\stackrel{~}{p}_0\times \stackrel{~}{p}_0`$ matrix-valued coefficients and $$\widehat{Z}_a^{}=\text{e}^{2\pi i\beta _{a\mu }\widehat{x}_\mu }\underset{b=1}{\overset{D}{}}(\mathrm{\Gamma }_b)^{B_{ab}}.$$ (3.44) Hermiticity of $`\widehat{𝒜}_\mu `$ requires $`\stackrel{~}{a}_\mu (\stackrel{}{m})=\stackrel{~}{a}_\mu (\stackrel{}{m})^{}`$. The operators (3.44) obey the commutation relations $`\widehat{Z}_a^{}\widehat{Z}_b^{}`$ $`=`$ $`\text{e}^{2\pi i\mathrm{\Theta }_{ab}^{}}\widehat{Z}_b^{}\widehat{Z}_a^{},`$ $`[\widehat{}_\mu ^{(0)},\widehat{Z}_a^{}]`$ $`=`$ $`2\pi i\left(\mathrm{\Sigma }^{}_{}{}^{}1\right)_{a\mu }\widehat{Z}_a^{},`$ (3.45) where $`\mathrm{\Theta }^{}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Theta }\left(\stackrel{~}{Q}+\frac{1}{p}\gamma \stackrel{~}{P}\right)\stackrel{~}{P}}}\left[\left(A{\displaystyle \frac{1}{p}}B\gamma \right)\mathrm{\Theta }+B\right]^{},`$ (3.46) $`\mathrm{\Sigma }^{^{}}`$ $`=`$ $`\mathrm{\Sigma }\left[\mathrm{\Theta }\left(\stackrel{~}{Q}+{\displaystyle \frac{1}{p}}\gamma \stackrel{~}{P}\right)\stackrel{~}{P}\right].`$ (3.47) The commutation relations (3.45) are the same as (2.42) with the replacements $`\widehat{Z}\widehat{Z}^{}`$, $`\widehat{}\widehat{}^{(0)}`$, $`\mathrm{\Theta }\mathrm{\Theta }^{}`$ and $`\mathrm{\Sigma }\mathrm{\Sigma }^{}`$. We may therefore define a basis $`\widehat{\mathrm{\Delta }}^{}(x^{})`$ for the mapping of fields into operators as in (2.40). The $`x_\mu ^{}\text{}`$ are interpreted as coordinates on a new, dual torus with period matrix $`\mathrm{\Sigma }^{}`$ given by (3.47). Furthermore, the dimensionful noncommutativity parameters $`\theta _{\mu \nu }^{}`$ which define the star-product are given by the dimensionless parameters (3.46) as $$\theta ^{}=\frac{1}{2\pi }\mathrm{\Sigma }^{}\mathrm{\Theta }^{}\mathrm{\Sigma }^{}.$$ (3.48) The expansion (3.43) then becomes $$\widehat{𝒜}_\mu =\text{d}^Dx^{}\widehat{\mathrm{\Delta }}^{}(x^{})𝒜_\mu ^{}(x^{}),$$ (3.49) where $`𝒜_\mu ^{}(x^{})`$ can be regarded as a single-valued U($`\stackrel{~}{p}_0`$) gauge field on the noncommutative torus $`\text{𝕋}_\mathrm{\Theta }^{}^D`$ of total volume $`\|det\mathrm{\Sigma }^{}\|`$. The operator trace $`\text{Tr }^{}`$, with $`\text{Tr }^{}\widehat{\mathrm{\Delta }}^{}(x^{})=1`$, over the coordinates of this new torus is related to the original trace Tr through $$\text{Tr }^{}\mathrm{tr}_{(\stackrel{~}{p}_0)}=\frac{\stackrel{~}{p}_0}{p}\left\|\frac{det\mathrm{\Sigma }^{}}{det\mathrm{\Sigma }}\right\|\text{Tr }\mathrm{tr}_{(p)}.$$ (3.50) Using (3.20) and (3.50) we arrive at the canonical form of the noncommutative Yang-Mills action for the gauge field in (3.49), $$S=\frac{1}{g^2}\text{d}^Dx^{}\mathrm{tr}_{(\stackrel{~}{p}_0)}(_{\mu \nu }^{}(x^{})^{}_{\mu \nu }(x^{})),$$ (3.51) where $$_{\mu \nu }^{}=_\mu ^{}𝒜_\nu ^{}_\nu ^{}𝒜_\mu ^{}+i\left(𝒜_\mu ^{}^{}𝒜_\nu ^{}𝒜_\nu ^{}^{}𝒜_\mu ^{}\right),$$ (3.52) $`^{}`$ denotes the star product defined using $`\theta ^{}`$ instead of $`\theta `$, and the new Yang-Mills coupling constant is given by<sup>3</sup><sup>3</sup>3The transformation (3.53) of the Yang-Mills coupling constant differs from the standard one which is derived using string theoretical T-duality. The source of this discrepency owes to the normalization of the trace Tr on the noncommutative torus $`\text{𝕋}_\mathrm{\Theta }^D`$. In string theory, this trace is normalized according to the standard formula in noncommutative geometry for the rank of the module of sections of the fundamental bundle associated to the given Chan-Paton gauge bundle over $`\text{𝕋}_\mathrm{\Theta }^D`$ (This module coincides with the Hilbert space of open string ground states) . In the present field theoretical case, which does not require the specification of any representation of the Weyl operators, we have chosen the more natural volume normalization $`\text{Tr }\mathrm{\hspace{0.17em}1}1=\|det\mathrm{\Sigma }\|`$. $$g^{}_{}{}^{}2=g^2\frac{\left\|det\left[\mathrm{\Theta }(Q+\gamma )p11_D\right]\right\|}{\stackrel{~}{p}_0p^{D1}},$$ (3.53) where we have used (3.36). We have therefore shown that the U$`(p)`$ gauge theory (3.10) on a bundle of topological charges $`Q_{ab}`$ and ’t Hooft fluxes $`\gamma _{ab}`$ over the noncommutative torus $`\text{𝕋}_\mathrm{\Theta }^D`$ (i.e. with gauge fields obeying twisted boundary conditions) is equivalent to the U$`(\stackrel{~}{p}_0)`$ gauge theory (3.51) on a trivial bundle over the noncommutative torus $`\text{𝕋}_\mathrm{\Theta }^{}^D`$ (i.e. with gauge fields obeying periodic boundary conditions) with noncommutativity parameter matrix defined by (3.46) and the reduced rank $`\stackrel{~}{p}_0`$ by (3.29). In particular, for $`\stackrel{~}{p}_0=1`$, the internal matrix structure of the gauge fields are completely absorbed into the operators $`\widehat{Z}_a^{}`$ which are regarded as the coordinate generators of $`\text{𝕋}_\mathrm{\Theta }^{}^D`$, thereby allowing one to reinterpret the original non-abelian gauge theory as a U(1) gauge theory on the noncommutative torus $`\text{𝕋}_\mathrm{\Theta }^{}^D`$ (Note that this statement is true even for $`\theta =0`$). On the other hand, when $`Q=0`$ we have $`\stackrel{~}{p}_0=p`$ and the rank of the gauge fields is unchanged. Note that from (3.46) and (3.47) it follows that, when $`\theta =0`$, the resulting dual gauge theory does not depend on $`\gamma `$ and $`\alpha `$ separately, but only on their combination given by $`Q`$ in (3.5). This means that keeping both $`\gamma `$ and $`\alpha `$ non-zero is in fact redundant in the commutative case. But for $`\theta 0`$, this is not quite the case. We note that eq. (3.40) and the antisymmetry of the matrices $`AB^{}`$ and $`\stackrel{~}{Q}^{}\stackrel{~}{P}`$ implies the matrix identity $$\left(\begin{array}{cc}A& B\\ \stackrel{~}{Q}^{}& \stackrel{~}{P}^{}\end{array}\right)\left(\begin{array}{cc}0& 11_D\\ 11_D& 0\end{array}\right)\left(\begin{array}{cc}A& B\\ \stackrel{~}{Q}^{}& \stackrel{~}{P}^{}\end{array}\right)^{}=\left(\begin{array}{cc}0& 11_D\\ 11_D& 0\end{array}\right).$$ (3.54) It follows that, when $`\gamma =0`$, the map $`\mathrm{\Theta }\mathrm{\Theta }^{}`$ in (3.46) is the usual SO$`(D,D;\text{})`$ transformation that relates Morita equivalent noncommutative tori . From the present point of view, Morita equivalence is therefore regarded as a change of basis $`\widehat{\mathrm{\Delta }}(x)\widehat{\mathrm{\Delta }}^{}(x^{})`$ for the mapping between Weyl operators and fields. The choice $`\gamma =0`$ is required whenever one wants to represent the operators $`\widehat{Z}_a^{}`$ on the Hilbert space of sections of the associated fundamental bundle, as in the case of D-brane applications wherein the open string wavefunctions transform in the fundamental representation of the Chan-Paton gauge group . A choice $`\gamma 0`$ is possible for representations on the Hilbert space of sections of the associated principal bundle. Notice also that for $`d>1`$ the calculation presented above accounts for only a subset of the possible noncommutative gauge theories, since in that instance, even in the commutative case, generic bundles over the torus do not always admit constant curvature connections such as (3.22). However, the derivation above does allow for arbitrary non-abelian ’t Hooft fluxes $`\gamma `$. ### 3.3 Star-gauge symmetry and transformation of observables We will now demonstrate explicitly how the star-gauge symmetries of two Morita equivalent Yang-Mills theories, as well as their star-gauge invariant observables, are related to each other. We begin with the U$`(\stackrel{~}{p}_0)`$ theory (3.20) in which all fields obey periodic boundary conditions. Its gauge symmetry is $$\widehat{𝒜}_\mu \widehat{g}\widehat{𝒜}_\mu \widehat{g}^{}i\widehat{g}[\widehat{}_\mu ^{(0)},\widehat{g}^{}],$$ (3.55) where $`\widehat{g}`$ may be written as $$\widehat{g}=\text{d}^Dx^{}\widehat{\mathrm{\Delta }}^{}(x^{})g^{}(x^{}).$$ (3.56) The gauge function $`g^{}(x^{})`$ is a single-valued, $`\stackrel{~}{p}_0\times \stackrel{~}{p}_0`$ star($`^{}`$)-unitary matrix field, which parametrizes the usual gauge transformation of the noncommutative Yang-Mills theory (3.51). To see how this is interpreted as the gauge symmetry of the U($`p`$) theory (3.19), we use (3.22) to rewrite (3.55) as $$\widehat{A}_\mu \widehat{g}\widehat{A}_\mu \widehat{g}^{}i\widehat{g}[\widehat{}_\mu ,\widehat{g}^{}],$$ (3.57) where $`\widehat{A}_\mu `$ is given by (3.18). The invariance of the action (3.19) under (3.57) follows by construction. Note that by the definition (3.56), the operator $`\widehat{g}`$ obeys $$\text{e}^{\mathrm{\Sigma }_{\mu a}\widehat{}_\mu }\widehat{g}\text{e}^{\mathrm{\Sigma }_{\mu a}\widehat{}_\mu }=\widehat{\mathrm{\Omega }}_a\widehat{g}\widehat{\mathrm{\Omega }}_a^{},$$ (3.58) from which it follows that the gauge transformed operator in (3.57) also satisfies the required twisted boundary condition (3.23). Let us rewrite (3.58) in terms of fields on the noncommutative torus $`\text{𝕋}_\mathrm{\Theta }^D`$. For this, we expand the Weyl operator $`\widehat{g}`$ using the basis $`\widehat{\mathrm{\Delta }}(x)`$ as $$\widehat{g}=\text{d}^Dx\widehat{\mathrm{\Delta }}(x)g(x).$$ (3.59) Due to (3.58), the $`p\times p`$ star($``$)-unitary matrix field $`g(x)`$ is multi-valued and defines an adjoint section of the gauge bundle, $$g(x+\mathrm{\Sigma }_{\mu a}\widehat{\mu })=\mathrm{\Omega }_a(x)g(x)\mathrm{\Omega }_a(x)^{}.$$ (3.60) Let us now consider star($`^{}`$)-gauge invariant observables in the U$`(\stackrel{~}{p}_0)`$ theory. Since the fields of this theory are single-valued functions on the torus, we may use the construction of Sections 2.3 and 2.4. Given the parallel transport operators $$𝒰^{}(x^{};C)=\mathrm{P}\mathrm{exp}_{^{}}\left(i\underset{C}{}\text{d}\xi ^\mu 𝒜_\mu ^{}(x^{}+\xi )\right),$$ (3.61) where notation is as in Section 2.3, we may define a star($`^{}`$)-gauge invariant observable by $$𝒪(C)=\text{d}^Dx^{}\mathrm{tr}_{(\stackrel{~}{p}_0)}\left(𝒰^{}(x^{};C)^{}S_v^{}(x^{})\right),$$ (3.62) where $$S_v^{}(x^{})=\text{e}^{ik_\mu ^{}x_\mu ^{}}\mathrm{\hspace{0.17em}1}1_{\stackrel{~}{p}_0}.$$ (3.63) The total momentum $`k_\mu ^{}`$ is quantized as $$k_\mu ^{}=2\pi \left(\mathrm{\Sigma }^{}_{}{}^{}1\right)_{a\mu }m_a^{}$$ (3.64) while the constraint on the relative separation vector between the two ends of the contour $`C`$ is $$v_\mu =\theta _{\mu \nu }^{}k_\nu ^{}+\mathrm{\Sigma }_{\mu a}^{}n_a^{},$$ (3.65) for some integer-valued vectors $`m_a^{}`$ and $`n_a^{}`$. We now rewrite the above observables in terms of Weyl operators. Introducing $`\widehat{U}^{}(C)`$ $`=`$ $`\mathrm{P}\mathrm{exp}\left({\displaystyle \underset{C}{}}\text{d}\xi ^\mu \left(\widehat{}_\mu ^{(0)}+i\widehat{𝒜}_\mu \right)\right),`$ $`\widehat{D}^{}(C)`$ $`=`$ $`\mathrm{P}\mathrm{exp}\left({\displaystyle \underset{C}{}}\text{d}\xi ^\mu \widehat{}_\mu ^{(0)}\right),`$ (3.66) we can define a star($`^{}`$)-gauge invariant observable as in (2.37) by $$𝒪(C)=\text{Tr }^{}\mathrm{tr}_{(\stackrel{~}{p}_0)}\left(\widehat{U}^{}(C)\widehat{D}^{}(C)^{}\widehat{S}_v\right),$$ (3.67) where the unitary operator $`\widehat{S}_v`$ may be expanded as $$\widehat{S}_v=\text{d}^Dx^{}\widehat{\mathrm{\Delta }}^{}(x^{})S_v^{}(x^{}),$$ (3.68) and it satisfies $$\widehat{S}_v\widehat{\mathrm{\Delta }}^{}(x^{})\widehat{S}_v^{}=\widehat{\mathrm{\Delta }}^{}(x^{}+v)11_{\stackrel{~}{p}_0}.$$ (3.69) Let us now see how to interpret the operator (3.67) as a sensible quantity also in the Morita equivalent U($`p`$) theory. We first define a parallel transport operator $`𝒰(x;C)`$ in the U($`p`$) theory using Eqs. (2.31) and (2.32). Note that $`\widehat{U}(C)=\widehat{U}^{}(C)`$. The observable defined in (3.67) can then be rewritten as $$𝒪(C)=\text{Tr }\mathrm{tr}_{(p)}\left[\left(\widehat{U}(C)\widehat{D}(C)^{}\right)\left(\widehat{D}(C)\widehat{D}^{}(C)^{}\right)\widehat{S}_v\right].$$ (3.70) We now expand the operator $`\widehat{S}_v`$ in terms of the basis $`\widehat{\mathrm{\Delta }}(x)`$ similarly to (3.59) and define the $`p\times p`$ star($``$)-unitary matrix field $`S_v(x)`$ as an adjoint section. Then, using (2.31), (3.66) and (3.22), we can write the quantity (3.70) as $$𝒪(C)=\text{d}^Dx\mathrm{tr}_{(p)}\left(𝒰(x;C)S_v^{(0)}(x)\right),$$ (3.71) where we have defined the $`p\times p`$ star($``$)-unitary matrix field $$S_v^{(0)}(x)=\left[\mathrm{P}\mathrm{exp}_{}\left(i\underset{C}{}\text{d}\xi ^\mu A_\mu ^{(0)}(x+\xi )\right)\right]^{}S_v(x).$$ (3.72) By construction, (3.71) is star($``$)-gauge invariant, which can also be checked by recalling the star($``$)-gauge transformation (2.29) of the parallel transport operator $`𝒰(x;C)`$ and the translation-generating property of the field (3.72), $$S_v^{(0)}(x)g(x)S_v^{(0)}(x)^{}=g(x+v),$$ (3.73) where $`g(x)`$ is an arbitrary adjoint section as in (3.60). Note that the integrand of (3.71) is a single-valued function because of the boundary conditions $`𝒰(x+\mathrm{\Sigma }_{\mu a}\widehat{\mu };C)`$ $`=`$ $`\mathrm{\Omega }_a(x)𝒰(x;C)\mathrm{\Omega }_a(x+v)^{},`$ (3.74) $`S_v^{(0)}(x+\mathrm{\Sigma }_{\mu a}\widehat{\mu })`$ $`=`$ $`\mathrm{\Omega }_a(x+v)S_v^{(0)}(x)\mathrm{\Omega }_a(x)^{},`$ (3.75) where the relation (3.75) is required for consistency of the equations (3.73) and (3.60). It follows that star($`^{}`$)-gauge invariant observables of the U($`\stackrel{~}{p}_0`$) theory can be interpreted as star($``$)-gauge invariant observables of the U($`p`$) theory. As in Section 2.3, we can show that the field $`S_v^{(0)}(x)`$ satisfying (3.73) on the noncommutative torus $`\text{𝕋}_\mathrm{\Theta }^D`$ is in fact unique, so that the Morita equivalence relation provides a one-to-one mapping between the quantum correlation functions of the two noncommutative Yang-Mills theories. An instructive example which shows how this correspondence works is the $`\theta =0`$ case. Then, one finds that Polyakov lines of commutative gauge theory map to open loops in the Morita equivalent noncommutative gauge theory. The observables (3.71,3.72) are the appropriate modifications of the observables constructed in Sections 2.3 and 2.4 to the case of multi-valued gauge fields on the noncommutative torus. We remark that the observables constructed here using the formalism of Weyl operators differ from the noncommutative holonomy operators constructed recently in which are star-gauge invariant without the need of introducing the operator trace Tr, but which do not immediately generalize to gauge field configurations of non-vanishing magnetic flux. The present noncommutative Wilson lines are defined for multi-valued gauge fields, and moreover, as we will see in Section 6, they are the ones which arise in the effective actions for charges propagating in the background of noncommutative gauge fields. ## 4 Lattice regularization In this Section we will consider a lattice regularization of noncommutative field theories. Through a general analysis of the noncommutative algebra generated by the spacetime coordinates, we will find that the discretization of the spacetime inevitably requires that it be compact in some restricted way depending on the noncommutativity. Due to this restriction, the commutative limit is not commutable with the continuum limit, which demonstrates the UV/IR mixing property of noncommutative field theories at a fully nonperturbative level. We will show that the discrete formulation allows a nonperturbative regularization of generic noncommutative field theories with single-valued fields. Given the Morita equivalence property discussed in the previous Section, this means that noncommutative Yang-Mills theories with gauge fields obeying twisted boundary conditions can also be regularized in terms of the Morita equivalent Yang-Mills theories with fields satisfying periodic boundary conditions, insofar as defining regularized correlation functions of star-gauge invariant observables are concerned. In this way, we obtain a nonperturbative definition of noncommutative Yang-Mills theories in arbitrary even dimensions with multi-valued or single-valued gauge fields. ### 4.1 General construction To describe the lattice regularization of noncommutative Yang-Mills theory, there are some subtle technical points concerning the transcription of the continuum theory onto a lattice that we first need to address. To this end, let us go back to the construction of Section 2.1. To define a lattice field theory, we restrict the spacetime points to $`x_\mu ϵ\text{}`$, where $`ϵ`$ is the lattice spacing. It follows that the lattice momentum must be identified under the shift $$k_\mu k_\mu +\frac{2\pi }{ϵ}\delta _{\mu \nu },\nu =1,\mathrm{},D.$$ (4.1) Correspondingly, there is the operator identity $$\text{e}^{i(k_\mu +\frac{2\pi }{ϵ}\delta _{\mu \nu })\widehat{x}_\mu }=\text{e}^{ik_\mu \widehat{x}_\mu },\nu =1,\mathrm{},D.$$ (4.2) By acting with the operator $`\text{e}^{ik_\mu \widehat{x}_\mu }`$ on both sides of (4.2), we find that there is also the operator identity $$\text{e}^{2\pi i\widehat{x}_\mu /ϵ}=11,\mu =1,\mathrm{},D,$$ (4.3) and, moreover, that the momentum $`k_\mu `$ is quantized according to $$\theta _{\mu \nu }k_\nu 2ϵ\text{}.$$ (4.4) For $`k_\mu `$ satisfying (4.4), the commutation relation $$\text{e}^{ik_\mu \widehat{x}_\mu }\text{e}^{2\pi i\widehat{x}_\nu /ϵ}=\text{e}^{2\pi i\theta _{\mu \nu }k_\mu /ϵ}\text{e}^{2\pi i\widehat{x}_\nu /ϵ}\text{e}^{ik_\mu \widehat{x}_\mu },\mu =1,\mathrm{},D$$ (4.5) is compatible with the identity (4.3). On the other hand, the compatibility of the commutation relation $$\text{e}^{v_\nu \widehat{}_\nu }\text{e}^{2\pi i\widehat{x}_\mu /ϵ}\text{e}^{v_\nu \widehat{}_\nu }=\text{e}^{2\pi i(\widehat{x}_\mu +v_\mu )/ϵ}$$ (4.6) with the identity (4.3) requires $`v_\mu ϵ\text{}`$. Thus we are led to use instead the lattice shift operator $$\widehat{D}_\mu =\text{e}^{ϵ\widehat{}_\mu },$$ (4.7) but not the derivation $`\widehat{}_\mu `$ itself anymore. The restriction to the lattice shift operator (4.7) is standard in lattice field theory, whereas the restriction (4.4) on the lattice momentum simply disappears in the commutative case $`\theta _{\mu \nu }=0`$ and is quite characteristic of the noncommutative geometry. The momentum quantization (4.4) requires that the spacetime be compactified in some restricted way. Let us consider the case when the fields have the periodic boundary conditions $$\varphi (x+\mathrm{\Sigma }_{\mu a}\widehat{\mu })=\varphi (x),a=1,\mathrm{},D,$$ (4.8) where the periods $`\mathrm{\Sigma }_{\mu a}`$ are integer multiples of the lattice spacing $`ϵ`$. Because of (4.8), the momentum is quantized as $$k_\mu =2\pi (\mathrm{\Sigma }^1)_{a\mu }m_a,m_a\text{}.$$ (4.9) The momentum periodicity (4.1) due to the lattice discretization $`x_\mu ϵ\text{}`$ can be recast in terms of the integers $`m_a`$ as $$m_am_a+\frac{1}{ϵ}\mathrm{\Sigma }_{\mu a},\mu =1,\mathrm{},D.$$ (4.10) Now the restriction on the momentum given by (4.4) implies that there exists a $`D\times D`$ integer-valued matrix $`M_{\mu a}`$ which satisfies $$M_{\mu a}\mathrm{\Sigma }_{\nu a}=\frac{\pi }{ϵ}\theta _{\mu \nu }.$$ (4.11) We have therefore discovered the remarkable fact that lattice regularization of noncommutative field theory forces the spacetime to be compact. For fixed $`M`$, say $`M=11_D`$, the infrared cutoff disappears as $`\frac{1}{ϵ}`$ in the continuum limit $`ϵ0`$. Note that the commutative limit $`\theta _{\mu \nu }0`$ does not commute with the continuum limit $`ϵ0`$. Thus the UV/IR mixing discovered in perturbative analyses of noncommutative field theories is demonstrated here at a fully nonperturbative level. On the other hand, in order to obtain a continuum spacetime of finite volume one has to send $`M\mathrm{}`$ as one takes the $`ϵ0`$ limit. It is important to note that for any given period matrix $`\mathrm{\Sigma }_{\mu a}`$ and noncommutativity parameter $`\theta _{\mu \nu }`$ in the continuum, one can construct a family of lattice geometries satisfying the restriction (4.11) and approaching the target continuum theory in the $`ϵ0`$ limit. This means that in spite of the restriction (4.11) for the regularized field theories, one can obtain the most general noncommutative geometries parametrized by $`\mathrm{\Sigma }_{\mu a}`$ and $`\theta _{\mu \nu }`$ in the continuum limit. Due to the momentum quantization (4.9), we are led to use the coordinate operators $$\widehat{Z}_a=\text{e}^{2\pi i(\mathrm{\Sigma }^1)_{a\mu }\widehat{x}_\mu },$$ (4.12) but not the $`\widehat{x}_\mu `$ themselves anymore. The commutation relations of the operators (4.12) and (4.7) are $`\widehat{Z}_a\widehat{Z}_b`$ $`=`$ $`\text{e}^{2\pi i\mathrm{\Theta }_{ab}}\widehat{Z}_b\widehat{Z}_a`$ (4.13) $`\widehat{D}_\mu \widehat{Z}_a\widehat{D}_\mu ^{}`$ $`=`$ $`\text{e}^{2\pi iϵ(\mathrm{\Sigma }^1)_{a\mu }}\widehat{Z}_a,`$ (4.14) where the dimensionless noncommutativity parameter $$\mathrm{\Theta }_{ab}=2\pi \left(\mathrm{\Sigma }^1\right)_{a\mu }\theta _{\mu \nu }\left(\mathrm{\Sigma }^1\right)_{b\nu }$$ (4.15) is necessarily rational-valued on the lattice, since the restriction (4.11) implies that $$M_{\mu a}=\frac{1}{2ϵ}\mathrm{\Sigma }_{\mu b}\mathrm{\Theta }_{ba}$$ (4.16) is an integer-valued matrix. Because of (4.3), we also have the identity $$\underset{a=1}{\overset{D}{}}\left(\widehat{Z}_a\right)^{\frac{1}{ϵ}\mathrm{\Sigma }_{\mu a}}=\mathrm{exp}\left[\pi i\underset{a<b}{}\mathrm{\Sigma }_{\mu a}\mathrm{\Theta }_{ab}\mathrm{\Sigma }_{\mu b}\right]11_D,\mu =1,\mathrm{},D.$$ (4.17) By construction, the identity (4.17) is compatible with the commutation relations (4.13) and (4.14). We can define a map $`\widehat{\mathrm{\Delta }}(x)`$ between lattice fields and Weyl operators by $$\widehat{\mathrm{\Delta }}(x)=\frac{1}{\left\|det\frac{1}{ϵ}\mathrm{\Sigma }\right\|}\underset{\stackrel{}{m}}{}\left(\underset{a=1}{\overset{D}{}}\left(\widehat{Z}_a\right)^{m_a}\right)\text{e}^{\pi i_{a<b}\mathrm{\Theta }_{ab}m_am_b}\text{e}^{2\pi i(\mathrm{\Sigma }^1)_{a\mu }m_ax_\mu },$$ (4.18) where the sum goes over $`\text{}^D`$ modulo the periodicity (4.10) and $`x_\mu `$ is a point on the spacetime lattice $`ϵ\text{}`$ with periodicity (4.8). Note that $$\frac{1}{\left\|det\frac{1}{ϵ}\mathrm{\Sigma }\right\|}\underset{\stackrel{}{m}}{}\text{e}^{2\pi i(\mathrm{\Sigma }^1)_{a\mu }m_ax_\mu }=\delta _{x,0(\text{mod}\mathrm{\Sigma })}.$$ (4.19) As we did in (2.14), we can define the lattice star-product using the map $`\widehat{\mathrm{\Delta }}(x)`$ defined in (4.18). Explicitly it can be given as $$\varphi _1(x)\varphi _2(x)=\underset{y,z}{}K(xy,xz)\varphi _1(y)\varphi _2(z),$$ (4.20) where the sums go over the spacetime lattice points $`ϵ\text{}`$ modulo the periodicity (4.8). The kernel $`K`$ in (4.20) is given by $$K(y,z)=\frac{1}{\left\|det\frac{1}{ϵ}\mathrm{\Sigma }\right\|^2}\underset{\stackrel{}{m},\stackrel{}{n}}{}\mathrm{exp}\left[2\pi i(\mathrm{\Sigma }^1)_{a\mu }(m_ay_\mu +n_az_\mu )+i\pi \mathrm{\Theta }_{ab}m_an_b\right].$$ (4.21) The map (4.18) and the star-product (4.20,4.21) possess all the algebraic properties that their continuum counterparts have. If $`2ϵ(\mathrm{\Sigma }\mathrm{\Theta })^1=M^1`$ is an integer-valued matrix, where $`M`$ is the integral matrix introduced in (4.11), then the sums in (4.21) can be done explicitly yielding $$K(y,z)=\frac{1}{\left\|det\frac{1}{ϵ}\mathrm{\Sigma }\right\|}\text{e}^{2i(\theta ^1)_{\mu \nu }y_\mu z_\nu }.$$ (4.22) The corresponding formula (4.20) for the lattice star-product is then analogous to the second expression in (2.14) for the continuum star-product. Let us now turn to the construction of field theories on a discrete noncommutative torus. In the case of a scalar lattice field $`\varphi (x)`$, one can define the operator $$\widehat{\varphi }=\underset{x}{}\varphi (x)\widehat{\mathrm{\Delta }}(x)$$ (4.23) where the sum runs over lattice points. We may then write down an action $$S\left[\widehat{\varphi }\right]=\text{Tr }\left(\frac{1}{2}\underset{\mu }{}\left(\widehat{D}_\mu \widehat{\varphi }\widehat{D}_\mu ^{}\widehat{\varphi }\right)^2+\frac{1}{2}\widehat{\varphi }^2+\frac{1}{4!}\widehat{\varphi }^4\right),$$ (4.24) which leads to the usual lattice action of the field $`\varphi (x)`$ with star-interaction term analogous to (LABEL:four-int). A huge advantage of the lattice regularization is that it allows one to construct finite dimensional representations of the noncommutative geometry. This is apparent already at the level of the lattice operators (4.18). It is straightforward to compute that they obey the commutation relations $$[\widehat{\mathrm{\Delta }}(x),\widehat{\mathrm{\Delta }}(y)]=\underset{z}{}𝒦(xz,yz)\widehat{\mathrm{\Delta }}(z),$$ (4.25) where $$𝒦(x,y)=\frac{2i}{\left\|det\frac{1}{ϵ}\mathrm{\Sigma }\right\|^2}\underset{\stackrel{}{m},\stackrel{}{n}}{}\mathrm{sin}\left(\pi \mathrm{\Theta }_{ab}m_an_b\right)\text{e}^{2\pi i(\mathrm{\Sigma }^1)_{a\mu }(m_ax_\mu +n_ay_\mu )}.$$ (4.26) The lattice operators $`\widehat{\mathrm{\Delta }}(x)`$ therefore generate a finite dimensional Lie algebra of dimension $`\|det\frac{1}{ϵ}\mathrm{\Sigma }\|=2^D\|detM/det\mathrm{\Theta }\|`$. In the lattice field theory, the operators $`\widehat{\mathrm{\Delta }}(x)`$ can in this way be regarded as maps from the $`\|det\frac{1}{ϵ}\mathrm{\Sigma }\|`$ lattice points to a finite dimensional matrix representation, thereby illustrating how the degrees of freedom are mapped into each other in the one-to-one correspondence between Weyl operators and fields. Indeed, the algebra (4.13) can always be represented by finite dimensional matrices which admit a tensor product decomposition into blocks depending on the detailed forms of the rational numbers $`\mathrm{\Theta }_{ab}`$ (See for some specific examples). We shall return to this point in the next Section. The lattice regularization can therefore also be thought of as an approximation to the algebra of Weyl operators in the continuum by finite dimensional matrix algebras. This constitutes the standard description of noncommutative lattices by approximately finite algebras in noncommutative geometry . The rigorous definition of the continuum limit described above within such an algebraic description of the noncommutative geometry is given in Ref. . ### 4.2 Noncommutative Yang-Mills theory on the lattice In order to construct a lattice regularization of noncommutative Yang-Mills theory, we need to maintain star-gauge invariance on the lattice. As in the case of ordinary lattice gauge theory , this is achieved by putting the gauge fields on the links of the lattice , $$\widehat{U}_\mu =\underset{x}{}\widehat{\mathrm{\Delta }}(x)U_\mu (x),$$ (4.27) where $`\widehat{U}_\mu `$ is a unitary operator and $`U_\mu (x)`$ is a $`p\times p`$ matrix field on the lattice which is star-unitary, $$U_\mu (x)U_\mu (x)^{}=11_p.$$ (4.28) One can write an action $$S=\frac{1}{g^2}\underset{\mu \nu }{}\text{Tr }\mathrm{tr}_{(p)}\left[\widehat{U}_\mu \left(\widehat{D}_\mu \widehat{U}_\nu \widehat{D}_\mu ^{}\right)\left(\widehat{D}_\nu \widehat{U}_\mu ^{}\widehat{D}_\nu ^{}\right)\widehat{U}_\nu ^{}\right]$$ (4.29) which is invariant under the transformation $$\widehat{U}_\mu \widehat{g}\widehat{U}_\mu \left(\widehat{D}_\mu \widehat{g}^{}\widehat{D}_\mu ^{}\right).$$ (4.30) The action (4.29) can be written in terms of the lattice fields $`U_\mu (x)`$ as $$S=\frac{1}{g^2}\underset{x}{}\underset{\mu \nu }{}\mathrm{tr}_{(p)}\left[U_\mu (x)U_\nu (x+ϵ\widehat{\mu })U_\mu (x+ϵ\widehat{\nu })^{}U_\nu (x)^{}\right],$$ (4.31) which is invariant under the lattice star-gauge transformation $$U_\mu (x)g(x)U_\mu (x)g(x+ϵ\widehat{\mu })^{},$$ (4.32) where the gauge function $`g(x)`$ is defined by $$\widehat{g}=\underset{x}{}\widehat{\mathrm{\Delta }}(x)g(x)$$ (4.33) and it is star-unitary, $`g(x)g(x)^{}=11_p`$. This discrete version of noncommutative Yang-Mills theory has several benefits. First of all, it provides a concrete definition of the quantum gauge theory path integral in the continuum. The unitary operators (4.27) live in a finite-dimensional operator algebra and are therefore elements of a compact unitary group. Note that the trace $`\text{Tr }\mathrm{tr}_{(p)}`$ which appears in (4.29) corresponds to the trace in the fundamental representation of this unitary group. The lattice gauge theory path integral may then be defined by the measure $`\text{d}\widehat{U}`$ which is the Haar measure of this compact Lie group that is invariant under the left and right actions $$\widehat{U}_\mu \widehat{g}\widehat{U}_\mu ,\widehat{U}_\mu \widehat{U}_\mu \widehat{g}.$$ (4.34) In terms of the star-unitary lattice fields $`U_\mu (x)`$, the path integral measure $`𝒟U_\mu (x)`$ is given by the Haar measure of the compact Lie group which is formed by the $`U_\mu (x)`$ with multiplication given by the star-product, and which is invariant under the left and right actions $`U_\mu (x)g(x)U_\mu (x)`$ and $`U_\mu (x)U_\mu (x)g(x)`$. This measure preserves star-gauge invariance. In the commutative limit $`\theta _{\mu \nu }=0`$, the discrete noncommutative Yang-Mills theory reduces to ordinary lattice gauge theory , including the path integral measure because of the uniqueness of the Haar measure. The lattice regularization also gives an approximation to the star-gauge symmetry group in the continuum. An algebraic description of this continuum gauge group is given in Ref. . Let us now describe the discrete analogs of the star-gauge invariant observables introduced in Section 2.3. We first define the discrete analog of the parallel transport operator $`𝒰(x;C)`$. We introduce an oriented contour $`C`$ on the lattice specified by a collection of links, $$C=\{\mu _1,\mu _2,\mathrm{},\mu _n\},$$ (4.35) where $`\mu _j=\pm 1,\pm 2,\mathrm{},\pm D`$. We define $`U_\mu (x)=U_\mu (xϵ\widehat{\mu })^{}`$. The parallel transport operator can then be defined as $$𝒰(x;C)=U_{\mu _1}(x)U_{\mu _2}(x+ϵ\widehat{\mu }_1)\mathrm{}U_{\mu _n}\left(x+ϵ\underset{j=1}{\overset{n1}{}}\widehat{\mu }_j\right).$$ (4.36) As in the continuum, in order to construct star-gauge invariant observables out of the operator $`𝒰(x;C)`$, we need a star-unitary function $`S_v(x)`$ with the property (2.25) for arbitrary functions $`g(x)`$ on the periodic lattice, where $$v=ϵ\underset{j=1}{\overset{n}{}}\widehat{\mu }_j.$$ (4.37) One finds that a necessary and sufficient condition is given again by (2.46) with the solution $`S_v(x)=\text{e}^{ik_\mu x_\mu }\mathrm{\hspace{0.17em}1}1_p`$, where the loop momentum $`k_\mu `$ satisfies (2.47). In the present case, both $`k`$ and $`v`$ are quantized and periodic. Modding out by the periodicity, there are only the same, finite number of values that $`k`$ and $`v`$ can take. Therefore, it makes sense to ask whether or not (2.47) gives a one-to-one correspondence between $`k`$ and $`v`$. The answer is affirmative if and only if there exist $`D\times D`$ integer-valued matrices $`J`$ and $`K`$ which satisfy $$\frac{1}{ϵ}\mathrm{\Sigma }J2MK=11_D,$$ (4.38) where $`M`$ is the integral matrix introduced in (4.11). If this condition is not met, then there exist loop separation vectors $`v`$ for which there is no momentum $`k`$ satisfying (2.47), and for the other $`v`$ there is more than one value of $`k`$ satisfying (2.47). For example, let us consider the case with period matrix $`\frac{1}{ϵ}\mathrm{\Sigma }=L11_D`$. If $`L`$ is an even integer, then (4.38) cannot be satisfied. If $`L`$ is odd and $`M^1`$ is an integral matrix, then one can satisfy (4.38) by taking $`J=11_D`$ and $`K=\frac{L1}{2}M^1`$. Thus, for periodic gauge fields, one can construct a lattice regularization of noncommutative Yang-Mills theory with arbitrary gauge group rank $`p`$, period $`\mathrm{\Sigma }`$, and noncommutativity parameter $`\theta `$. Due to Morita equivalence in the continuum, this means that we have a nonperturbative formulation of the general class of (twisted) continuum gauge theories that were described in Section 3.1. In this regard, it is important that Morita equivalence also holds at the level of observables as we discussed in Section 3.3. The problem with defining a discrete noncommutative U$`(p)`$ gauge theory with multi-valued gauge fields directly is that one immediately encounters an obstacle to constructing a background abelian gauge field (3.6) on the lattice. As we discussed above, the relation (4.4) imposes a quantization constraint on the allowed momenta in the lattice discretization of noncommutative geometry. This would imply that the background tensor field $`F`$ be quantized proportionally to $`\theta ^1`$, and so the configuration (3.6) cannot be constructed on the lattice in general. However, as we have discussed, we circumvent this difficulty by using the noncommutative lattice gauge theory with single-valued gauge fields and Morita equivalence in the continuum to obtain a non-perturbative definition of all noncommutative Yang-Mills theories associated with constant curvature vector bundles over tori. Other problems with discretizing noncommutative geometry, such as the construction of an appropriate lattice Dirac operator, are discussed in Ref. . ## 5 Explicit realizations of discrete noncommutative gauge theory The construction of the previous Section has been quite general, and we will now describe some concrete examples of discrete noncommutative Yang-Mills theory. Within the lattice formalism, we will demonstrate the Morita equivalence between commutative U($`p`$) lattice gauge theory with fields obeying twisted boundary conditions and a noncommutative U($`\stackrel{~}{p}_0`$) lattice gauge theory with fields obeying periodic boundary conditions. Using this property, we will further show that noncommutative gauge theory can be regularized by means of commutative lattice gauge theory with ’t Hooft flux. As a special case, this construction includes a previous proposal for a concrete definition of noncommutative Yang-Mills theory using large $`N`$ reduced models . This will further lead to the explicit finite dimensional matrix representations of the noncommutative torus that were discussed in the previous Section. ### 5.1 Discrete Morita equivalence In Section 4.2, we have shown that noncommutative Yang-Mills theories can be nonperturbatively regularized. The regularized theory is described by star-unitary gauge fields (4.28) on a lattice and the action (4.31) is defined in terms of lattice star-products. Although the theory is explicitly given in terms of a finite number of degrees of freedom, as it stands it is not very suitable for practical purposes, say for numerical studies such as Monte Carlo simulations. In this Subsection, we show further that noncommutative Yang-Mills theory can be nonperturbatively regularized by means of commutative lattice gauge theories with multi-valued gauge fields . For this, we consider the lattice analog of Morita equivalence which we derived in Section 3.2. As we discussed in Section 4.2, there is a technical obstruction to constructing noncommutative Yang-Mills theory directly on the lattice for gauge fields of non-vanishing topological charge. We recall that the only obstacle was that the constant abelian background gauge field (3.6) does not generally have a momentum compatible with the restriction (4.4). This obstacle disappears, of course, for the commutative case $`\theta _{\mu \nu }=0`$. What we will prove in the lattice formulation is the Morita equivalence between commutative Yang-Mills theory with gauge fields obeying twisted boundary conditions and noncommutative Yang-Mills theory with periodic gauge fields, in arbitrary even dimension $`D=2d`$. In fact, we will find that, for a given noncommutative Yang-Mills theory with arbitrary period matrix $`\mathrm{\Sigma }`$ and deformation parameter $`\mathrm{\Theta }`$ in the continuum, we can construct a family of commutative U($`p`$) lattice gauge theories with ’t Hooft flux whose sequence of Morita equivalent noncommutative theories converges to the target noncommutative field theory in the continuum limit. When $`\mathrm{\Theta }`$ is irrational, the rank $`p`$ of the gauge group of the commutative theory must be sent to infinity as one takes the continuum limit. Combining this with the continuum Morita equivalence which connects two noncommutative Yang-Mills theories with periodic and twisted boundary conditions on the gauge fields, we will find that continuum noncommutative Yang-Mills theories, with either rational or irrational noncommutativity parameters and with or without twists in the boundary conditions on the gauge fields, can be regularized by means of commutative lattice gauge theory with ’t Hooft flux. We start with a commutative U$`(p)`$ lattice gauge theory with fields obeying twisted boundary conditions. The action is $$S=\frac{1}{g^2}\underset{x}{}\underset{\mu \nu }{}\mathrm{tr}_{(p)}\left[U_\mu (x)U_\nu (x+ϵ\widehat{\mu })U_\mu (x+ϵ\widehat{\nu })^{}U_\nu (x)^{}\right],$$ (5.1) where $`U_\mu (x)`$ are U($`p`$) gauge fields satisfying the twisted boundary conditions $$U_\mu (x+\mathrm{\Sigma }_{\nu a}\widehat{\nu })=\mathrm{\Omega }_a(x)U_\mu (x)\mathrm{\Omega }_a(x+ϵ\widehat{\mu })^{}$$ (5.2) with period matrix $`\mathrm{\Sigma }`$. The transition functions $`\mathrm{\Omega }_a(x)`$ are U($`p`$) matrices which, for consistency of the constraints (5.2), must satisfy the cocycle condition $$\mathrm{\Omega }_a(x+\mathrm{\Sigma }_{\mu b}\widehat{\mu })\mathrm{\Omega }_b(x)=𝒵_{ab}\mathrm{\Omega }_b(x+\mathrm{\Sigma }_{\mu a}\widehat{\mu })\mathrm{\Omega }_a(x),$$ (5.3) where $`𝒵_{ab}=\text{e}^{2\pi i\gamma _{ab}/p}\text{}_p`$. The antisymmetric matrix $`\gamma `$ has elements $`\gamma _{ab}\text{}`$ representing the ’t Hooft fluxes. As an explicit form of the transition functions, let us take the gauge choice $$\mathrm{\Omega }_a(x)=1\mathrm{\Gamma }_a,$$ (5.4) where $`\mathrm{\Gamma }_a`$ are the SU$`(p)`$ twist eaters obeying the Weyl-’t Hooft commutation relations (3.4). In order to satisfy (5.3) we must have $`Q=\gamma `$. Note that the gauge choice (5.4) corresponds to setting $`\alpha =0`$ in the continuum expression (3.2). The reason for this is that, as we discussed in Section 3.1, keeping both $`\alpha `$ and $`\gamma `$ non-zero in the commutative case is redundant. In other words, the abelian magnetic flux of a U$`(p)`$ gauge field, which plays a very important role in Morita equivalences of noncommutative gauge theories , arises in the commutative case only via the corresponding ’t Hooft flux . A similar consideration with non-vanishing $`\alpha `$ would lead us to the same final results. Due to the constraint (5.2), the gauge theory can be expressed in terms of lattice gauge fields $`U_\mu (x)`$ with $`x_\mu `$ lying in a unit cell of period $`\mathrm{\Sigma }`$. We will now show that the lattice field theory (5.1) with the constraint (5.2) is equivalent to a noncommutative U($`\stackrel{~}{p}_0`$) lattice gauge theory with fields obeying periodic boundary conditions and the reduced rank $`\stackrel{~}{p}_0`$ defined by (3.29). For this, we rewrite the lattice field theory (5.1) in terms of (finite dimensional) commutative operators. We write $`\widehat{U}_\mu `$ $`=`$ $`{\displaystyle \underset{x}{}}\widehat{\mathrm{\Delta }}(x)U_\mu (x)`$ (5.5) $`\widehat{\mathrm{\Omega }}_a`$ $`=`$ $`{\displaystyle \underset{x}{}}\widehat{\mathrm{\Delta }}(x)\mathrm{\Omega }_a(x),`$ (5.6) where the map $`\widehat{\mathrm{\Delta }}(x)`$ is defined by (4.18) and the operator $`\widehat{\mathrm{\Omega }}_a`$ is given explicitly by $$\widehat{\mathrm{\Omega }}_a=11\mathrm{\Gamma }_a.$$ (5.7) The action (5.1) can then be written as $$S=\frac{1}{g^2}\underset{\mu \nu }{}\text{Tr }\mathrm{tr}_{(p)}\left[\widehat{U}_\mu \left(\widehat{D}_\mu \widehat{U}_\nu \widehat{D}_\mu ^{}\right)\left(\widehat{D}_\nu \widehat{U}_\mu ^{}\widehat{D}_\nu ^{}\right)\widehat{U}_\nu ^{}\right],$$ (5.8) where $`\widehat{D}_\mu `$ are the lattice shift operators (4.7), and the constraint (5.2) becomes $$\left(\widehat{D}_\nu \right)^{\frac{1}{ϵ}\mathrm{\Sigma }_{\nu a}}\widehat{U}_\mu \left(\widehat{D}_\nu ^{}\right)^{\frac{1}{ϵ}\mathrm{\Sigma }_{\nu a}}=\widehat{\mathrm{\Omega }}_a\widehat{U}_\mu \widehat{\mathrm{\Omega }}_a^{}.$$ (5.9) Proceeding exactly as in Section 3.1, one finds that the general solution to the constraint (5.9) is given by $$\widehat{U}_\mu =\underset{\stackrel{}{m}}{}\left(\underset{a=1}{\overset{D}{}}\left(\widehat{Z}_a^{}\right)^{m_a}\right)\text{e}^{\pi i_{a<b}\mathrm{\Theta }_{ab}^{}m_am_b}u_\mu (\stackrel{}{m}),$$ (5.10) where $`u_\mu (\stackrel{}{m})`$ is a $`\stackrel{~}{p}_0\times \stackrel{~}{p}_0`$ matrix and $$\widehat{Z}_a^{}=\text{e}^{2\pi i\left(\mathrm{\Sigma }^{}_{}{}^{}1\right)_{a\mu }\widehat{x}_\mu }\underset{b=1}{\overset{D}{}}(\mathrm{\Gamma }_b)^{B_{ab}}.$$ (5.11) The $`D\times D`$ matrices $`\mathrm{\Sigma }^{}`$ and $`\mathrm{\Theta }^{}`$ are given by $`\mathrm{\Sigma }^{}`$ $`=`$ $`\mathrm{\Sigma }\stackrel{~}{P}`$ (5.12) $`\mathrm{\Theta }^{}`$ $`=`$ $`\stackrel{~}{P}^1B^{},`$ (5.13) where the integral matrices $`B`$ and $`\stackrel{~}{P}`$ are defined in Section 3.2. Because of their dependence on the twist eating solutions, the operators (5.11) obey the commutation relations $`\widehat{Z}_a^{}\widehat{Z}_b^{}`$ $`=`$ $`\text{e}^{2\pi i\mathrm{\Theta }_{ab}^{}}\widehat{Z}_b^{}\widehat{Z}_a^{}`$ $`\widehat{D}_\mu \widehat{Z}_a^{}\widehat{D}_\mu ^{}`$ $`=`$ $`\text{e}^{2\pi i\left(\mathrm{\Sigma }^{}_{}{}^{}1\right)_{a\mu }}\widehat{Z}_a^{}.`$ (5.14) The commuation relations (5.14) are the same as (4.13) and (4.14) with the replacements $`\widehat{Z}\widehat{Z}^{}`$, $`\mathrm{\Theta }\mathrm{\Theta }^{}`$ and $`\mathrm{\Sigma }\mathrm{\Sigma }^{}`$. The two matrices $`\mathrm{\Sigma }^{}`$ and $`\mathrm{\Theta }^{}`$ must satisfy the general constraint (4.16), which implies that $`M=\frac{1}{2ϵ}\mathrm{\Sigma }B^{}`$ must be an integer-valued matrix. The sum over $`\stackrel{}{m}`$ in (5.10) can then be taken over $`\text{}^D`$ modulo the periodicity $`m_am_a+\frac{1}{ϵ}\mathrm{\Sigma }_{\mu a}^{}`$ with $`\mu =1,\mathrm{},D`$. We now introduce a corresponding map $`\widehat{\mathrm{\Delta }}^{}(x^{})`$ analogously to (4.18) and decompose the operator $`\widehat{U}_\mu `$ in this new basis as $$\widehat{U}_\mu =\underset{x^{}}{}\widehat{\mathrm{\Delta }}^{}(x^{})U_\mu ^{}(x^{}).$$ (5.15) Substituting (5.15) into (5.8), we arrive at the action $$S=\frac{1}{g^2}\underset{x^{}}{}\underset{\mu \nu }{}\mathrm{tr}_{(\stackrel{~}{p}_0)}\left[U_\mu ^{}(x^{})^{}U_\nu ^{}(x^{}+ϵ\widehat{\mu })^{}U_\mu ^{}(x^{}+ϵ\widehat{\nu })^{}^{}U_\nu ^{}(x^{})^{}\right]$$ (5.16) where $$g^2=\frac{p}{\stackrel{~}{p}_0}g^2.$$ (5.17) This shows that an ordinary U($`p`$) lattice gauge theory with multi-valued gauge fields is equivalent to noncommutative U($`\stackrel{~}{p}_0`$) lattice gauge theory with single-valued gauge fields and deformation parameter matrix (5.13). We now show that for any noncommutative Yang-Mills theory with given $`\mathrm{\Theta }^{}`$ and $`\mathrm{\Sigma }^{}`$ in the continuum and with gauge group U($`\stackrel{~}{p}_0`$), we can construct a family of commutative U($`p`$) lattice gauge theories with multi-valued gauge fields whose sequence of Morita equivalent theories converges to the target noncommutative gauge theory in the continuum limit. Let us first consider the case where the noncommutativity parameter matrix $`\mathrm{\Theta }^{}`$ given in the continuum is rational-valued. By using the SL($`D,\text{}`$) transformation $`\mathrm{\Theta }^{}\mathrm{\Lambda }\mathrm{\Theta }^{}\mathrm{\Lambda }^{}`$, where $`\mathrm{\Lambda }\text{SL}(D,\text{})`$, we can rotate $`\mathrm{\Theta }^{}`$ into the canonical skew-diagonal form $$\mathrm{\Theta }^{}=\left(\begin{array}{ccccc}0& \vartheta _1^{}& & & \\ \vartheta _1^{}& 0& & & \\ & & \mathrm{}& & \\ & & & 0& \vartheta _d^{}\\ & & & \vartheta _d^{}& 0\end{array}\right),$$ (5.18) where the skew-eigenvalues $`\vartheta _i^{}`$ are also rational numbers. Let us denote them by $`\vartheta _i^{}=\frac{b_i}{\stackrel{~}{p}_i}`$, where $`b_i`$ and $`\stackrel{~}{p}_i`$ are relatively prime integers for each $`i=1,\mathrm{},d`$. Then there exist integers $`a_i`$ and $`\stackrel{~}{q}_i`$ which satisfy (3.38). We can define the lattice rank $`p`$ by (3.29) and the twist matrix $`Q`$ by (3.36). Since the matrix $`\stackrel{~}{P}`$ is invertible, we can approximate the period $`\mathrm{\Sigma }^{}`$ to arbitrary precision by the matrix $`\mathrm{\Sigma }\stackrel{~}{P}`$ in the continuum limit $`ϵ0`$, while keeping fixed the lattice period $`\frac{1}{ϵ}\mathrm{\Sigma }`$ to an integer-valued matrix. When $`\mathrm{\Theta }^{}`$ is irrational, one has to consider a rational-valued deformation parameter matrix $`\mathrm{\Theta }_ϵ^{}`$ for each lattice spacing $`ϵ`$, which converges to the given irrational numbers $`\mathrm{\Theta }_ϵ^{}\mathrm{\Theta }^{}`$ as $`ϵ0`$. This means that the corresponding integers $`b_{ϵ,i}`$ and $`\stackrel{~}{p}_{ϵ,i}`$ should go to infinity in the continuum limit. Thus, the corresponding $`p_ϵ`$’s given by (3.29) should also go to infinity. We conclude that, in order to define a continuum noncommutative gauge theory with irrational noncommutativity parameters, the rank $`p`$ of the gauge group of the approximating commutative lattice gauge theory should be sent to infinity as one takes the continuum limit. Note that if one takes the ’t Hooft (planar) limit by fixing $`\lambda =g^2p`$ as $`p\mathrm{}`$, then the coupling constant $`g^{}`$ of the Morita equivalent noncommutative U($`\stackrel{~}{p}_0`$) gauge theory is fixed according to (5.17). Therefore, in the ’t Hooft limit, the lattice spacing $`ϵ`$ should be taken to zero only after one takes the limit $`p\mathrm{}`$ and the limit should be accompanied by an appropriate renormalization of the ’t Hooft coupling constant $`\lambda `$. However, in this limit the dimensionful noncommutativity parameters $`\theta ^{}`$ inevitably diverge. On the other hand, in order to have finite $`\theta ^{}`$ one should take a “double scaling limit” by sending $`ϵ0`$ together with the $`p\mathrm{}`$ limit in a correlated way. In this limit, non-planar Feynman diagrams survive in addition to the planar diagrams. Thus, although we have to take the large $`p`$ limit in order to obtain irrational $`\mathrm{\Theta }^{}`$, this limit should not be confused with the ’t Hooft limit in general. ### 5.2 Twisted Eguchi-Kawai model We will now consider the construction of the previous Subsection in the special case where the period matrix is $`\mathrm{\Sigma }=ϵ11_D`$. A one-site U($`p`$) lattice gauge theory is just the Eguchi-Kawai model , and the fact that the boundary conditions are twisted as in (5.2) means that it is actually the twisted Eguchi-Kawai model . To see this explicitly, we reduce the action (5.1) to a single point $`x=0`$ by using the constraints (5.2) to get $$U_\mu (ϵ\delta _{\nu a}\widehat{\nu })=\mathrm{\Gamma }_aU_\mu (0)(\mathrm{\Gamma }_a)^{}.$$ (5.19) Substituting (5.19) into (5.1) at $`x=0`$ and using the commutation relations (3.4), we arrive at the action $$S=\frac{1}{g^2}\underset{\mu \nu }{}𝒵_{\mu \nu }\mathrm{tr}_{(p)}\left(V_\mu V_\nu V_\mu ^{}V_\nu ^{}\right)$$ (5.20) where $`V_\mu =U_\mu (0)\mathrm{\Gamma }_\mu `$, $`\mu =1,\mathrm{},D`$, are $`p\times p`$ unitary matrices. This is the action of the twisted Eguchi-Kawai model, where the phase factor $`𝒵_{\mu \nu }=\text{e}^{2\pi i\gamma _{\mu \nu }/p}`$ is called the “twist”. Thus, the recent proposal that the twisted large $`N`$ reduced model serves as a concrete definition of noncommutative Yang-Mills theory can be interpreted as the simplest example of Morita equivalence. The possibility of such an interpretation has also been suggested in Ref. . Conversely, one can derive the explicit map from matrices in the twisted Eguchi-Kawai model to U($`\stackrel{~}{p}_0`$) gauge fields on the noncommutative torus using the present formalism. For this, we decompose the matrices $`V_\mu `$ of the action (5.20) using the SU($`p`$) twist eaters $`\mathrm{\Gamma }_\mu `$ as $$V_\mu =U_\mu \mathrm{\Gamma }_\mu .$$ (5.21) We then use the twisted boundary conditions (5.2) to generate from the $`p\times p`$ unitary matrix $`U_\mu (0)=U_\mu `$ a commutative U$`(p)`$ lattice gauge field $`U_\mu (x)`$, and introduce the operator $`\widehat{U}_\mu `$ by (5.5). Then, using the expansion (5.15) and remembering that the original field theory is reduced to a point $`x=0`$, the U$`(\stackrel{~}{p}_0)`$ gauge field $`U_\mu ^{}(x^{})`$ on the noncommutative torus can be given as $`U_\mu ^{}(x^{})`$ $`=`$ $`\text{Tr }^{}\left(\widehat{\mathrm{\Delta }}^{}(x^{})\widehat{U}_\mu \right)`$ (5.22) $`=`$ $`{\displaystyle \frac{p}{\stackrel{~}{p}_0}}{\displaystyle \underset{\stackrel{}{m}}{}}\text{e}^{\pi i_{\mu <\nu }\mathrm{\Theta }_{\mu \nu }^{}m_\mu m_\nu }\text{e}^{2\pi i\left(\mathrm{\Sigma }^{}_{}{}^{}1\right)_{\mu \nu }m_\mu x_\nu ^{}}`$ $`\times \mathrm{tr}_{(\stackrel{~}{p}_1\mathrm{}\stackrel{~}{p}_d)}\left[V_\mu (\mathrm{\Gamma }_\mu )^{}{\displaystyle \underset{\nu =1}{\overset{D}{}}}\left({\displaystyle \underset{\rho =1}{\overset{D}{}}}(\mathrm{\Gamma }_\rho )^{B_{\nu \rho }}\right)^{m_\nu }\right].`$ Eq. (5.22) gives the mapping between the $`p\times p`$ unitary matrices $`V_\mu `$ of the twisted Eguchi-Kawai model and noncommutative U$`(\stackrel{~}{p}_0)`$ lattice gauge fields $`U_\mu ^{}(x^{})`$ with noncommutativity parameter matrix (5.13) and period matrix (5.12). This shows explicitly how, via Morita equivalence, large $`N`$ reduced gauge theories serve as a nonperturbative formulation of noncommutative Yang-Mills theory. As in the case of general periods $`\mathrm{\Sigma }`$ which we discussed in the previous Subsection, one can reproduce arbitrary dimensionless noncommutativity parameters $`\mathrm{\Theta }^{}`$ in the continuum limit. However, since $`\mathrm{\Sigma }=ϵ11_D`$ in the present case, the induced period matrix $`\mathrm{\Sigma }^{}`$ in (5.12) is not the most general one. For example, in $`D=2`$ dimensions, one can obtain only a square torus, i.e. $`\mathrm{\Sigma }^{}=\mathrm{}\mathrm{\hspace{0.17em}1}1_2`$, as the continuum spacetime. The rank $`p`$ of the gauge group of the reduced model must be sent to infinity even for rational $`\mathrm{\Theta }`$, unlike the case of general $`\mathrm{\Sigma }`$. We note also that the twisted Eguchi-Kawai model was originally proposed as a model which reproduces large $`N`$ gauge theory in the ’t Hooft limit (or the planar limit). In the present interpretation, however, we have to take a different large $`p`$ limit, as we discussed in the previous subsection, whereby there are non-planar contributions. Therefore, although the matrix model is the same, the limiting procedure that one should take is different in the two interpretations. Indeed, the scaling of Wilson loops in such a non-planar large $`p`$ limit has been observed numerically for the two-dimensional Eguchi-Kawai model in Ref. . Let us also comment on the compatibility of finite noncommutativity parameters and finite spacetime volume in the continuum limit. In Ref. , a particular choice of twist which corresponds to $`q_i=L^{d1}`$ in eq. (3.27) with $`p=L^d`$ was taken, where $`D=2d`$ is the dimension of spacetime. In that case, the integers $`\stackrel{~}{p}_i`$ and $`\stackrel{~}{q}_i`$ defined in (3.28) are given by $`\stackrel{~}{p}_i=L`$ and $`\stackrel{~}{q}_i=1`$. One can then take $`a_i=0`$ and $`b_i=1`$ to satisfy (3.38). Therefore, the volume of the torus is $`(ϵL)^D`$ and the dimensionful noncommutativity parameters are given by $`\theta _{\mu \nu }^{}Lϵ^2`$. As a result, in order to keep the noncommutativity parameters finite in the continuum limit, one has to send the volume of the discrete periodic lattice to infinity.<sup>4</sup><sup>4</sup>4Indeed, in the two-dimensional Eguchi-Kawai model studied in Ref. , the scaling of Wilson loops was observed in the limit of large $`L`$ while keeping the parameter $`Lϵ^2`$ fixed. In Ref. , this restriction was removed by considering a restricted twisted Eguchi-Kawai model obtained via the introduction of a quotient condition on the matrices which is the discrete analog of that which appears in the Connes-Douglas-Schwarz formalism . However, the present general construction suggests a much simpler way to avoid this constraint, namely by choosing the twist as described below eq. (5.18). We close this Section by remarking that the above derivations are independent of any representation of the operators $`\widehat{Z}_a^{}`$ and $`\widehat{D}_\mu `$. This refines the approach of whereby a specific finite dimensional representation of these operators was used to interpret the twisted Eguchi-Kawai model as noncommutative Yang-Mills theory. All we have used in the above construction is the equivalence of commutative and noncommutative gauge theories on the lattice at the field theoretical level, without ever having to specify the representation of the operators used. To this end, we note that the above construction could have been obtained directly from the action (4.29) by using the representation of the operators in terms of the $`p\times p`$ matrices $`\widehat{Z}_a^{}`$ $`=`$ $`{\displaystyle \underset{b=1}{\overset{D}{}}}(\mathrm{\Gamma }_b)^{B_{ab}},`$ (5.23) $`\widehat{D}_\mu `$ $`=`$ $`\mathrm{\Gamma }_\mu .`$ (5.24) These matrices clearly satisfy the commutation relations (5.14), and thereby constitute a finite dimensional representation of the noncommutative torus. ## 6 Coupling to fundamental matter fields In this last Section we will consider noncommutative gauge theory coupled to matter fields in the fundamental representation of the gauge group. One advantage of the lattice formulation is that it allows a hopping parameter (or large mass) expansion of this theory which will enable us to clarify various aspects of the Wilson loops in noncommutative gauge theories that were constructed in Section 2.3. A peculiar property of these Wilson loops, in contrast to their commutative counterparts, is that star-gauge invariance does not require that the loops be closed, but it does require that the separation between their two endpoints is proportional to the total momentum of the loop. In the following we will clarify the physical meaning of these observables. This analysis will also demonstrate explicitly how these observables reduce smoothly to ordinary closed Wilson loops in the commutative limit. For this, we will first show that the effective action for the gauge field induced by integrating over the matter fields can be written as a sum over all star-gauge invariant observables associated with closed loops. We will then construct star-gauge invariant local operators out of the matter fields and evaluate their correlation functions, thereby demonstrating that the star-gauge invariant observables associated with open loops appear very naturally in this way. We will also show that Morita equivalence holds in noncommutative Yang-Mills theory when it is coupled to fundamental matter fields. This Morita equivalence clarifies the interpretation of the perturbative beta-function which was calculated in Ref. . It also allows us to nonperturbatively define noncommutative Yang-Mills theory with fundamental matter fields, as we did in the matter-free case in Section 5.1. As a special case, we obtain the twisted Eguchi-Kawai model with fundamental matter fields as constructed in Ref. . The main idea of these constructions is to introduce a number of flavours $`N_f`$ which is equal to (or, more generally, a multiple of) the number of colours, i.e. the rank of the gauge group. This enables us to impose a boundary condition on the matter fields using a rotation in flavour space to mimick the boundary condition for adjoint representation fields. For simplicity, we will consider only complex scalar matter fields. Fermionic matter fields can be treated in an analogous way. ### 6.1 Properties of star-gauge invariant observables In this Subsection we will consider a discrete noncommutative torus and periodic boundary conditions on all fields. We introduce the operator $$\widehat{\varphi }=\underset{x}{}\widehat{\mathrm{\Delta }}(x)\varphi (x),$$ (6.1) where $`\varphi (x)`$ is a complex scalar field in the fundamental representation of the gauge group U($`p`$), and as always the lattice operator $`\widehat{\mathrm{\Delta }}(x)`$ is defined in (4.18). The action for the matter field can be written as $$S_{\mathrm{matter}}=\kappa \underset{\mu }{}\left\{\text{Tr }\left(\widehat{\varphi }^{}\widehat{U}_\mu \widehat{D}_\mu \widehat{\varphi }\widehat{D}_\mu ^{}\right)+\text{c.c.}\right\}+\text{Tr }\left(\widehat{\varphi }^{}\widehat{\varphi }\right),$$ (6.2) and it is invariant under the transformation $`\widehat{\varphi }`$ $``$ $`\widehat{g}\widehat{\varphi },`$ $`\widehat{\varphi }^{}`$ $``$ $`\widehat{\varphi }^{}\widehat{g}^{},`$ $`\widehat{U}_\mu `$ $``$ $`\widehat{g}\widehat{U}_\mu \widehat{D}_\mu \widehat{g}^{}\widehat{D}_\mu ^{}.`$ (6.3) In terms of lattice fields, the action can be written as $$S_{\mathrm{matter}}=\kappa \left\{\underset{x,\mu }{}\varphi (x)^{}U_\mu (x)\varphi (x+ϵ\widehat{\mu })+\text{c.c.}\right\}+\underset{x}{}\varphi (x)^{}\varphi (x),$$ (6.4) and it is invariant under the star-gauge transformation $`\varphi (x)`$ $``$ $`g(x)\varphi (x),`$ $`\varphi (x)^{}`$ $``$ $`\varphi (x)^{}g(x)^{},`$ $`U_\mu (x)`$ $``$ $`g(x)U_\mu (x)g(x+ϵ\widehat{\mu })^{}.`$ (6.5) When we integrate over the matter field $`\varphi (x)`$, we make an expansion with respect to the hopping parameter $`\kappa `$. This will require the computation of various correlation functions of the scalar fields, which can be calculated by using the formulae $`{\displaystyle \underset{x,y}{}}\varphi (x)^{}F(x)\varphi (x+v)\varphi (y)^{}_{\kappa =0}G(y)\varphi (y+u)`$ $`={\displaystyle \underset{x}{}}\varphi (x)^{}F(x)G(x+v)\varphi (x+v+u),`$ (6.6) $$\underset{x}{}\varphi (x)^{}F(x)\varphi (x+v)^{}_{\kappa =0}=\delta _{v,0}\underset{x}{}F(x).$$ (6.7) Throughout this section, the brackets $`\mathrm{}`$ denote the vacuum expectation value for fixed gauge background, i.e. we integrate over the matter fields only. The suffix $`\kappa =0`$ means that the hopping parameter is set to zero in the matter field action (6.4). Star-gauge invariance requires that the lattice fields $`F(x)`$ and $`G(x)`$ transform as $`F(x)`$ $``$ $`g(x)F(x)g(x+v)^{}`$ $`G(y)`$ $``$ $`g(y)G(y)g(y+u)^{}`$ (6.8) under the star-gauge transformation (6.5). Let us first consider the effective action $`\mathrm{\Gamma }_{\mathrm{eff}}[U]`$ for the gauge field $`U_\mu (x)`$ induced by the integration over $`\varphi (x)`$. Using the hopping parameter expansion, it can be given as $$\mathrm{\Gamma }_{\mathrm{eff}}[U]=\mathrm{ln}\left[\underset{n=0}{\overset{\mathrm{}}{}}\frac{\kappa ^n}{n!}\left(\underset{x,\mu }{}\varphi (x)^{}U_\mu (x)\varphi (x+ϵ\widehat{\mu })+\text{c.c.}\right)^n_{\kappa =0}\right].$$ (6.9) Integrating over the matter field using Wick’s theorem and the formulae (6.6) and (6.7), we obtain $$\mathrm{\Gamma }_{\mathrm{eff}}[U]=\underset{C}{}\frac{\kappa ^{L(C)}}{L(C)}\underset{x}{}\mathrm{tr}_{(p)}\left(𝒰(x;C)\right),$$ (6.10) where $`_C`$ is the sum over all closed loops (with and without self-intersection) on the lattice, and $`L(C)`$ denotes the number of links in the contour $`C`$. In this way, we encounter the star-gauge invariant observables associated with closed loops. Let us now consider star-gauge invariant observables involving matter fields such as $$G_1[f]=\underset{x}{}\varphi (x)^{}\varphi (x)f(x),$$ (6.11) which is star-gauge invariant for arbitrary functions $`f(x)`$ on the lattice. The lattice field $`f(x)`$ can be regarded as the wavefunction of the composite operator $`\varphi (x)^{}\varphi (x)`$. We first perform a Fourier transformation of $`f(x)`$ and express the observable (6.11) as $$G_1[f]=\underset{\stackrel{}{k}}{}\stackrel{~}{f}(k)\underset{x}{}\varphi (x)^{}\text{e}^{ik_\mu x_\mu }\mathrm{\hspace{0.17em}1}1_p\varphi (x+v),$$ (6.12) where $$v_\mu =\theta _{\mu \nu }k_\nu .$$ (6.13) Here we have used the fact that the plane wave $`\text{e}^{ik_\mu x_\mu }\mathrm{\hspace{0.17em}1}1_p`$ acts as a translation operator in the noncommutative field theory as explained in (2.25)–(2.27). Integrating over the matter field using the hopping parameter expansion, Wick’s theorem, and the formulae (6.6) and (6.7), we arrive at $$G_1[f]=\underset{\stackrel{}{k}}{}\stackrel{~}{f}(k)\underset{C_v}{}\kappa ^{L(C_v)}\underset{x}{}\mathrm{tr}_{(p)}\left(𝒰(x;C_v)\text{e}^{ik_\mu x_\mu }\mathrm{\hspace{0.17em}1}1_p\right),$$ (6.14) where $`_{C_v}`$ denotes the sum over all loops on the lattice starting from the origin and ending at the lattice point (6.13). Thus we find the star-gauge invariant observables encountered in Section 2.3 associated with open loops. In the commutative case $`\theta _{\mu \nu }=0`$, the separation vector (6.13) of the loops vanishes independently of their momenta $`k_\mu `$. Then the sum over loops in (6.14) contains only closed loops. The sum over $`\stackrel{}{k}`$ can therefore be done explicitly reproducing the wavefunction $`f(x)`$. In the noncommutative case, however, the sum over loops depends on $`k_\mu `$ and we cannot reverse the order of the first two sums in (6.14). The separation vector $`v_\mu `$ of the two ends of the loop $`C_v`$ grows with its momentum $`k_\mu `$. This is a characteristic phenomenon in noncommutative field theories, and it is another signal of the UV/IR mixing. If one would like to have a higher resolution in one direction, say in the $`\mu =1`$ direction, by increasing the momentum $`k_1`$, then the object will extend in the other directions proportionally to $`\theta _{\mu 1}k_1`$. Thus the size of the object depends on its momentum. On the other hand, if one considers the case where $`\stackrel{~}{f}(k)`$ is supported on finite momenta $`k_\mu `$, then one can take the $`\theta _{\mu \nu }0`$ limit smoothly, reproducing the commutative case. Let us next consider a two-point function of the composite operators which is defined as $$G_2[f,g]=\underset{x}{}\varphi (x)^{}\varphi (x)f(x)\underset{y}{}\varphi (y)^{}\varphi (y)g(y).$$ (6.15) The connected part of this correlation function can be readily computed as before to give $`G_2^{(\mathrm{conn})}[f,g]`$ $`\stackrel{\mathrm{def}}{=}`$ $`G_2[f,g]G_1[f]G_1[g]`$ (6.16) $`=`$ $`{\displaystyle \underset{\stackrel{}{k},\stackrel{}{p}}{}}\stackrel{~}{f}(k)\stackrel{~}{g}(p){\displaystyle \underset{C_u,C_v^{}}{}}\kappa ^{L(C_u)+L(C_v^{})}`$ $`\times {\displaystyle \underset{x,y}{}}\mathrm{tr}_{(p)}(𝒰(x;C_u)\text{e}^{ip_\nu y_\nu }\mathrm{\hspace{0.17em}1}1_p𝒰(y;C_v^{})\text{e}^{ik_\mu x_\mu }\mathrm{\hspace{0.17em}1}1_p),`$ where the double sum over contours in (6.16) is restricted in such a way that the parallel transport operator $`𝒰(x;C_u)`$ goes from $`x`$ to $`y+u`$ and $`𝒰(y;C_v^{})`$ from $`y`$ to $`x+v`$. The separation vectors $`v_\mu `$ and $`u_\mu `$ are related to the momenta $`k_\mu `$ and $`p_\mu `$ respectively by $`u_\mu =\theta _{\mu \nu }p_\nu `$ and $`v_\mu =\theta _{\mu \nu }k_\nu `$. As with the one-point function (6.14), the sums over $`\stackrel{}{k}`$ and $`\stackrel{}{p}`$ can be done explicitly in the commutative limit, reproducing the wavefunctions $`f(x)`$ and $`g(y)`$. In the noncommutative case, the sums over contours in (6.16) depend on the momenta $`\stackrel{}{k}`$ and $`\stackrel{}{p}`$ and one cannot interchange the momentum and loop sums. Note that the two parallel transport operators in (6.16) do not form a closed loop. The difference between the endpoint of one and the starting point of the other are given by $`u_\mu `$ and $`v_\mu `$. If one increases the momenta $`k_\mu `$ or $`p_\mu `$, then the sizes $`u_\mu `$ or $`v_\mu `$ also increase. The considerations in this section illustrate that the star-gauge invariant observables constructed in Section 2.3 do indeed play the same fundamental role as ordinary Wilson loops do in commutative gauge theories. As the knowledge of all Wilson loop correlators would give all the information about the quantum gauge theory coupled to matter fields, so do the star-gauge invariant observables in the noncommutative case. In particular, this result implies that the issue of whether or not a noncommutative gauge theory is confining can be addressed by searching for the area law behaviour of a star-gauge invariant observable associated with a closed loop. We have also seen explicitly how these observables reduce smoothly to ordinary Wilson loops in the commutative limit for fixed gauge background. ### 6.2 Morita equivalence with fundamental matter fields In this section, we will describe Morita equivalence of noncommutative Yang-Mills theories coupled to fundamental matter fields. We consider the setup of Section 5.1. We start with commutative U($`p`$) gauge theory and introduce $`N_f`$ flavours of matter fields in the fundamental representation of the gauge group. We represent it as a matrix $`\mathrm{\Phi }(x)_{iJ}`$, where $`i=1,\mathrm{},p`$ is the colour index and $`J=1,\mathrm{},N_f`$ is the flavour index. The spacetime is discretized as $`x_\mu ϵ\text{}`$. For the present purposes, it is essential to take $`N_f`$ to be an integer multiple of $`p`$. In what follows, we will assume that $`N_f=p`$ for simplicity. Then, $`\mathrm{\Phi }(x)_{iJ}`$ becomes a $`p\times p`$ complex-valued matrix, i.e. an element of $`\mathrm{gl}(p,\text{})`$. The action for the gauge field is given by (5.1). The action for the matter field is $$S_{\mathrm{matter}}=\kappa \left\{\underset{x,\mu }{}\mathrm{tr}_{(p)}\left(\mathrm{\Phi }(x)^{}U_\mu (x)\mathrm{\Phi }(x+ϵ\widehat{\mu })\right)+\text{c.c.}\right\}+\underset{x}{}\mathrm{tr}_{(p)}\left(\mathrm{\Phi }(x)^{}\mathrm{\Phi }(x)\right).$$ (6.17) This matter-coupled gauge theory is invariant under the gauge transformation $`U_\mu (x)`$ $``$ $`g(x)U_\mu (x)g(x+ϵ\widehat{\mu })`$ $`\mathrm{\Phi }(x)`$ $``$ $`g(x)\mathrm{\Phi }(x)`$ $`\mathrm{\Phi }(x)^{}`$ $``$ $`\mathrm{\Phi }(x)^{}g(x)^{}.`$ (6.18) The action (6.17) also possesses the global U($`N_f`$) flavour symmetry $$\mathrm{\Phi }(x)\mathrm{\Phi }(x)g^{};\mathrm{\Phi }(x)^{}g^{}\mathrm{\Phi }(x)^{},$$ (6.19) where $`g^{}\text{U}(N_f)`$. We impose twisted boundary conditions on the fields $`\mathrm{\Phi }(x)`$ and $`U_\mu (x)`$ given by $`U_\mu (x+\mathrm{\Sigma }_{\nu a}\widehat{\nu })`$ $`=`$ $`\mathrm{\Gamma }_aU_\mu (x)(\mathrm{\Gamma }_a)^{}`$ $`\mathrm{\Phi }(x+\mathrm{\Sigma }_{\nu a}\widehat{\nu })`$ $`=`$ $`\mathrm{\Gamma }_a\mathrm{\Phi }(x)(\mathrm{\Gamma }_a)^{},`$ (6.20) where $`\mathrm{\Sigma }`$ is the period matrix and $`\mathrm{\Gamma }_a`$ are the twist-eaters satisfying the commutation relations (3.4). Note that in (6.20), the $`\mathrm{\Gamma }_a`$ acting on the left of $`\mathrm{\Phi }(x)`$ represent a (global) gauge transformation, whereas the $`(\mathrm{\Gamma }_a)^{}`$ acting on the right represent a rotation in flavour space. This is the trick to introducing the fundamental representation of the gauge group in Morita equivalence transformations. Even though the matter fields define fundamental sections of the given gauge bundle over the noncommutative torus (see (6.18)), we can exploit the global SU($`N_f`$) flavour symmetry to mimick the boundary conditions for adjoint representation fields. Such an idea first appeared in Ref. in the context of supersymmetric field theories. Rewriting the fields in terms of operators and solving the constraints (6.20) as we did in Section 5.1, we find that the resulting field theory is noncommutative U($`\stackrel{~}{p}_0`$) lattice gauge theory, on a lattice with period matrix $`\mathrm{\Sigma }^{}=\mathrm{\Sigma }\stackrel{~}{P}`$ and with dimensionless noncommutativity parameters given by $`\mathrm{\Theta }^{}=\stackrel{~}{P}^1B^{}`$, coupled to $`\stackrel{~}{p}_0`$ flavours of matter fields in the fundamental representation of U($`\stackrel{~}{p}_0`$). When $`N_f=n_fp`$, we obtain $`n_f\stackrel{~}{p}_0`$ flavours in the Morita equivalent noncommutative field theory. The one-loop beta-function of this theory for $`\stackrel{~}{p}_0=1`$ has been calculated in Ref. . The contribution of fundamental matter fields to the one-loop beta-function of noncommutative U($`\stackrel{~}{p}_0`$) Yang-Mills theory with $`n_f`$ flavours coincides up to a finite rescaling of the Yang-Mills coupling constant with that in commutative U($`p`$) Yang-Mills theory with $`N_f=n_fp`$ flavours, since one-loop Feynman diagrams involving fundamental matter fields are always planar. This coincidence can now be understood as a manifestation at the one-loop level of the Morita equivalence we have found in this Section. As we did in Section 5.2, we can set $`\mathrm{\Sigma }=ϵ11_D`$ to obtain the twisted Eguchi-Kawai model coupled to matter fields. This is the model which was introduced in Ref. as a model which reproduces large $`N`$ gauge theory with $`N_f`$ flavours of matter in the Veneziano limit $`N_fN\mathrm{}`$. Here we have found that the same model with $`N_f=n_fN`$ can be interepreted as noncommutative U($`\stackrel{~}{p}_0`$) gauge theory with periodic fields including $`n_f\stackrel{~}{p}_0`$ flavours of matter fields in the fundamental representation of the gauge group U($`\stackrel{~}{p}_0`$). ### Acknowledgements The work of J.A. and Y.M. is supported in part by MaPhySto founded by the Danish National Research Foundation. Y.M. is sponsored in part by the Danish National Bank. J.N. is supported by the Japan Society for the Promotion of Science as a Research Fellow Abroad. The work of R.J.S. is supported in part by the Danish Natural Science Research Council.
warning/0004/hep-ph0004068.html	ar5iv	text	# 1 Introduction ## 1 Introduction The experimental discovery of the Higgs boson is crucial for the understanding of the mechanism of electroweak symmetry breaking. The search for Higgs particles is one of the main goals for the LEP2 and Tevatron experiments and is one of the major motivations for the future Large Hadron Collider (LHC) and Linear $`e^+e^{}`$ Collider (LC). Once the Higgs boson is discovered, it will be of primary importance to determine in a model independent way its tree-level and one-loop induced couplings, spin, parity, $`CP`$-nature, and its total width. In this respect the $`\gamma \gamma `$ Compton Collider option of the LC offers a unique opportunity to produce both Standard Model (SM) Higgs boson and neutral Higgs states $`h`$, $`H`$, $`A`$ of the Minimal Supersymmetric Standard Model (MSSM) or general two Higgs Doublet Model (2HDM) as $`s`$-channel resonance decaying into $`b\overline{b}`$, $`WW^{}`$, $`ZZ`$ or $`t\overline{t}`$: $$\gamma \gamma h^0,H^0,A^0b\overline{b},WW^{},ZZ,t\overline{t}.$$ The ability to control the polarizations of back-scattered photons provides a powerful means for exploring the $`CP`$ properties of any single neutral Higgs boson that can be produced with reasonable rate at the Photon Linear Collider . $`CP`$-even Higgs $`0^{++}`$ bosons $`h^0`$, $`H^0`$ couple to the combination $$\stackrel{}{\epsilon _1}\stackrel{}{\epsilon _2}=\frac{1}{2}(1+\lambda _1\lambda _2),$$ (1) while a $`CP`$-odd $`0^+`$ Higgs boson $`A^0`$ couples to $$[\stackrel{}{\epsilon _1}\times \stackrel{}{\epsilon _2}]\stackrel{}{k_\gamma }=\frac{\omega _\gamma }{2}i\lambda _1(1+\lambda _1\lambda _2),$$ (2) where $`\stackrel{}{\epsilon _i}`$ and $`\lambda _i`$ are photon polarization vectors and helicities. The first of these structures couples to linearly polarized photons with the maximal strength if the polarizations are parallel, the letter if the polarizations are perpendicular. Moreover, if the Higgs boson is a mixture of $`CP`$-even and $`CP`$-odd states, as can occur e.g. in a general 2HDM with $`CP`$-violating neutral sector, the interference of these two terms gives rise to $`CP`$-violating asymmetries . Since MSSM Higgs particles $`h^0`$, $`H^0`$, $`A^0`$ decay predominantly into $`b\overline{b}`$ or $`t\overline{t}`$ quark pairs depending on the mass of the Higgs boson, the heavy quark pair background in $`\gamma \gamma `$ collisions has been studied in great detail. One-loop QCD corrections were calculated for the photon helicity states corresponding to projection of total angular momentum on beam axes $`J_z=0`$ and $`J_z=\pm 2`$ . Virtual one-loop QCD corrections for $`J_z=0`$ were found to be especially large due to the double-logarithmic enhancement factor, so that the corrections are comparable or even larger than the Born contribution for the two-jet final topologies . In order to solve this theoretical problem leading QCD corrections for $`J_z=0`$ have been calculated at the two-loop level and recently these leading double-logarithmic QCD corrections were resummed to all orders . The account of non-Sudakov form factor to higher orders makes the $`J_z=0`$ cross section well defined and positive definite in all regions of the phase space . All these studies of the influence of QCD corrections on heavy quark production in $`\gamma \gamma `$ collisions were concentrated on circularly polarized initial photons. However, for the direct measurements of the parity of states of Higgs bosons, (1-2), linear polarization of photon beams is needed . In the present paper we consider the QCD corrections to heavy quark-antiquark pair production in photon-photon collision for the general case of initial photon polarizations. We mainly concentrate on QCD corrections for linearly polarized photon-photon collisions. The production cross sections and spin asymmetries for $`t\overline{t}`$-pair production are calculated for linearly polarized photon collisions. The measurement of spin asymmetries is necessary to determine the $`CP`$ parity of the Higgs boson. In the scattering of linearly polarized photon off circularly polarized ones at one-loop level azimuthal asymmetries arise in the production of heavy fermion pairs. This is a pure quantum effect which does not exist at the Born level. These type of asymmetries are suppressed by factor $`m_Q^2/s`$, $`s`$ is the c.m.s energy of colliding photons, and are sizeable only for $`t\overline{t}`$-pair production. The paper is organized as follows. In the next Section we recall the basic definitions and consider the Born cross section of heavy quark-antiquark pair production in polarized photon collisions. Calculations of virtual corrections are presented in Section 3. The real gluon emission part is discussed in Section 4. The numerical results for top-antitop pair production cross sections and expected asymmetries for linearly polarized photon beams are discussed in Section 5. ## 2 Born Cross Section The cross section of heavy quark-antiquark pair production in polarized photon-photon collision $$\gamma (p_1)+\gamma (p_2)Q(p_3)+\overline{Q}(p_4)$$ (3) can be written in the most general form using the Stokes parameters which describe the polarizations of initial photons. The covariant density matrix of polarized photon with arbitrary polarization can be written in the following form $`\rho _{\mu \nu }^{(1,2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(e_\mu ^xe_\nu ^x+e_\mu ^ye_\nu ^y\right)\pm {\displaystyle \frac{\xi _1^{(1,2)}}{2}}\left(e_\mu ^xe_\nu ^y+e_\mu ^ye_\nu ^x\right)`$ (4) $`{\displaystyle \frac{\mathrm{i}\xi _2^{(1,2)}}{2}}\left(e_\mu ^xe_\nu ^ye_\mu ^ye_\nu ^x\right)+{\displaystyle \frac{\xi _3^{(1,2)}}{2}}\left(e_\mu ^xe_\nu ^xe_\mu ^ye_\nu ^y\right).`$ Here $`\xi _i^{(1,2)}`$ are three Stokes parameters describing polarization of the photon with momentum $`p_{1,2}`$ and $`e^x`$ and $`e^y`$ denote ort vectors in $`x`$ and $`y`$ directions. The momenta of the particles involved in the reaction in the c.m.s. of the initial photons are given by $`p_1=E(1;0,0,1),p_2=E(1;0,0,1),`$ (5) $`p_3=E(1;\beta \mathrm{sin}\theta \mathrm{cos}\varphi ,\beta \mathrm{sin}\theta \mathrm{sin}\varphi ,\beta \mathrm{cos}\theta ),p_4=E(1;\beta \mathrm{sin}\theta \mathrm{cos}\varphi ,\beta \mathrm{sin}\theta \mathrm{sin}\varphi ,\beta \mathrm{cos}\theta ),`$ where $`E=\sqrt{s}/2`$ is the photon beam energy and $`\beta =\sqrt{14m_Q^2/s}`$ is the quark (antiquark) velocity. With this definitions, the Born cross section of heavy quark pair production in photon-photon collisions has the form $`{\displaystyle \frac{\mathrm{d}\sigma (\gamma \gamma b\overline{b})}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{6\pi \alpha ^2e_Q^4}{s^2}}\{({\displaystyle \frac{t_1}{u_1}}+{\displaystyle \frac{u_1}{t_1}})(1+\stackrel{~}{\xi }_2^{(1)}\stackrel{~}{\xi }_2^{(2)}(12{\displaystyle \frac{Y}{t_1u_1}}))+4(1+\xi _3^{(1)}+\stackrel{~}{\xi }_3^{(2)}){\displaystyle \frac{m_Q^2sY}{t_1^2u_1^2}}`$ (6) $`+2(\stackrel{~}{\xi }_1^{(1)}\stackrel{~}{\xi }_1^{(2)}\stackrel{~}{\xi }_3^{(1)}\stackrel{~}{\xi }_3^{(2)})(12{\displaystyle \frac{Y}{t_1u_1}})4\stackrel{~}{\xi }_3^{(1)}\stackrel{~}{\xi }_3^{(2)}{\displaystyle \frac{Y^2}{t_1^2u_1^2}}\},`$ where $$Y=tum_Q^4,t_1=tm_Q^2\text{and}u_1=um_Q^2,$$ and $`s`$, $`t`$ and $`u`$ are Mandelstam variables for the process under consideration, $$s=(p_1+p_2)^2,t=(p_1p_3)^2,u=(p_1p_4)^2.$$ In the above formulae instead of the original Stokes parameters their combinations are used. The expressions of the cross sections are shorter and more convenient for integration if one includes the dependence on the azimuthal angle in the Stokes parameters, $`\stackrel{~}{\xi }_i^{(1,2)}`$. They can be expressed by original Stokes parameters through the equations: $`\stackrel{~}{\xi }_1^{(1)}`$ $`=`$ $`\xi _1^{(1)}\mathrm{cos}(2\varphi )\xi _3^{(1)}\mathrm{sin}(2\varphi ),`$ $`\stackrel{~}{\xi }_1^{(1)}`$ $`=`$ $`\xi _2^{(1)},`$ $`\stackrel{~}{\xi }_3^{(1)}`$ $`=`$ $`\xi _1^{(1)}\mathrm{sin}(2\varphi )+\xi _3^{(1)}\mathrm{cos}(2\varphi ),`$ $`\stackrel{~}{\xi }_1^{(2)}`$ $`=`$ $`\xi _1^{(2)}\mathrm{cos}(2\varphi )+\xi _3^{(2)}\mathrm{sin}(2\varphi ),`$ $`\stackrel{~}{\xi }_1^{(2)}`$ $`=`$ $`\xi _2^{(2)},`$ $`\stackrel{~}{\xi }_3^{(2)}`$ $`=`$ $`\xi _1^{(2)}\mathrm{sin}(2\varphi )+\xi _3^{(2)}\mathrm{cos}(2\varphi ),`$ (7) here $`\varphi `$ is the azimuthal angle. Parameters $`\stackrel{~}{\xi }_i^{(1,2)}`$ describe photon polarization with respect to unit vectors $`\stackrel{~}{e^x}`$, $`\stackrel{~}{e^y}`$, where $`\stackrel{~}{e^x}`$ lies in the reaction plane and $`\stackrel{~}{e^y}`$ is orthogonal to the reaction plane. ## 3 Virtual Corrections The first order QCD corrections to the cross section are determined by the interference between the tree level and one-loop diagrams. At the one-loop level no three-gluon vertex enters and the calculation is analogous to calculations of QED corrections to Compton scattering for finite electron mass (see e.g. ). The calculations of virtual corrections were done by using the symbolic manipulation program FORM . To regularise the infrared singularities we introduced an infinitesimal mass of the gluon $`\lambda `$. In the basis of the Stokes parameters the one-loop corrections have the form $`{\displaystyle \frac{\mathrm{d}\sigma (\gamma \gamma Q\overline{Q})}{\mathrm{d}t}}=`$ $`_0+\stackrel{~}{\xi }_1^{(1)}\stackrel{~}{\xi }_1^{(2)}_{11}+\mathrm{i}(\stackrel{~}{\xi }_1^{(1)}\stackrel{~}{\xi }_2^{(2)}+\stackrel{~}{\xi }_2^{(1)}\stackrel{~}{\xi }_1^{(2)})_{12}+\stackrel{~}{\xi }_2^{(1)}\stackrel{~}{\xi }_2^{(2)}_{22}+\stackrel{~}{\xi }_3^{(1)}\stackrel{~}{\xi }_3^{(2)}_{33}`$ (8) $`+(\stackrel{~}{\xi }_3^{(1)}+\stackrel{~}{\xi }_3^{(2)})_{03}.`$ In addition to the Born level expression, Eq. 2, there is the new term proportional to the non-diagonal product $`(\xi _1^{(1)}\xi _2^{(2)}+\xi _2^{(1)}\xi _1^{(2)})`$, which corresponds to the scattering of linearly polarized photon on the circularly polarized one. The functions $`_i`$ can be expressed through scalar one-loop integrals $`B`$, $`C`$ and $`D`$. $`_0`$ $`=`$ $`{\displaystyle \frac{\alpha _S\alpha ^2e_Q^4}{2\pi s}}\beta \{2D(s,t)(2s{\displaystyle \frac{s^2t_1+4m_Q^4(s2t)}{t_1u_1}}+stss_1+8m_Q^4)`$ (9) $`4C_1(t)\left(2{\displaystyle \frac{3m_Q^4s^2m_Q^2s}{u_1}}+2m_Q^4{\displaystyle \frac{st_1}{t_1^2}}5m_Q^2s+t\right)`$ $`+2B(t)(4m_Q^4({\displaystyle \frac{s+4s_1}{st_1u_1}}{\displaystyle \frac{s_0}{t_1^3}})+4m_Q^2{\displaystyle \frac{3st2m_Q^2s_1}{st_1^2}}`$ $`({\displaystyle \frac{m_Q^2}{t}}{\displaystyle \frac{s}{t_1}})(12{\displaystyle \frac{t_2}{u_1}})2{\displaystyle \frac{ss_1}{u_1t_1}}+1)`$ $`+C_1(s)s_0{\displaystyle \frac{s^2+2m_Q^2s+2t_1u_1}{t_1u_1}}C(s)s{\displaystyle \frac{3s^22t_1u_18m_Q^4}{t_1u_1}}`$ $`+4B(s){\displaystyle \frac{sY}{s_0t_1u_1}}+2\mathrm{ln}\left({\displaystyle \frac{\lambda ^2}{m_Q^2}}\right)\left(4{\displaystyle \frac{m_Q^4s^2}{t_1^2u_1^2}}s{\displaystyle \frac{s+4m_Q^2}{t_1u_1}}+2\right)+16{\displaystyle \frac{s^2m_Q^4}{t_1^2u_1^2}}`$ $`+4{\displaystyle \frac{m_Q^4}{t_1u_1}}m_Q^2(t+u){\displaystyle \frac{2m_Q^4+t_1u_1}{tt_1uu_1}}4m_Q^2s{\displaystyle \frac{t_1^2+u_1^2}{t_1^2u_1^2}}3s{\displaystyle \frac{6m_Q^2+s}{t_1u_1}}`$ $`+6+(tu)\},`$ $`_{11}`$ $`=`$ $`{\displaystyle \frac{\alpha _S\alpha ^2e_Q^4}{2\pi s}}\beta \{2D(s,t)(2s{\displaystyle \frac{8m_Q^4+3s^27m_Q^2s}{u_1}}+{\displaystyle \frac{s^2}{Y}}(s(s+m_Q^27t)`$ (10) $`+2m_Q^2(3tm_Q^2)+{\displaystyle \frac{s^3}{u_1}})+st3m_Q^2s+8m_Q^4)`$ $`4C_1(t)\left(s{\displaystyle \frac{s(s_05t_1)u_1+4m_Q^2t_1u_1+s^3}{u_1Y}}+4{\displaystyle \frac{s^22m_Q^2s+2m_Q^4}{u_1}}5m_Q^2+t\right)`$ $`+2B(t)\left(8{\displaystyle \frac{m_Q^2t}{t_1^2}}5{\displaystyle \frac{m_Q^2s}{t_1u_1}}+{\displaystyle \frac{m_Q^2Y}{tt_1u_1}}2{\displaystyle \frac{s}{u_1}}+1\right)`$ $`+C_1(s)\left(8s^2{\displaystyle \frac{s2m_Q^2}{Y}}{\displaystyle \frac{s^5}{t_1u_1Y}}2s{\displaystyle \frac{3s^25m_Q^2s+4m_Q^4}{t_1u_1}}+2s_0\right)`$ $`C(s)s\left(4{\displaystyle \frac{s^2+6m_Q^43m_Q^2s}{t_1u_1}}2s{\displaystyle \frac{3s4m_Q^2}{Y}}+{\displaystyle \frac{s^4}{t_1u_1Y}}+2\right)`$ $`+2B(s)s{\displaystyle \frac{2t_1u_1s^2+2m_Q^2s}{t_1u_1s_0}}+4\mathrm{ln}\left({\displaystyle \frac{\lambda ^2}{m_Q^2}}\right)\left(12{\displaystyle \frac{m_Q^2s}{t_1u_1}}\right)`$ $`+{\displaystyle \frac{m_Q^4}{t_1u}}+{\displaystyle \frac{m_Q^4}{tu_1}}19{\displaystyle \frac{m_Q^2s}{t_1u_1}}+6+(tu)\},`$ $`_{12}`$ $`=`$ $`{\displaystyle \frac{\alpha _S\alpha ^2e_Q^4}{2\pi s}}\beta \{2D(s,t)m_Q^2s_0(2{\displaystyle \frac{ss_0t_1}{u_1Y}})`$ (11) $`4C_1(t)\left(s_0{\displaystyle \frac{ts2m_Q^2t_1}{Y}}{\displaystyle \frac{(s3m_Q^2)^2}{u_1}}m_Q^2{\displaystyle \frac{m_Q^2t_1ts}{t_1^2}}\right)`$ $`+2B(t)m_Q^2\left({\displaystyle \frac{(t_1+2m_Q^2)(5t_1s)}{t_1^3}}+{\displaystyle \frac{((u2t)t_1+2m_Q^4}{tt_1u_1}}\right)`$ $`+C_1(s)ss_0\left({\displaystyle \frac{s2m_Q^2}{t_1u_1}}{\displaystyle \frac{s_0}{Y}}\right)+C(s)\left(s_0^2s\left({\displaystyle \frac{1}{t_1u_1}}{\displaystyle \frac{1}{Y}}\right)8{\displaystyle \frac{m_Q^2Y}{t_1u_1}}\right)`$ $`Y{\displaystyle \frac{2(t+u)m_Q^2(tu)^2}{tt_1uu_1}}+(tu)\},`$ $`_{22}`$ $`=`$ $`{\displaystyle \frac{\alpha _S\alpha ^2e_Q^4}{2\pi s}}\beta \{2D(s,t)(2s{\displaystyle \frac{s^2t+m_Q^2s^2+4m_Q^4(ut)}{t_1u_1}}+s(stm_Q^2)8m_Q^4)`$ (12) $`4C_1(t)\left(2{\displaystyle \frac{s^2+m_Q^2s5m_Q^4}{u_1}}2{\displaystyle \frac{m_Q^4(t_1s)}{t_1^2}}t+s+5m_Q^2\right)`$ $`+2B(t)\left(2m_Q^2{\displaystyle \frac{4tu+3m_Q^2}{t_1u_1}}4m_Q^2{\displaystyle \frac{3stm_Q^2t_1}{t_1^3}}+3{\displaystyle \frac{m_Q^2t_1st}{tt_1}}2{\displaystyle \frac{m_Q^2u}{tu_1}}1\right)`$ $`C_1(s)\left({\displaystyle \frac{ss_0(s+2m_Q^2)}{t_1u_1}}+2s_0\right)C(s)\left(3s{\displaystyle \frac{s^28m_Q^4}{t_1u_1}}2s\right)`$ $`4B(s){\displaystyle \frac{sY}{s_0t_1u_1}}+2\mathrm{ln}\left({\displaystyle \frac{\lambda ^2}{m_Q^2}}\right)\left(s{\displaystyle \frac{s+4m_Q^2}{t_1u_1}}2{\displaystyle \frac{m_Q^2s^3}{t_1^2u_1^2}}2\right)`$ $`s{\displaystyle \frac{8m_Q^2s^218m_Q^2t_1u_13t_1su_1}{t_1^2u_1^2}}+m_Q^2(u+t){\displaystyle \frac{2m_Q^4+t_1u_1}{tt_1uu_1}}4m_Q^4t_1u_1`$ $`6+(tu)\},`$ $`_{33}`$ $`=`$ $`{\displaystyle \frac{\alpha _S\alpha ^2e_Q^4}{2\pi s}}\beta \{2D(s,t)({\displaystyle \frac{s^4t_1}{u_1Y}}+s^2{\displaystyle \frac{7ts_14m_Q^2t_1m_Q^2(s+t)}{Y}}`$ (13) $`2s{\displaystyle \frac{t_1(3s^2+m_Q^4)7m_Q^2t_1s_14m_Q^4s}{t_1u_1}}16{\displaystyle \frac{sm_Q^6}{t_1u_1}}+3sm_Q^2st8m_Q^4)`$ $`+4C_1(t)\left({\displaystyle \frac{s^3t_1}{u_1Y}}4{\displaystyle \frac{(sm_Q^2)^2+m_Q^4}{u_1}}+{\displaystyle \frac{5s^2t_1+4m_Q^2s(u_1+2s)}{Y}}+5m_Q^2t\right)`$ $`+2B(t)\left({\displaystyle \frac{(t+m_Q^2)(2t+m_Q^2)}{tu_1}}2m_Q^2{\displaystyle \frac{t5m_Q^2}{t_1^2}}+16m_Q^2\left({\displaystyle \frac{usm_Q^4}{st_1u_1}}+{\displaystyle \frac{m_Q^2t}{t_1^3}}{\displaystyle \frac{m_Q^4}{st_1^2}}\right)+3\right)`$ $`C_1(s)\left(2{\displaystyle \frac{ss_0s_1}{t_1u_1}}+{\displaystyle \frac{s^3(4Y+s^2)}{t_1u_1Y}}8{\displaystyle \frac{s^2(s2m_Q^2)}{Y}}2s_0\right)`$ $`C(s)s(2s{\displaystyle \frac{4m_Q^23s}{Y}}+{\displaystyle \frac{s^4}{t_1u_1Y}}+4s_1{\displaystyle \frac{s2m_Q^2}{t_1u_1}}+2)`$ $`2B(s)s{\displaystyle \frac{s^22sm_Q^22t_1u_1}{s_0u_1t_1}}4\mathrm{ln}\left({\displaystyle \frac{\lambda ^2}{m_Q^2}}\right)\left(2{\displaystyle \frac{m_Q^2sY}{t_1^2u_1^2}}1\right)`$ $`m_Q^2s{\displaystyle \frac{11t_1u_116sm_Q^2}{t_1^2u_1^2}}+m_Q^4{\displaystyle \frac{tu_1+t_1u}{tt_1uu_1}}+6+(tu)\},`$ $`_{03}`$ $`=`$ $`{\displaystyle \frac{\alpha _S\alpha ^2e_Q^4}{2\pi s}}\beta \{2D(s,t)({\displaystyle \frac{sts_0^2}{Y}}{\displaystyle \frac{s(s2m_Q^2)(st_18m_Q^2t)}{t_1u_1}}2m_Q^2s_0)`$ (14) $`C_1(t)\left(4s_0{\displaystyle \frac{st2t_1m_Q^2}{Y}}4s_1{\displaystyle \frac{s7m_Q^2}{u_1}}4m_Q^2{\displaystyle \frac{stt_1m_Q^2}{t_1^2}}\right)`$ $`2B(t)(18{\displaystyle \frac{m_Q^4}{t_1u_1}}16{\displaystyle \frac{m_Q^6}{st_1u_1}}m_Q^2{\displaystyle \frac{(s8m_Q^2)(st_1+2st2t_1m_Q^2)}{st_1^3}}`$ $`m_Q^2{\displaystyle \frac{t_1u_111tu_12tt_1t_1m_Q^2}{tt_1u_1}})C(s)({\displaystyle \frac{ss_0^2}{Y}}s{\displaystyle \frac{ss_08m_Q^2s1}{t_1u_1}}8m_Q^2)`$ $`+C_1(s)ss_0\left({\displaystyle \frac{s2m_Q^2}{t_1u_1}}{\displaystyle \frac{s_0}{Y}}\right)+8\mathrm{ln}\left({\displaystyle \frac{\lambda ^2}{m_Q^2}}\right){\displaystyle \frac{m_Q^2sY}{t_1^2u_1^2}}`$ $`2m_Q^4{\displaystyle \frac{8s^2+t_1u_1}{t_1^2u_1^2}}+m_Q^2(t+u){\displaystyle \frac{m_Q^4+t_1u_1}{tt_1uu_1}}+2m_Q^2s{\displaystyle \frac{t_1^2+u_1^2}{t_1^2u_1^2}}s{\displaystyle \frac{s15m_Q^2}{t_1u_1}}`$ $`+4+(tu)\}.`$ Here $`t_2=t+m_Q^2`$, $`s_0=s4m_Q^2`$, and $`s_1=sm_Q^2`$. The definitions and expressions for the scalar integrals are given in the Appendix A of . It is worth mentioning that all functions $`_i`$ are expressed through four- and three-point functions and ultraviolet finite combinations of two-point functions. ## 4 Real Gluon Emission The contribution of the real gluon emission to the total cross section is separated in two parts, soft gluon emission which cancels out the infrared divergences of virtual corrections and hard gluon emission. The cross section of the soft gluon emission can be reproduced in a factorized form as a product of Born level cross section and the infrared divergent factor: $`{\displaystyle \frac{d\sigma ^{soft}}{dt}}={\displaystyle \frac{d\sigma ^{tree}}{dt}}R,`$ (15) where $`R`$ $`=`$ $`{\displaystyle \frac{8\alpha _s}{3\pi }}\{[1+{\displaystyle \frac{1}{\beta }}(1{\displaystyle \frac{2m_Q^2}{s}})\mathrm{ln}{\displaystyle \frac{1+\beta }{1\beta }}]\mathrm{ln}{\displaystyle \frac{2k_c}{\lambda }}`$ (16) $`+{\displaystyle \frac{1}{2\beta }}\left(1{\displaystyle \frac{2m_Q^2}{s}}\right)\left[\mathrm{Sp}\left({\displaystyle \frac{2\beta }{1\beta }}\right)\mathrm{Sp}\left({\displaystyle \frac{2\beta }{1+\beta }}\right)\right]`$ $`+{\displaystyle \frac{1}{2\beta }}\mathrm{ln}{\displaystyle \frac{1+\beta }{1\beta }}\},`$ Here $`k_c`$ is the soft photon energy cut and the velocity $`\beta `$ is defined in Section 2. The dependence on the gluon mass $`\lambda `$ is exactly canceled while adding soft gluon emission part and virtual corrections for any sets of Stokes parameters. As can be seen from Eqs. (6) and (8) the term proportional to the product $`\xi _1^{(1)}\xi _2^{(2)}`$ appears only in the virtual part and therefore should be infrared finite, although one-loop scalar functions $`D(s,t)`$, $`D(s,u)`$ and $`C_1(s)`$ do contain infrared logarithms. One can easily check that infrared divergent contributions from these scalar functions cancel each other in the $`_{12}`$-term of Eq. 8. The hard gluon emission part is also calculated using FORM program. The expressions for the squared matrix elements are lengthy and we do not reproduce them in the paper. The integration over three-particle phase space is done by using Monte-Carlo integration routine VEGAS . After adding the soft and hard gluon emission cross sections the final numerical results do not depend on the imposed energy cut of emitted gluon, $`k_c`$. Special care is taken to handle sharp peaks of the cross section while the gluon is soft or is emitted along the quark or antiquark three momenta (see detailed discussion in ). These peaks correspond to the infrared and collinear singularities in the case of massless fermions and in our case become essential when mass of quark is small compared to the c.m.s. energies of photons. To arrange the infrared and collinear singularities along the integration axis we take the gluon energy, denominator of the quark propagator, polar angles of quark and antiquark and azimuthal angle of reaction as integration variables. Integration over azimuthal angle is necessary in the case of linearly polarized photons, as in this case the cross section does depend on the $`\varphi `$ angle. The infrared singularity lies on the axis of integration over gluon energy, whereas collinear singularities are located along the axis of quark propagator. There are also additional peaks when quark or antiquark are produced at zero angle with respect to the beam direction. Such singularities are treated by integration over the polar angles of quark and antiquark. ## 5 Results and Discussion As was mentioned above, the helicity cross sections $`\sigma (J_z=0)`$, $`\sigma (J_z=2)`$ for the heavy quark pair production in the circularly polarized photon-photon collisions were considered in . In this paper we mainly present numerical results for the production of heavy quark-antiquark pair for the linearly polarized photons. We consider two cases of linear polarizations of initial photons, when $`\mathrm{\Delta }\gamma =0`$ and $`\pi /2`$, where $`\mathrm{\Delta }\gamma `$ is the angle between the directions of polarization vectors of the photons. The $`\mathrm{\Delta }\gamma =0,\pi /2)`$ correspond to the collision of linearly polarized photons with parallel and perpendicular polarizations, respectively. For the measurements of the Higgs boson parity it is necessary to consider collisions of linearly polarized photons in order to measure the polarization asymmetry $$A=\frac{\sigma _{}\sigma _{}}{\sigma _{}+\sigma _{}}.$$ (17) In fact if inclusive two-jet final states are studied then after averaging over azimuthal angles and spins of the final particles only three independent cross sections remain for arbitrary polarization states of initial photons. These independent cross sections can be taken as $`\sigma _{tot}`$, $`\sigma (J_z=0)\sigma (J_z=2)`$ and $`\sigma _{}\sigma _{}`$ . For the study of the Higgs boson signal in photon-photon collisions it was essential that the background from $`b\overline{b}`$ quark production is suppressed by a factor of $`m_b^2/s`$ for $`J_z=0`$ at the Born level. However at the next-to-leading order the cross section of the $`b\overline{b}g`$ production for $`J_z=0`$ is not suppressed any more. Therefore experimental cuts selecting only two-jet final states were important to suppress the $`b\overline{b}g`$ background . In this Section we show, that the difference of $`\sigma _{}\sigma _{}`$ is suppressed by a factor of $`m_Q^2/s`$ even at the next-to-leading order. The cross section for heavy quark pair production can be cast in the form : $`\sigma _{\gamma \gamma Q\overline{Q}(g)}`$ $`=`$ $`{\displaystyle \frac{\alpha ^2Q^4N_c}{s}}\left(f_{\gamma \gamma }^{(0)}+{\displaystyle \frac{4}{3}}{\displaystyle \frac{\alpha _s}{\pi }}f_{\gamma \gamma }^{(1)}\right),`$ (18) where the functions $`f_,^{(0,1)}`$ depend on the dimensionless variable $`s/(4m_Q^2)`$ only. The numerical values of the functions $`f_,^{(0,1)}`$ are presented in the Figure 1 and Table I. Because of the Sommerfeld rescattering correction, the function $`f_{}^{(1)}`$ is nonzero at the threshold, as one can see from the table I. In the both cases, $`\mathrm{\Delta }\gamma =0`$ and $`\mathrm{\Delta }\gamma =\pi /2`$, the functions corresponding to QCD corrections are positive and rising at high energies. Taking the average of values of these two functions one obtains the corresponding function for the unpolarized cross section. Our results are in agreement with previous calculations of QCD corrections for heavy quark-antiquark production in unpolarized photon-photon collisions . As in the case of Born level functions $`f_{}^{(0)}`$ and $`f_{}^{(0)}`$, in the asymptotic regime the difference between $`f_{}^{(1)}`$ and $`f_{}^{(1)}`$ vanishes and each of the function tends to the unpolarized one, $`f_{unpol}^1`$. Such an asymptotic behavior of the corrections can be understood considering the helicity amplitudes for massless quarks. The difference of the cross sections with parallel and orthogonal polarized photons can be expressed via the interference term of the following helicity amplitudes $`\mathrm{\Delta }\sigma (\gamma \gamma Q\overline{Q})=\sigma _{}\sigma _{}`$ $``$ $`\mathrm{Re}{\displaystyle M_{++}^{Born}(\gamma \gamma Q\overline{Q})M_{}^{oneloop}{}_{}{}^{}(\gamma \gamma Q\overline{Q})}`$ $`+\mathrm{Re}{\displaystyle M_{++}^{oneloop}(\gamma \gamma Q\overline{Q})M_{}^{Born}{}_{}{}^{}(\gamma \gamma Q\overline{Q})}`$ $`+\mathrm{Re}{\displaystyle M_{++}^{Born}(\gamma \gamma Q\overline{Q}g)M_{}^{Born}{}_{}{}^{}(\gamma \gamma Q\overline{Q}g)}`$ here the sum over the helicities of the final state particles $`Q\overline{Q}(g)`$is implied. The Born amplitude of the $`Q\overline{Q}`$ pair production in the photon-photon collisions is known to vanish like $`m_Q^2/s`$ for equal photon helicities and massless quarks . Therefore first two terms in the Eq. 19 vanish in the high energy limit. In addition, helicity amplitudes for the process of massless $`Q\overline{Q}`$ pair production with the additional gluon emission $`\gamma \gamma Q\overline{Q}g`$ identically vanish for photon and gluon helicities $`\lambda _1=\lambda _2=\lambda _g=\pm 1`$ and arbitrary quark helicities . Consequently, in the third term of the Eq. 19 amplitudes $`M_{++}^{Born}(\gamma \gamma Q\overline{Q}g)`$ and $`M_{}^{Born}(\gamma \gamma Q\overline{Q}g)`$ are nonzero only in the case when emitted gluons have different polarizations. Therefore there is no interferention between corresponding amplitudes and third term of Eq. 19 term also vanishes at high energies. As result, the difference of the cross sections for $`(Q\overline{Q})`$-pair production, $`\mathrm{\Delta }\sigma (\gamma \gamma Q\overline{Q})`$ is suppressed by a factor of $`O(m_Q^2/s)`$. For $`b\overline{b}`$ production the relative difference of the cross section for parallel and orthogonal polarized photons is less than $`1\%`$ for $`\sqrt{s}200`$ GeV, i.e. in the whole range of the PLC energies. On the other hand, for top-antitop production there is no strong suppression of $`\mathrm{\Delta }\sigma `$ by the mass of quark at energies $`\sqrt{s}500÷800`$ GeV. The production cross sections for $`t\overline{t}`$-pair production are illustrated in the Figure 2 for different helicity states of initial photons. The usual cut for suppression the Higgs background is imposed, $`\|\mathrm{cos}\theta \|<0.7`$. The solid lines correspond to the Born level cross sections and the dashed lines to the QCD corrected ones. As one can see from Figure 2 the corrections are large near the threshold and decrease very rapidly with increasing the photon-photon c.m.s energy, $`W_{\gamma \gamma }`$. In the Figure 3a the difference of two cross sections, $`\mathrm{\Delta }\sigma (\gamma \gamma t\overline{t})`$, is shown. The correction to the $`\mathrm{\Delta }\sigma `$ are rather large near the threshold, up to $`W_{\gamma \gamma }400`$ GeV, and decreases rapidly. However, the asymmetry, $`(\sigma _{}\sigma _{})/(\sigma _{}+\sigma _{})`$ gets only small corrections in the whole range of energies. The azimuthal asymmetry of heavy quark pair production is the specific effect which occurs only at one-loop level. The relevant term in Eq. 8 proportional to imaginary unit $`i`$ corresponds to the scattering when one photon is polarized linearly and the other circularly, $`\xi _1^{(1)}=\pm 1`$ and $`\xi ^{(2)}=\pm 1`$ or other way round (all other Stokes parameters are zero). This term gives the contributions with opposite signs to the cross section being integrated in two different ranges of azimuthal angle, $`\pi /2<\varphi <0`$ and $`0<\varphi <\varphi /2`$. The value $`\varphi =0`$ corresponds to the direction of polarization vector of linearly polarized photon. It is obvious that the total contribution of $`_{12}`$-term to the cross section is zero. We define the azimuthal asymmetry of $`(Q\overline{Q})`$-pair production in the linearly polarized photon scattering off circularly polarized one in the following way $`A_\varphi ={\displaystyle \frac{\sigma (\frac{\pi }{2}<\varphi <0)\sigma (0<\varphi <\frac{\pi }{2})}{\sigma (\frac{\pi }{2}<\varphi <0)+\sigma (0<\varphi <\frac{\pi }{2})}}.`$ (20) This asymmetry can be sizeable only for top-antitop production at PLC energies, while for the $`b\overline{b}`$-pair production azimuthal asymmetries are also suppressed by factor $`O(m_b^2/s)`$. The expectations for the $`t\overline{t}`$-pair production asymmetry are shown in Figure 3. The effect is $`1.5\%`$ at energies of the PLC. ## 6 Conclusion In the presented paper we derived the compact expressions for the $`\alpha _s`$-corrections to the squared matrix element $`\gamma \gamma Q\overline{Q}`$ for the general case of initial photon polarizations. The total cross sections up to the order $`\alpha ^2\alpha _s`$ are calculated, which are given by the sum of the tree level cross section, contribution of the interference term between the QCD one-loop $`\alpha _s`$-correction and tree level amplitude and tree level cross sections of quark pair production accompanied by the real gluon emission. The numerical results are obtained for heavy quark pair production cross sections in the case of linearly polarized photon collisions. The difference of the cross sections of heavy quark pair production for parallel and perpendicular polarized photon collisions is suppressed by factor $`m_Q^2/s`$. We show that the QCD correction for the $`b\overline{b}`$ production asymmetry is less than $`1\%`$ in the whole energy range of the PLC and practically does not change the background for the measurement of $`CP`$ parity of Higgs boson. At the same time, there is no such suppression for top-antitop production due to the large mass of top quark. The relevant cross sections for $`t\overline{t}`$-pair productions are calculated. The QCD corrections are large near the threshold and decrease rapidly with increasing the c.m.s energy of colliding photons. We also calculated the azimuthal asymmetries for $`t\overline{t}`$ production. This effect arises only at one-loop level and the asymmetry is about $`1.5\%`$ in the energy range of the PLC. Acknowledgments We would like to thank J.I. Illana for useful discussions.
warning/0004/astro-ph0004182.html	ar5iv	text	# Cold dust in the starburst galaxy M 82 ## 1 Introduction The prototypical starburst galaxy M 82 is the prime example of a galaxy in which the violent star formation activity gives rise to the formation of a bipolar outflow and an associated extended halo, which is visible in various regimes of the electro–magnetic spectrum. Outflowing material was first evident in the H$`\alpha `$ light (Lynds & Sandage 1963, McCarthy et al. 1987, Bland & Tully 1988). The kinematics of the outflowing material and the geometry of the cone have been thoroughly worked out by McKeith et al. (1995). The recent detection of H$`\alpha `$ emission $``$11 kpc away from the plane of M 82 (Devine & Bally 1999) determines the cap of this outflow. The whole scenario is further corroborated by the existence of vertical magnetic field lines (Reuter et al. 1994), along which relativistic particles partake in the outflow, thus forming a radio halo (Seaquist & Odegard 1991). This halo exhibits a filamentary structure away from the plane (Reuter et al. 1992). Hot gas has been also found in some kind of halo which extends for several arcminutes along the minor axis (Fabbiano 1988, Schaaf et al. 1989, Bregman et al. 1995). This X–ray emission correlates well with the H$`\alpha `$ if observed with sufficiently high angular resolution (Watson et al. 1984). In this context, the search for neutral gas and dust away from the plane of M 82 has seen a number of attempts. Owing to the close interaction between M 81 and M 82, neutral hydrogen is seen enveloping M 82, but whether this gas stems from the M 81 or has been expelled from M 82 is still a matter of debate (Yun et al. 1994). Molecular gas associated with the dusty filaments outside the plane of M 82 has been reported by Stark & Carlson (1984) and has been observed by Sofue et al. (1992) out to a projected distance of $`\pm `$2 kpc. Clear evidence for the existence of scattering dust away from the plane of M 82 comes from measurements of optical polarization (Bingham et al. 1976, Notni et al. 1981, Neininger et al. 1990). Owing to its high luminosity in all spectral bands, M 82 has been a prime target for first–light experiments in the mm– and submm–regime (Elias et al. 1978, Jura et al. 1978). With the improvement of bolometric measurements detailed studies of the distribution of the cold dust in this galaxy have become feasible; thus the recent past has seen an increasing number of such investigations ($`\lambda `$2 mm: Kuno & Matsuo 1997; $`\lambda `$1 mm: Thronson et al. 1989, Hughes et al. 1990, Krügel et al. 1990a, Alton et al. 1999; $`\lambda `$0.5 mm: Jaffe et al. 1984, Smith et al. 1990, Alton et al. 1999). Carlstrom & Kronberg (1991) have shown that at $`\lambda `$3 mm thermal free–free emission still dominates the overall spectrum. Apart from these dedicated bolometric measurements, continuum maps at millimeter wavelengths have been produced as by-products in various observations of spectral lines (e.g. Neininger et al. 1998). As a result of the discoveries of a galactic wind and halo in M 82, the more recent bolometric measurements aimed at detecting a dust halo at mm wavelengths. CO emission away from the plane was reported by Stark & Carlson (1984), Nakai et al. (1986) and by Sofue et al. (1992). An outflow of molecular gas was claimed by Nakai et al. (1987). A dust continuum halo was first mentioned by Hughes et al. (1990). Observations with higher resolution (Kuno & Matsuo 1997) suggest emission at $`\lambda `$2 mm out to 400 pc from the plane. Recently Alton et al. (1999) reported a dust outflow from the central region of M 82, based on their sub-mm images. Here we report observations of M 82 in the $`\lambda 1.2`$ mm continuum performed with the MPIfR 19-channel bolometer at the IRAM 30-m telescope. In Sect. 2 we describe the observations and data analysis, with consideration of the possible influence of the error beam. In Sect. 3 the distribution of the cold dust in M 82 will be presented, with a discussion conducted in Sect. 4, along with a comparison with other published measurements. In Sect. 5 we give a short summary of our results. ## 2 Observations and data reduction ### 2.1 Observations We used the 19-channel bolometer of the MPIfR (Kreysa et al. 1993) in the Nasmyth focus of the IRAM 30-m telescope during the spring session in 1997. The weather conditions were fair during our night-time runs, with rather stable zenith opacities between 0.2 and 0.3 at 230 GHz. The 19 channels of the bolometer are arranged in a closely packed hexagonal array, with beamsizes of 11<sup>′′</sup> FWHM and spacings of 20<sup>′′</sup>. The calibration was performed by mapping Uranus every morning. The pointing and focus checks were made at regular intervals during the observations using nearby quasars (in particular 1308+326). The map was centered on the near–infrared peak at $`\alpha _{50}=09^\mathrm{h}51^\mathrm{m}43.^{\prime \prime }4`$, $`\mathrm{\Delta }_{50}=69{}_{}{}^{}55{}_{}{}^{}01^{\prime \prime }`$ (Joy et al. 1987). The data were taken in the standard mapping mode where the field is scanned at a constant speed of 4<sup>′′</sup> per second in azimuth, with subsequent scans being spaced by 4<sup>′′</sup> in elevation. The map size was chosen large enough (320<sup>′′</sup>$`\times `$ 270<sup>′′</sup>) to assure a good baseline determination in the presence of extended emission. We thus made sure that the field was mapped sufficiently far out perpendicular to the major axis of M 82, such as to trace the full extent of any dust component in the halo. While scanning the subreflector was wobbled at a 2 Hz rate with an amplitude of 45<sup>′′</sup>, which is small enough to ensure a stable beam pattern. Due to the small wobbler throw the OFF position was not necessarily free of emission, but using the algorithm of Emerson et al. (1979) this should have no effect on the resulting map. We also mapped Mars once and the point source 3C 273 twice each in order to obtain a detailed beam pattern. In addition, more maps of Mars and other calibration sources were used to properly check the calibration. ### 2.2 Standard data reduction The whole field was covered six times. For each coverage we performed the data reduction separately using NIC (GILDAS software package), starting with the subtraction of a third–order baseline. The scanning mode described in the previous section leads to maps which contain a superposition of a positive and a negative image of the source (double-beam maps). In principle the deconvolution can be done by dividing the Fourier transform (FT) of the measured intensity distribution by the FT of the wobble function and transforming back to image space. As the FT of the wobble function is a sine wave with zeroes at the origin and at harmonics of the inverse wobbler throw, it would cause problems at these spatial frequencies. To avoid these problems we used the algorithm of Emerson et al. (1979) to restore the double–beam maps into equivalent single–beam maps. Thereafter we combined the 19 channels and calibrated the resulting map with respect to Mars and Uranus (assuming brightness temperatures of 198.5 K and 97.2 K respectively). In a last step the four best coverages were averaged and regridded onto equatorial (B1950) coordinates. (Two coverages were ignored because of bad S/N ratios.) This final map has a sensitivity of 5 mJy per beam and a angular resolution of 12<sup>′′</sup> (HPBW). ### 2.3 Influence of the error beam In order to check the influence of the error beam on the extended emission we applied a CLEAN algorithm developed for single–dish maps (Klein & Mack 1995) to each coverage using a radially symmetric antenna pattern. This pattern was obtained in the following way: A field of 4$`.^{}`$5 $`\times `$ 3$`.^{}`$5 size was mapped centred on 3C 273. This antenna pattern was then rotated in small steps ($`5^{}`$), and these individual (rotated) maps were averaged to yield a nearly radially symmetric pattern. The subtraction of the CLEANed map from the original one showed that the error beam contributes less than 25% of the rms noise level to the extended emission ($`20`$<sup>′′</sup> away from the central position). Even the sidelobes can be ignored in our case, because at $`\lambda `$ 1.2 mm they are of the order of 2%, corresponding to roughly 5 mJy/beam in the M82 maps. Averaging four coverages with different parallactic angles (and accordingly different positions of the sidelobes) the effect of the antenna pattern on the extended emission falls well below the rms noise of our final map. This extra analysis make sure that the extended emission seen in our map is real and not affected by the error beam. ### 2.4 Correction for CO emission Using the 19-channel bolometer, with its central frequency of 240 GHz and its effective bandwidth of $`\mathrm{\Delta }\nu _{bol}`$ 80 GHz, the <sup>12</sup>CO(2$``$1) line at 230 GHz might noticeably contribute to the total intensity in the bandpass. To get a reliable information on the distribution of the cold dust, we applied the following correction to our continuum map: We used a $`3^{}\times 3^{}`$ map of M 82 made with the same telescope in the <sup>12</sup>CO(2-1) line (obtained to provide the zero spacings for an investigation with the Plateau de Bure interferometer, Weiß et al.in prep.). This map was measured in the on–the–fly mode. The on–the–fly technique ensures a uniform sampling of the emission over the whole field with good signal-to-noise ratio. Since we knew from our continuum map that the effects of the error beam can be neglected at a frequency of 240 GHz, we did not apply any error beam correction to the CO map. Integrating over the effective bandwidth of the bolometer, the total continuum flux in the beam at the peak position is given by $$\mathrm{S}^{\mathrm{bol}}=\mathrm{\Delta }\nu _{\mathrm{bol}}\mathrm{S}_{240}^{\mathrm{bol}}=2.610^{16}\mathrm{Wm}^2$$ where $`\mathrm{S}_{240}^{\mathrm{bol}}=330`$ mJy is the peak flux at 240 GHz. The integrated intensity of the <sup>12</sup>CO(2$``$1) line at the same position $`\mathrm{T}_{\mathrm{mb}}\mathrm{dv}=625\mathrm{K}\mathrm{km}\mathrm{s}^1`$ was taken from our CO map. Expressing this in terms of frequency instead of velocity, the integrated intensity reads $`\mathrm{T}_{\mathrm{mb}}d\nu =500\mathrm{K}\mathrm{MHz}`$. To obtain the total flux in the <sup>12</sup>CO(2$``$1) line we used the Rayleigh–Jeans Approximation and assumed that T$`{}_{\mathrm{mb}}{}^{}\mathrm{T}_{\mathrm{ex}}`$: $$\mathrm{S}^{\mathrm{CO}}=\mathrm{\Delta }\nu _{\mathrm{CO}}\mathrm{\Omega }\frac{2\mathrm{kT}_{\mathrm{ex}}}{\mathrm{c}^2}\nu ^2\mathrm{\Omega }\frac{2\mathrm{k}\nu ^2}{\mathrm{c}^2}\mathrm{T}_{\mathrm{mb}}d\nu $$ With the beam solid angle $`\mathrm{\Omega }=4.510^9`$sr and the integrated <sup>12</sup>CO(2$``$1) intensity from above we get $`\mathrm{S}^{\mathrm{CO}}=3.610^{17}\mathrm{Wm}^2`$. Thus, $`\mathrm{S}^{\mathrm{CO}}/\mathrm{S}^{\mathrm{bol}}0.14`$ at the position of the continuum peak. If we assume that this value holds for the whole galaxy, which of course is a simplification, we can obtain a CO–corrected continuum map by subtracting the appropriately scaled CO map from our continuum map. The uncorrected and the corrected map are shown in Fig. 1. ## 3 Results ### 3.1 Overall distribution of the dust emission Fig. 2 shows the CO-corrected continuum map of M 82 at $`\lambda `$1.2 mm, superimposed onto an optical image. At $`\lambda `$1.2 mm M 82 has a more or less oval shape (ellipticity e $``$ 0.6), with the major axis parallel to the galactic plane (molecular gas). This corresponds to a position angle of approximately 70. The maximum of the continuum emission is located at $`\alpha _{50}=9^\mathrm{h}51^\mathrm{m}41^\mathrm{s}`$, $`\mathrm{\Delta }_{50}=69`$5455<sup>′′</sup>. Within the errors this position is identical to the peaks of the distributions at $`\lambda `$800 $`\mu `$m and $`\lambda `$1.1 mm (Hughes et al. 1990). Besides this regular shape there are some spur–like features extending above and below the galactic plane, which might be associated with the outflow seen at optical and IR wavelengths. ### 3.2 Major– and minor–axis profiles The contours in Fig. 2 indicate some extended emission also along the minor axis of the galaxy, but at first glance the extent in this direction is not as large as suggested by the maps of Kuno & Matsuo (1997) and Alton et al. (1999). We have therefore computed the intensity as a function of distance from the galactic center along the minor (major) axis, averaging over some 30<sup>′′</sup> in the major-axis (minor-axis) direction. These profiles are displayed in Fig. 3. For comparison the <sup>12</sup>CO(2$``$1) line intensity, the 850 $`\mu `$m (Alton et al. 1999) and 20 cm continuum profiles are shown, too. The $`\lambda `$1.2 mm emission shows the most concentrated distribution, significantly smaller than that at 850 $`\mu `$m. Especially the scale height along the minor axis is very small with respect to the other wavelengths and the molecular emission. Contrary to the expectation the molecular emission is the most extended component. Possible explanations for this phenomenon are discussed in Sect. 4.2. ### 3.3 Integrated flux density Integrating over the continuum map we obtain a total flux density of S$`{}_{240}{}^{}=1.91\pm 0.13`$ Jy. This is lower than the value measured by Krügel et al. (1990a), but their flux density was much more uncertain, owing to the lower sensitivity of their 1-channel bolometer. Fig. 4 shows a radio-to-IR spectrum of M 82 with the individual values listed in Tab. 1. Our new 240 GHz measurement is completely consistent with the slope of the spectrum at this frequency (indicated by the dotted line), whereas the value derived by Kuno & Matsuo (1997) at 157 GHz lies much too high with respect to all adjacent data points. The CO–corrected flux density is significantly lower than the flux density in our original map. This demonstrates that a line correction for continuum measurements in the relevant bands is indispensible. The total flux can be turned into a gas mass via $$\mathrm{M}_\mathrm{g}=\frac{\mathrm{S}_{240}\mathrm{D}^2}{\mathrm{B}(\mathrm{T}_\mathrm{d})\kappa _{240}^{}}\frac{\mathrm{M}_\mathrm{g}}{\mathrm{M}_\mathrm{d}}.$$ With a distance of D = 3.25 Mpc, and using the same parameters as Krügel et al. (1990a), i.e. a dust temperature of T<sub>d</sub> = 30 K (Chini et al. 1989) and a dust absorption coefficient $`\kappa _{240}=\kappa _{240}^{}\frac{\mathrm{M}_\mathrm{d}}{\mathrm{M}_\mathrm{g}}=0.003\mathrm{cm}^2\mathrm{g}^1`$ (Krügel et al. 1990b), we derive a gas mass of $`\mathrm{M}_\mathrm{g}=7.510^8\mathrm{M}_{}`$. If we adopted the Draine & Lee value $`\kappa _{240}=0.005\mathrm{cm}^2\mathrm{g}^1`$ (Draine & Lee 1984) the resulting gas mass would be somewhat lower ($`\mathrm{M}_\mathrm{g}=4.510^8\mathrm{M}_{}`$). In this calculation we assumed a gas-to-dust ratio of approximately M<sub>g</sub>/M$`{}_{\mathrm{d}}{}^{}100`$. The corresponding dust mass is M$`{}_{\mathrm{d}}{}^{}=7.510^6M_{}`$. ## 4 Discussion ### 4.1 A dust halo? The $`\lambda `$2 mm map (Kuno & Matsuo 1997) shows continuum emission with a box–like shape, indicating a dust halo around M 82. Although the emission at $`\lambda `$1.2 mm and $`\lambda `$2 mm should be similar, we do not see any /\*-evidence for a significant halo contribution to the cold dust. This difference could be due to the error beam of the Nobeyama 45-m telescope. Effects of the error beam and the antenna pattern have been investigated and largely removed from our map (see Sect. 2.3), while this is not mentioned in case of the Nobeyama measurements. The total flux density in the Nobeyama map is 1.7 Jy, which is much more than expected at this wavelength (see Fig. 4). This extra flux might be due to emission entering through sidelobes and the error beam. The 850 $`\mu `$m map made by Alton et al. (1999) with the JCMT shows continuum emission at the 50 mJy level up to 40<sup>′′</sup> above the galactic plane. This translates to 15 mJy at $`\lambda `$1.2 mm, corresponding to the second contour (3$`\sigma `$) in our map. Considering only the overall extent of the emission, the JCMT map is consistent with ours. The different shapes might be due to problems with the baseline subtraction, because the area mapped by Alton et al. is only slightly larger than the extent of the continuum emission. ### 4.2 Comparison of dust and CO In non–active galaxies like NGC 4565 (Neininger et al. 1996) and NGC 5907 (Dumke et al. 1997) the molecular line emission is more concentrated towards the central region than the thermal dust emission. As can be seen from Fig. 3 and 5 the galaxy M 82 exhibits the opposite behaviour. The CO emission appears to be much more extended than the continuum. Since the maps in Fig. 5 are made with the same telescope at the same frequency, the difference can not be due to different beam patterns and error beams. The dust is heated by the UV radiation field, which is concentrated towards the dense central region, because the UV photons can hardly escape without being absorbed. Only a small fraction of the UV photons makes its way out to kpc distances from the galactic plane, where it might contribute to the observed polarized radiation. The molecular gas can be excited by the radiation field, but also by low–energy cosmic rays and by soft X–ray photons (see e.g. Glassgold & Langer 1973), which have a much larger scale height than the UV photons (Shopbell & Bland-Hawthorn 1998). Fig. 3 shows that the synchrotron intensity at $`\lambda `$20 cm drops faster with distance from the galaxy plane than the CO and even the dust emission. However, at low radio frequencies there may be a significant radio halo (low–energy cosmic rays have longer lifetimes), which is already indicated by the 327–MHz map of Reuter et al. (1992). Hence, there may be abundant cosmic rays for heating. Such cosmic rays (E $``$100 MeV) remain, however, invisible in the radio window, since even in the strong ($`\mathrm{B}10\mathrm{}50\mu `$G) magnetic field of M 82 they produce synchrotron radiation at frequencies below about 10 MHz. Soft X–rays are also seen far out of the plane of M 82 (e.g. Bregman et al. 1995). These circumstances altogether provide the likely heating sources for any molecular material transported out of the disk of M 82 into the halo, rendering the CO ”overluminous”. This is also indicated by the low conversion ratio of CO line intensity to molecular gas column density, X$`{}_{\mathrm{CO}}{}^{}=\mathrm{N}_{\mathrm{H}_2}/\mathrm{I}_{\mathrm{CO}}`$ derived by Smith et al. (1991). The existence of a soft X–ray and a low–frequency radio halo around M 82 readily explains the larger extent of the CO emission as compared to the cold dust. ## 5 Summary In this paper we presented a new high–sensitivity map of the continuum emission at $`\lambda `$1.2 mm in the starburst galaxy M 82. To get a reliable map of the extended thermal dust emission we applied a CLEAN algorithm to our map (error beam correction) and subsequently subtracted the <sup>12</sup>CO(2-1) line emission. Besides this the contribution of the sidelobes of the IRAM 30-m telescope has been estimated. The total mass of dust in the inner 3 kpc of the galaxy derived from the integrated $`\lambda `$1.2 mm flux density is $`5.010^6`$ M, the inferred total mass of gas is $`7.510^8`$ M . Although the continuum emission is not confined to the galactic disk, our map rules out the existence of any pronounced halo of cold dust around M 82. We find the CO(2$``$1) line emission to be clearly more extended than the $`\lambda `$1.2 mm dust emission. This unexpected observational result now waits for a conclusive theoretical explanation. ###### Acknowledgements. We thank Dr. A. Sievers from IRAM (Granada) for his support with additional calibration. We are very grateful to Dr. P.B. Alton for making his JCMT data available to us.
warning/0004/astro-ph0004088.html	ar5iv	text	# A kinematic and morphological investigation of the asymmetric nebula around the LBV candidate WRA 751 ## 1 Introduction With today’s initial mass function, stars with masses of $`M50100`$ M and luminosities of $`L10^{56}`$ L top the Hertzsprung-Russell diagram (HRD) and represent the most massive stars known. These massive stars show a remarkable evolutionary behavior (e.g., Schaller et al. 1992, Stothers & Chin 1996). After spending a ‘normal’ life as O stars on the main sequence they evolve towards cooler temperatures while entering a phase with very high mass loss (up to 10<sup>-4</sup>M yr<sup>-1</sup>), they become Luminous Blue Variables (LBVs). This phase starts when the stars reach the Humphreys-Davidson limit (Humphreys & Davidson 1979, 1994, Langer 1994) in the HRD. Analyzing HRDs of the Galaxy and the LMC, Humphreys (1978, 1979) and Humphreys & Davidson (1979) found a lack of very luminous red supergiants. Apperently the most massive stars do not evolve into red supergiants but instead their evolution is reversed towards the blue supergiant part in the HRD. The turning points form the empirical Humphreys-Davidson limit. Around this Humphreys-Davidson limit in the HRD LBVs are found. One of the most prominent characteristics of the unstable LBV phase are strong stellar winds and possible giant eruptions, which lead to the peeling off of parts of the stellar envelope and the formation of small circumstellar nebulae, the so-called LBV nebulae (LBVN, e.g. Nota et al. 1995). Only very few LBVs are known in our Galaxy (8 classified and candidate objects) and a few in other galaxies, for instance in the LMC, SMC, M31, and M33 (Humphreys & Davidson 1994). Altogether nearly 40 LBVs and good candidates are currently known. Not all of them, but many show circumstellar nebulae. In this paper we analyze the LBVN around the galactic star WRA 751 and put it into context with the well known and better studied LBVs $`\eta `$ Carinae and HR Carinae. WRA 751 ($`=`$ He 3-591) was first identified as a possible Wolf-Rayet star by Henize (see Roberts 1962) and then appeared in the Carlson and Henize (1979) sample of southern peculiar emission-line stars, where the authors already recognized the strong \[Fe ii\] lines and classified it as Bep. In addition they proposed a similarity between WRA 751 and HR Car (which in the meantime is known to be a LBV). Using spectroscopic and photometric observations Hu et al. (1990) noted for the first time that WRA 751 is a good LBV candidate. In addition to the characteristically strong \[Fe ii\] lines this star shows an irregular light variation typical for LBVs. In their analysis the star was classified as O9.5 with $`M_{\mathrm{bol}}=9.6`$<sup>m</sup>. They derived $`T_{\mathrm{eff}}=\mathrm{30\hspace{0.17em}000}`$ K, $`E(\mathrm{B}\mathrm{V})=1.8,M=50\mathrm{M}_{\mathrm{}}`$, and a distance of about 5 kpc. One year later Hutsemékers & Van Drom (1991a) found that WRA 751 was surrounded by a nebula roughly 22″ in diameter. They found the nebula to be nearly spherically symmetric and to show a diffuse structure. Their low-dispersion spectra revealed a nebula of low excitation with an electron temperature of $`T_{\mathrm{e},\mathrm{neb}}<\mathrm{15\hspace{0.17em}000}`$K, an electron density of $`n_\mathrm{e}400\mathrm{cm}^3`$, and an expansion velocity of $`v_{\mathrm{exp}}=26`$km s<sup>-1</sup>. Assuming a distance of 7 kpc (derived from radial velocities and making use of the galactic rotation curve) their radius of the nebula is 0.38 pc. Analysis of infrared data (see de Winter 1992) show a NIR excess from which an estimate for the mass loss of $`\dot{M}10^{5.5\mathrm{}6}M_{\mathrm{}}\mathrm{yr}^1`$ with a velocity of $`v_{\mathrm{wind}}=500`$km s<sup>-1</sup> were derived. They found a cool dusty circumstellar shell with strong emission in the FIR to surround WRA 751. In 1992, from interstellar and circumstellar reddening, van Genderen measured a distance of $`45`$kpc for WRA 751 as a lower limit. García-Lario et al. (1998) re-determined the characteristic parameters and found $`T_{\mathrm{eff}}\mathrm{25\hspace{0.17em}000}`$K, $`v_{\mathrm{exp}}=24`$km s<sup>-1</sup>, and $`n_\mathrm{e}200\mathrm{cm}^3`$. In summary, the published observations and their interpretations point towards WRA 751 being a LBV star with a slowly expanding nebula. ## 2 Observation and data reduction ### 2.1 Imaging To compare kinematics and morphology we used images of the nebula around WRA 751 taken with the Super Seeing Imager (SUSI) on the New Technology Telescope (NTT) of the European Southern Observatory (ESO). The images were retrieved from the ESO NTT archive and reduced in a standard way, using flat-fields and bias frames from the same observing night. Besides of the emission-line images in H<sub>α</sub> and \[N ii\] a broad band red continuum filter image was used for the continuum subtraction in order to be able to distinguish the pure line emission from the continuum<sup>1</sup><sup>1</sup>1 Due to missing documentation in the NTT archive and especially the lack of information in the header of the images a more detailed specification of the filters can not be given here.. Fig. 1 shows a 1050 s exposure H<sub>α</sub> image (not continuum subtracted) with a round field of view, about 50″ in diameter. From the image, we determined the seeing as 0$`\stackrel{}{.}`$8. As the image was taken using a coronograph, an occulting bar obscures the central star and some parts of the nebula. A non-occulted image of the nebula can be found in Hutsemékers & van Drom (1991a). ### 2.2 Long-slit echelle spectroscopy To analyze the nebula around WRA 751 in more detail we preformed high-resolution long-slit echelle spectroscopy. We used the echelle spectrograph on the 4 m telescope at the Cerro Tololo Inter-American Observatory. Inserting a post-slit H<sub>α</sub> filter (6563/75 Å) and replacing the cross-disperser by a flat mirror we selected the H<sub>α</sub> line and two \[N ii\] lines at 6548 Å and 6583 Å. We picked the 79 l mm<sup>-1</sup> echelle grating and a slit-width of 150 $`\mu `$ (corresponding to 1″), which lead to an instrumental FWHM at the H<sub>α</sub> line of about 10 km s<sup>-1</sup>. All data were recorded with the long focus red camera and a $`2048\times 2048`$ Tek2 CCD.The pixel size was 0.08 Å pixel<sup>-1</sup> along the dispersion and 0$`\stackrel{}{.}`$26 pixel<sup>-1</sup> on the spatial axis. Vignetting limited the slit length to $`4^{}`$ . Seeing was $`2`$″ during the observations and the weather was not photometric. Thorium-Argon comparison lamp frames were taken for wavelength calibration and geometric distortion correction. The slit was oriented in two different position angles (PA) perpendicular to each other (Fig. 2). We observed 5 positions with $`\mathrm{PA}=140`$°. One position was centered on the star (in our nomenclature Slit center), two positions were offset to the north by 8″ and 12″ from the central star (Slit 8N and Slit 12N) and two positions were offset by 8″ and 16″ south of the central star (Slit 8S and Slit 16S). For $`\mathrm{PA}=50`$° we observed at two positions offset from the star by 9″ to the north and south, respectively (Slit 9N and Slit 9S). The left column of Fig. 3 shows the echellograms of the slits at $`\mathrm{PA}=140`$°, Fig. 4 those at $`\mathrm{PA}=50`$°. All echellograms cover a spectral range of 75 Å centered on H<sub>α</sub> and extend about 2′ in the spatial direction, centered on the projected position of the central star. The top of each echellogram in Fig. 3 points to the north-west, and in Fig. 4 to the south-west. Some telluric lines are visible and were used to improve the absolute wavelength calibration. ## 3 Morphology of the nebula around WRA 751 The first image taken of the nebula around WRA 751 was published by Hutsemékers & van Drom (1991a). They discovered a spherical diffuse nebula about 22″ in size. We retrieved and re-analyzed \[N ii\] and H<sub>α</sub> images retrieved from the ESO NTT archive in order to compare the morphology and the kinematics of the nebula. Figure 1 shows that the nebula around WRA 751 indeed appears nearly round and almost spherically symmetric (see also Nota 1998). We measure a diameter of the main body of the nebula of 22$`\stackrel{}{.}`$8, which corresponds to 0.50 pc assuming the lower limit distance by van Genderen et al. (1992; $`4.5`$kpc). Appearing rather homogeneous in its surface density and closely attached to the central star, it resembles very much a Stromgren sphere. However, this is in contrast to the morphological appearance of most other LBV nebulae like that of AG Car or HR Car (e.g. Hutsemékers & van Drom 1991b, Nota et al. 1995, Smith et al. 1997) which show a more intense detached ring structure, or as in the case of HR Car at least parts of a bipolar shell (Weis et al. 1997, Nota et al. 1997). This is mainly due to an increased matter density along the line of sight when looking through the edge of a shell. A closer inspection of the images allows us to locate spatial variations in the surface density of the nebula around WRA 751: a brighter circular half-shell shows up in the eastern part, while—in comparison—in the western part the brightness decreases by about a factor of two. Even more remarkable are two extensions of the spherical shell, one nearly exactly to the north at a position angle of 340° which we will call the Northern Cap (Fig. 1) and the other to the south (Southern Cap) at $`\mathrm{PA}=165`$°. This means that the Caps are almost along an axis through the star. Morphologically both appear roughly triangular in shape. The surface brightness of the Northern Cap is somewhat larger, that of the Southern Cap somewhat smaller than the low surface part (the west side) of the central spherical nebula. The Northern Cap extends by 2$`\stackrel{}{.}`$6 beyond the main body of the nebula, the Southern Cap by 4$`\stackrel{}{.}`$7. ## 4 Kinematic indication of a bipolar structure Using high-resolution long-slit echelle spectra an analysis of the kinematics of the LBVN around WRA 751 was performed. With the 7 slit positions our mapping covers the entire nebula. The right columns of both, Figs. 3 and 4, show the position-velocity diagrams ($`pv`$-diagrams) corresponding to the echellograms in the same line. For the $`pv`$-diagrams we used the stronger \[Nii\] line at 6583Å which is less contaminated by background emission. Additionally all measurements were compared to those of H<sub>α</sub> to ensure that no kinematic difference exists between the two lines. All velocities were corrected to the Local Standard of Rest (LSR) system. The 0″ position in all $`pv`$-diagrams corresponds to the projected position of the central star (WRA 751) onto the slit, from there positive offsets are to the north-west for the $`\mathrm{PA}=140`$° slits (or south-west for $`\mathrm{PA}=50`$°) and negative to the south-east (or north-east, respectively). All spectra (except Slit 16S) clearly resolve the Doppler ellipse, and indicate a predominantly spherical expansion of the nebula. The largest radial expansion velocity $`v_{\mathrm{exp}}`$ was found in the central position (0″) of Slit 9S ($`\mathrm{PA}=50`$°) with a value of 29.4 km s<sup>-1</sup>$`\pm `$ 2 km s<sup>-1</sup>, the second largest in Slit 8N at 28.3 km s<sup>-1</sup>$`\pm `$ 2 km s<sup>-1</sup>. A similar expansion velocity was derived at the same position in the central slit at 26.3 km s<sup>-1</sup>$`\pm `$ 2 km s<sup>-1</sup>. Since this slit crosses the star in the center, the velocity measurements are less certain at the 0″ position. There $`v_{\mathrm{exp}}`$ might be even higher as can be estimated from interpolating for the missing data points in the $`pv`$-diagram of Fig. 3. If a nebula is spherically expanding, the largest expansion velocity should occur at the position projected onto the star. Even though this is not the case here, the difference of the expansion velocity between the central slit and the Slit 8N is not significant and within the errors. However, a comparison with Slit 9S ($`\mathrm{PA}=50`$°) shows an asymmetric shape of the expansion structure indicating that the nebula is most likely not perfectly round but perhaps a tilted ellipsoid with the maximum expansion off-centered to the east of WRA 751. There the largest expansion velocities are found. Another possibility is that the nebula is spherical with a bump at this point. This is best visible in the asymmetric shapes of the expansion ellipses in Slits 8N (Fig. 3) and Slit 9S ($`\mathrm{PA}=50`$°; Fig. 4). While the largest negative velocity is found at the 0″ position, the largest positive velocity is located at $`2`$″, i.e., the Doppler ellipse is asymmetric with respect to the projected position of the star. All derived global expansion velocities are in good agreement with expansion velocities found in the literature (Hutsemékers & van Drom 1991a, García-Lario et al. 1998). In almost all spectra the redshifted line is more intense and broader than the blueshifted side. Since it is the redshifted component which is stronger, this brightness difference cannot be accounted by an absorption of the blueshifted component. It is either an intrinsic brightness difference or results from a thicker shell at the backside which leads to a higher emission measure. The FWHM of the redshifted wing of the expansion ellipse is resolved at 21 km s<sup>-1</sup> (instrumental FWHM $`10`$km s<sup>-1</sup>) and supports the interpretation as a thicker shell. The widths itself is indicative of a stratification of the radial velocities along a line of sight, presumably with the inner parts of the nebula moving slower. A wider redshifted component can be explained as due to a thicker shell there. Fig. 5 shows a composite $`pv`$-diagram of all slits with $`\mathrm{PA}=140`$°. It shows the trend of the expansion velocities across the nebula: * The split of the Doppler ellipse decreases as we move away from the geometric center of the bubble. The decrease, however, is not in agreement with a purely spherical expansion. While the redshifted part of the expansion ellipse decreases more like that of a spherical expansion, the blueshifted component stays nearly at a constant velocity for a given position along the slit. * For a spherical expansion not only the width of the expansion ellipse in velocity space shrinks, but it also becomes narrower in spatial extent. In the WRA 751 nebula, however, we find that the convergence points at positive positions migrate to smaller values as one proceeds to slits farther away from the star while they are at approximately the same location for negative values. Both points are in agreement with an asymmetry within the nebula and indicate deviations from a spherical expansion. Besides the—asymmetric—expansion of the shell, the spectra show additional kinematic components. At slit position 8N the redshifted side converges with the blueshifted side in the north-west, at the position of about 10″ as expected at this position for an expanding bubble. At the same point, however, a component which is more redshifted as the nebula appears, and extends to a position of 13″. The redshift of this component increases to 32 km s<sup>-1</sup>above the central expansion velocity as one moves away from the nebula. In Fig. 3 a dashed line in the $`pv`$-diagrams for Slits 8N, center, and 8S indicates the rest velocity of the kinematic center of the expansion at 10 km s<sup>-1</sup>, 14 km s<sup>-1</sup>, and 12 km s<sup>-1</sup>, respectively. A similar kinematic behavior can be found in Slits central and 8S, but this time at the south-eastern rim of the nebula. At about the position $`12`$″ in both spectra an additional blueshifted component appears with maximum velocities of 12 and 18 km s<sup>-1</sup> below the respective central expansion velocities. Again the velocities of these components reach their maximum at the location most distant from the star. The position of the two kinematic extensions is very symmetric with respect to the central star and thus indicative of a bipolar kinematic structure. At Slit 8N material at the rim of the inner (nearly) spherical nebula moves away from us while in Slit 8S material at the corresponding positions, namely the edge of the nebula, approaches us. With respect to the center of expansion of the central spherically expanding nebula these kinematic components are bipolar. The spatial distribution of the velocity field corresponds well to the morphology discussed above: The main part of the expansion (between the two convergence points of the expansion ellipse) measures 23″ in diameter and thus agrees with the diameter of the main body of the nebula as determined from the image to be 22$`\stackrel{}{.}`$8. Moreover, it is noteworthy that the sizes of the kinematic extensions as measured from the spectra and of the Northern and Southern Caps as determined from the images agree equally well. To ensure that the Caps are part of the nebula around WRA 751 and not just a background object onto which the nebula is projected, we have determined the line ratio \[N ii\] 6583Å/H<sub>α</sub> after subtracting contributions of the telluric H<sub>α</sub> line. Within the accuracy of our measurements, the entire nebula, including both Caps, shows a constant value of \[N ii\] 6583Å/H$`{}_{\alpha }{}^{}=1.2\pm 0.1`$, in contrast with the value for the background material of 0.3. The background line ratio is compatible with the value for a galactic H ii region (Shaver et al. 1983). This difference ensures that we can disentangle source and background contributions. ## 5 Discussion and conclusions A comparison of detailed kinematic data and the morphology of the LBVN around the galactic star WRA 751 indicates that the expansion is not as perfectly spherically symmetric as previously thought. The redshifted central part of the nebula expands nearly spherically while the blueshifted part appears stalled and moves at an almost constant radial velocity of about $`15`$km s<sup>-1</sup>(Fig. 5). A natural reason may be that the front side of the nebula expands into a denser medium and that thus the expansion is decelerated. The maximum expansion velocity was found off-centered from the central star adding a further piece of evidence that the nebula is neither exactly spherical nor expanding accordingly. The off-centered maximum velocity is either an evidence for an elongated (elliptical) shape and tilt or for a bump on the back side of the nebula. Two kinematic extensions, one redshifted, one blueshifted appear attached to the ridge of the nebula. A comparison with the morphology of the nebula around WRA 751 allows to identify the redshifted extension (Slit 8N) with the Northern Cap and the blueshifted one (Slits center and 8S) with the Southern Cap (see Cap and Slit positions in Figs. 1 and 2). The Northern and the Southern Cap therefore indicate morphologically and kinematically bipolar components in the nebula around WRA 751. While one may speculate about a jet like structure or—similar to planetary nebulae (for a discussion of similarities between LBVs and planetary nebulae, see, e.g., Frank 1998)—wider funnels, or an altogether elliptical form of the nebula creating the bipolar appearance, due to the limited spatial resolutions of both the available kinematic data and the available images, higher resolution observations will be necessary to further constrain the structure of the nebula around WRA 751. Nonetheless, from the data presented one can draw several conclusions: * The LBVN around WRA 751 consists of an expanding shell of 23″ in diameter, with a thicker shell at the back side. The radial velocity distribution deviates from that of a simple spherical expansion pattern. We find bipolar contributions to the morphology as well as radial velocity distribution in approximately north-south direction (Northern and Southern Caps). * The maximum expansion velocities are off-centered which most likely can be accounted for by a bump in the back side east of the central star. * The \[N ii\] 6583Å/H<sub>α</sub> ratio is considerably larger than that of the background material. This is typical for CNO-processed material. These three results are fully consistent with WRA 751 being a true member of the LBV class of stars. In many of the well investigated ones bipolar components to the morphological as well as kinematic structures of the nebula are found. $`\eta `$ Car (Duschl et al. 1995, Morse et al. 1998) and HR Car (Weis et al. 1997, Nota et al. 1997) are the most prominent examples. As stars in late phases of their evolution, LBVs quite naturally show CNO-processed material in their envelopes and ejecta (García-Segura et al. 1996, Smith et al. 1978). A large \[N ii\] 6583Å/H<sub>α</sub> ratio in a nebula is often used as one of the indicators for a star being a LBV. One finds, for instance, values of \[N ii\] 6583Å/H$`{}_{\alpha }{}^{}=3\mathrm{}7`$ for $`\eta `$ Car (Davidson et al. 1982, Meaburn et al. 1987, 1996, and Weis et al. 1999), $`=0.4\mathrm{}0.9`$ for HR Car (Hutsemékers & van Drom 1991b, and Weis et al. 1997), and $`=0.7`$ for AG Car (Thackeray 1977 and Smith et al. 1997). Combining our morphological, kinematic and spectroscopic results, we find mounting evidence that WRA 751 is a LBV, indeed. Moreover, there are bipolar components in its nebula, notably the two Caps. While the bipolarity in this object is less pronounced and less obvious than in other LBVs, it still strengthens our suspicion that bipolarity—albeit at different levels—is a genuine property of LBV nebulae. This makes it even more likely that modeling LBVNs as simple windblown spheres is an unjustified oversimplification. ###### Acknowledgements. The author is very grateful to Wolfgang J. Duschl and Dominik J. Bomans for many discussions on the subject of this paper, and for carefully reading and improving the manuscript. Sincere thanks go to You-Hua Chu for her support. Part of the work was carried out as a visiting graduate student at the Department of Astronomy of the University of Illinois. Its hospitality is gratefully acknowledged. The author thanks the referee, Damien Hutsemékers for helpful comments on the paper.
warning/0004/math0004009.html	ar5iv	text	# On products of harmonic forms ## 1. Introduction On a general Riemannian manifold, wedge products of harmonic forms are not usually harmonic. But there are some examples where this does happen, like compact globally symmetric spaces. For these the harmonic forms coincide with the invariant ones, and the latter are clearly closed under products. See , pp. 10-13. Sullivan observed that “There are topological obstructions for $`M`$ to admit a metric in which the product of harmonic forms is harmonic.” The reason Sullivan gave is that if the product of harmonic forms is harmonic, then the rational homotopy type is a formal consequence of the cohomology ring. Therefore, manifolds which are not formal in this sense cannot admit a metric for which the products of harmonic forms are harmonic. This motivates the following: ###### Definition 1. A Riemannian metric is called (metrically) formal if all wedge products of harmonic forms are harmonic. A closed manifold is called geometrically formal if it admits a formal Riemannian metric. Thus geometric formality implies formality in the sense of Sullivan. Compact globally symmetric spaces are metrically formal, as are arbitrary Riemannian metrics on rational homology spheres. Further examples can be generated by taking products, because the product of two formal metrics is again formal. In this paper we describe a number of elementary topological obstructions for geometric formality of closed oriented manifolds. These obstructions are independent of formality in the sense of rational homotopy theory, and are often nonzero on formal manifolds. The simplest obstruction is the product of the first Betti number and the Euler characteristic. In small dimensions these elementary obstructions are strong enough to imply: ###### Theorem 2. If $`M`$ is a closed oriented geometrically formal manifold of dimension $`4`$, then $`M`$ has the real cohomology algebra of a compact globally symmetric space. It is however not true that $`M`$ is a globally symmetric space, even up to homotopy. We give many examples of this in dimensions $`3`$ and $`4`$. We also give examples of $`4`$-manifolds which do have the real cohomology algebra of a compact symmetric space, but are not geometrically formal. This is detected by some more subtle obstructions coming from symplectic geometry and Seiberg-Witten gauge theory. The pattern of the arguments presented here is that metric formality is a weakening of a reduction of holonomy. For example, it implies that harmonic forms have constant length, though it does not imply that they are parallel. Nevertheless, the more harmonic forms there are, the stronger the constraints. In Section 7 we prove that every manifold which is not a rational homology sphere admits a non-formal Riemannian metric. This paper was motivated by an example pointed out by D. Toledo, which we describe in Section 2, and which is related to joint work in progress of the author with H. Endo. Further impetus came from a question posed by D. Huybrechts and U. Semmelmann concerning products of harmonic forms on Calabi-Yau manifolds, see . ## 2. Motivation: dimension two Let us consider first the case of a closed oriented surface $`\mathrm{\Sigma }`$. If its genus is $`0`$ or $`1`$, then there are globally symmetric Riemannian metrics. If the genus of $`\mathrm{\Sigma }`$ is $`2`$, there are nontrivial harmonic $`1`$-forms for all metrics, but every $`1`$-form has zeros. In this case every wedge product of $`1`$-forms also has zeros, but for cohomological reasons cannot vanish identically in all cases. The only harmonic $`2`$-forms are the constant multiples of the Riemannian volume form, so there cannot be any formal Riemannian metric. This proves Theorem 2 in the $`2`$-dimensional case. The argument above, pointed out to the author by D. Toledo, shows that harmonicity of products of harmonic forms can fail on compact locally rather than globally symmetric spaces, contradicting a statement in , p. 158. Note that $`\mathrm{\Sigma }`$ is formal in the sense of Sullivan, but that the nonvanishing of $`b_1(\mathrm{\Sigma })\chi (\mathrm{\Sigma })`$ obstructs geometric formality. On the sphere every metric is formal, because there are no interesting harmonic forms. However, when there are enough harmonic forms, the harmonicity of their products is a restriction on the metric enforcing rigidity. ###### Theorem 3. Every formal Riemannian metric on the two-torus is flat. ###### Proof. Let $`g`$ be a formal metric on $`T^2`$, and $`\alpha `$ a nontrivial harmonic $`1`$-form. Then $`\alpha `$ and (1) $$\alpha \alpha =\|\alpha \|^2dvol_g$$ are also harmonic and so $`\alpha `$ has constant length. In particular it has no zeros. It follows that $`\alpha (p)`$ and $`\alpha (p)`$ span $`T_p\mathrm{\Sigma }`$ for all $`p\mathrm{\Sigma }`$. As $`\|a\|`$ is constant for every constant linear combination $`a`$ of $`\alpha `$ and $`\alpha `$, the Bochner formula (2) $$\mathrm{\Delta }(a)=^{}a+Ric(a)$$ allows us to compute: $$0=\mathrm{\Delta }(\frac{1}{2}\|a\|^2)=g(^{}a,a)\|a\|^2=g(Ric(a),a)\|a\|^2.$$ This shows that the (Ricci) curvature is everywhere nonpositive. But by the Gauß-Bonnet theorem this implies that $`g`$ is flat. ∎ ## 3. Elementary obstructions Let $`M`$ be a closed oriented manifold of dimension $`n`$, and $`g`$ a formal Riemannian metric on $`M`$. As usual, we extend $`g`$ to spaces of differential forms. ###### Lemma 4. The inner product of any two harmonic $`k`$-forms is a constant function. In particular, the length of any harmonic form is constant. ###### Proof. That the length of any harmonic form is constant follows from equation (1). The more general statement follows by polarisation. ∎ ###### Lemma 5. Suppose $`\alpha _1,\mathrm{},\alpha _m`$ are orthogonal harmonic $`k`$-forms. Then $$\alpha =\underset{i=1}{\overset{m}{}}f_i\alpha _i$$ is harmonic if and only if the functions $`f_i`$ are all constant. ###### Proof. If $`\alpha `$ is harmonic, then $`g(\alpha ,\alpha _i)=f_i\|\alpha _i\|^2`$ is constant by Lemma 4. Using that the length of $`\alpha _i`$ is also constant by Lemma 4, we conclude that $`f_i`$ is constant. The converse is trivial. ∎ Lemma 4 implies that harmonic forms which are linearly independent at some point are linearly independent everywhere. Systems of linearly independent harmonic forms can be orthonormalised using constant coefficients. We can now generalise the discussion in Section 2 to higher dimensions. ###### Theorem 6. Suppose the closed oriented manifold $`M^n`$ is geometrically formal. Then 1. the real Betti numbers of $`M`$ are bounded by $`b_k(M)b_k(T^n)`$, and 2. if $`n=4m`$, then $`b_{2m}^\pm (M)b_{2m}^\pm (T^n)`$. 3. The first Betti number $`b_1(M)n1`$. ###### Proof. Fix a formal Riemannian metric on $`M`$. It follows from the above Lemmas that the number of linearly independent harmonic $`k`$-forms is at most the rank of the vector bundle $`\mathrm{\Lambda }^k`$. Similarly, when the dimension is $`4m`$, the number of self-dual or anti-self-dual harmonic forms in the middle dimension is bounded by the rank of $`\mathrm{\Lambda }_\pm ^{2m}`$. Suppose now that $`\alpha _1,\mathrm{},\alpha _{n1}`$ are linearly independent harmonic $`1`$-forms. Then $`(\alpha _1\mathrm{}\alpha _{n1})`$ is also a harmonic $`1`$-form, and is linearly independent of $`\alpha _1,\mathrm{},\alpha _{n1}`$. Thus $`b_1(M)n1`$ implies $`b_1(M)=n`$. ∎ There is an uncanny similarity here with the classification of flat Riemannian manifolds , which satisfy all the conclusions of Theorem 6. We can push this further: ###### Theorem 7. Suppose the closed oriented manifold $`M^n`$ is geometrically formal. If $`b_1(M)=k`$, then there is a smooth submersion $`\pi :MT^k`$, for which $`\pi ^{}`$ is an injection of cohomology algebras. In particular, if $`b_1(M)=n`$, then $`M`$ is diffeomorphic to $`T^n`$. In this case every formal Riemannian metric is flat. ###### Proof. Fix a formal Riemannian metric $`g`$ on $`M`$. We consider the Albanese or Jacobi map $`\pi :MT^k`$ given by integration of harmonic $`1`$-forms. As the harmonic $`1`$-forms have constant lengths and inner products, $`\pi `$ is a submersion. It induces an isomorphism on $`H^1`$, and products of linearly independent harmonic $`1`$-forms are never zero, but are harmonic because the metric is formal. In the case $`b_1(M)=n`$, we conclude that $`M`$ is a covering of $`T^n`$, and is therefore a torus itself. Every formal metric on $`T^n`$ must be flat because it admits an orthonormal framing by harmonic $`1`$-forms. ∎ In the case $`b_1(M)=1`$ we have a partial converse to Theorem 7: ###### Theorem 8. Let $`M`$ be a closed oriented $`n`$-manifold which fibers smoothly over $`S^1`$. If $`b_1(M)=1`$ and $`b_k(M)=0`$ for $`1<k<n1`$, then $`M`$ is geometrically formal. ###### Proof. Suppose $`M`$ fibers over $`S^1`$ with fiber $`F`$ and monodromy diffeomorphism $`\varphi :FF`$. By Moser’s Lemma, we may assume that $`\varphi `$ preserves a volume form $`ϵ`$ on $`F`$, so that its pullback to $`F\times `$ descends to $`M`$ as a closed form which is a volume form along the fibers. We can find a Riemannian metric on $`M`$ for which $`ϵ=\alpha `$ is the closed $`1`$-form defining the fibration over $`S^1`$, and has constant length. Then $`\alpha `$ and $`ϵ`$ generate the harmonic forms in degree $`1`$ and $`n1`$, and their product is harmonic. ∎ Flat manifolds satisfy further topological restrictions, for example their Euler characteristics vanish. In our present context we have: ###### Theorem 9. Suppose the closed oriented manifold $`M^n`$ is geometrically formal. 1. If $`b_k(M)0`$, then $`e(\mathrm{\Lambda }^k)=0`$, and 2. if $`n=4m`$ and $`b_{2m}^\pm (M)0`$, then $`e(\mathrm{\Lambda }_\pm ^{2m})=0`$. In particular the Euler characteristic of $`M`$ vanishes if $`b_1(M)0`$. ###### Proof. This follows from the obstruction-theory definition of the Euler class, and the fact that every nontrivial harmonic form has no zeros because of (1). ∎ ## 4. Dimension three If $`M`$ is a closed oriented geometrically formal $`3`$-manifold, then by Theorem 6 we have $`b_1(M)\{0,1,3\}`$. If the first Betti number is maximal, then Theorem 7 says that $`M`$ is the 3-torus. At the other extreme, if the first Betti number is zero, then $`M`$ is a rational homology sphere. Clearly every metric on every such manifold is formal. Thus, the only interesting case is that of first Betti number one. Then the real cohomology algebra is that of the globally symmetric space $`S^2\times S^1`$, so that Theorem 2 is proved in the $`3`$-dimensional case. Theorems 7 and 8 imply: ###### Corollary 10. Let $`M`$ be a closed oriented $`3`$-manifold with $`b_1(M)=1`$. Then $`M`$ is geometrically formal if and only if it fibers over $`S^1`$. This includes many non-symmetric manifolds. Thurston has proved that every $`3`$-manifold which fibers over $`S^1`$ carries a unique locally homogeneous geometry. It is not clear whether the induced metric is formal, even when the first Betti number is one. ## 5. Dimension four If $`M`$ is a closed oriented geometrically formal $`4`$-manifold, then by Theorem 6 we have $`b_1(M)\{0,1,2,4\}`$. If the first Betti number is maximal, then Theorem 7 says that $`M`$ is the 4-torus. ### 5.1. First Betti number $`=2`$. In this case the Euler characteristic vanishes by Theorem 9, and $`b_2(M)=2`$. Fix a formal Riemannian metric $`g`$. If $`\alpha `$ and $`\beta `$ are harmonic $`1`$-forms generating $`H^1(M)`$, then they are pointwise linearly independent. Therefore $`\omega =\alpha \beta `$ is a non-zero harmonic $`2`$-form with square zero. Thus the intersection form of $`M`$ is indefinite, and we conclude $`b_2^+=b_2^{}=1`$. This means that the real cohomology ring of $`M`$ is the same as that of the globally symmetric space $`S^2\times T^2`$. We know from Theorem 7 that $`M`$ fibers over $`T^2`$. The above discussion shows that the fiber is nontrivial in homology. There are many examples of such manifolds, other than $`S^2\times T^2`$. If $`N`$ is any $`3`$-manifold with $`b_1(N)=1`$ which fibers over the circle, then the product $`M=N\times S^1`$ is a $`4`$-manifold with the real cohomology ring of $`S^2\times T^2`$. By Corollary 10 it is geometrically formal, because we can take a product metric which on $`N`$ is the formal metric constructed in the proof of Theorem 8. ### 5.2. First Betti number $`=1`$. If the first Betti number is one, the Euler characteristic vanishes by Theorem 9, and therefore $`b_2(M)=0`$. In this case $`M`$ has the real cohomology algebra of the globally symmetric space $`S^3\times S^1`$. Theorems 7 and 8 imply: ###### Corollary 11. Let $`M`$ be a closed oriented $`4`$-manifold with $`b_1(M)=1`$ and $`b_2(M)=0`$. Then $`M`$ is geometrically formal if and only if it fibers over $`S^1`$. This includes many non-symmetric manifolds. The simplest example is a product of $`S^1`$ with a rational homology $`3`$-sphere which is not symmetric. ### 5.3. First Betti number $`=0`$. From Theorem 6 we know $`b_2^\pm 3`$. If $`b_2^+>0`$, then there are nontrivial self-dual harmonic forms. By (1) they have no zeros and so define almost complex structures compatible with the orientation of $`M`$. It follows that $`b_2^+`$ is odd. Similarly, if $`b_2^{}>0`$, then there are almost complex structures compatible with the orientation of $`\overline{M}`$ and $`b_2^{}`$ must be odd. Suppose now that $`b_2^+(M)=3`$. Then the self-dual harmonic forms trivialise $`\mathrm{\Lambda }_+^2`$, and each defines an almost complex structure with trivial first Chern class (because the pointwise orthogonal complement of each in $`\mathrm{\Lambda }_+^2`$ is trivial). Thus $`0=c_1^2(M)=2\chi (M)+3\sigma (M)=4+5b_2^+b_2^{}=19b_2^{}`$, which contradicts $`b_2^{}3`$. Therefore $`b_2^+=3`$ is not possible, and similarly $`b_2^{}=3`$ is not possible either. Thus the only possible values for $`b_2^\pm `$ are $`0`$ and $`1`$, and all combinations occur for the globally symmetric spaces $`S^4`$, $`P^2`$, $`\overline{P^2}`$ and $`S^2\times S^2`$. This finally completes the proof of Theorem 2 in the $`4`$-dimensional case. Any other example must have the same real cohomology ring as one of the above. In fact, other examples exist for each cohomological type. In the case of $`S^4`$ any rational homology $`4`$-sphere will do. In the case of $`P^2`$, there is the Mumford surface , an algebraic surface of the form $`H^2/\mathrm{\Gamma }`$ with the same rational cohomology as $`P^2`$. The Kähler form is of course harmonic, it generates the cohomology and its square is harmonic. Reversing the orientation of the Mumford surface we obtain an example with the cohomology ring of $`\overline{P^2}`$. Finally, in the case of $`S^2\times S^2`$, there is also a locally Hermitian symmetric algebraic surface $`M`$ of general type with the same real cohomology ring, due to Kuga, cf. p. 237. In this case $`M`$ is of the form $`(^2\times ^2)/\mathrm{\Gamma }`$, and the harmonic forms are generated by a self-dual and an anti-self-dual harmonic $`2`$-form (with respect to the locally symmetric metric). These are Kähler forms for $`M`$ and $`\overline{M}`$ respectively, and are therefore parallel and have harmonic products. ## 6. Obstructions from symplectic geometry In this section we discuss relations between harmonicity of products of harmonic forms on $`4`$-manifolds on the one hand, and existence of symplectic structures on the other. This leads to further obstructions to geometric formality. Let $`M`$ denote a closed oriented $`4`$-manifold with a Riemannian metric $`g`$. Suppose that $`b_2^+(M)>0`$. Then there is a nontrivial $`g`$-self-dual harmonic $`2`$-form $`\omega `$. If the product $$\omega \omega =\omega \omega =\|\omega \|^2dvol_g$$ is harmonic, then $`\omega `$ has constant length, and in particular has no zeros. It is then a symplectic form on $`M`$ compatible with the orientation, and $`g`$ is an almost Kähler metric. There are $`4`$-manifolds for which the elementary obstructions of Section 3 vanish, but which are not geometrically formal because they do not admit any symplectic structure: ###### Example 12. Let $`X`$ be $`P^2`$, $`S^2\times S^2`$, or the Kuga or Mumford surface. Let $`N`$ be a rational homology $`4`$-sphere whose fundamental group has a nontrivial finite quotient. Then $`M=X\mathrm{\#}N`$ has the real cohomology ring of the geometrically formal manifold $`X`$, but is not itself geometrically formal because it does not admit any symplectic structure by the result of . There is another application of the relationship between harmonicity of products of harmonic forms and symplectic structures. Namely we can show that on certain manifolds all products of certain harmonic forms are non-harmonic. This is obviously much stronger than geometric non-formality. For an example, consider the smooth manifold $`M`$ underlying a complex K3 surface. Then $`M`$ is simply connected with $`b_2^+=3`$ and $`b_2^{}=19`$. The elementary considerations in Section 5 already show that $`M`$ is not geometrically formal. We can sharpen this as follows: ###### Proposition 13. Let $`g`$ be an arbitrary Riemannian metric on the K3 surface $`M`$. If $`\alpha `$ is a $`g`$-anti-self-dual harmonic $`2`$-form, then it must have a zero. In particular the wedge product $`\alpha \beta `$ is not harmonic for any $`\beta `$ unless it vanishes identically. For example, if $`\alpha `$ is nontrivial then $`\alpha \alpha `$ is not harmonic. ###### Proof. Suppose $`\alpha `$ is nontrivial and anti-self-dual. We have $$\alpha \alpha =\alpha \alpha =\|\alpha \|^2dvol_g,$$ which is harmonic if and only if the norm of $`\alpha `$ is constant. If it is constant, it must be a non-zero constant, and then the above equation shows that $`\alpha `$ is a symplectic form inducing the opposite (non-complex) orientation on $`M`$. In particular, $`\overline{M}`$ must have non-trivial Seiberg-Witten invariants, see . But the K3 surface contains smoothly embedded (-2)-spheres, which become (+2)-spheres when the orientation is reversed, showing that all the Seiberg-Witten invariants vanish, see . ∎ ###### Remark 14. The vanishing of the Seiberg-Witten invariants for $`\overline{M}`$ can also be proved without appealing to the vanishing theorem for spheres of positive self-intersection. For a scalar-flat Calabi-Yau metric the Seiberg-Witten equations on $`\overline{M}`$ have no solution, though they do have a (unique) solution on $`M`$. This can be generalised quite substantially. If $`M`$ has an indefinite intersection form, there are both self-dual and anti-self-dual harmonic forms for all metrics. If the square of such a form is harmonic, it is a symplectic form on $`M`$, respectively $`\overline{M}`$. But by the results of , manifolds which are symplectic for both choices of orientation are quite rare. Thus, the above proof generalises to many cases to show that for all metrics on certain $`4`$-manifolds, all self-dual and/or all anti-self-dual harmonic forms must have zeros and non-harmonic squares. This generalises the existence of zeros of harmonic $`1`$-forms on surfaces, cf. Section 2. In the case of complex surfaces, Theorem 1 of implies the following: ###### Theorem 15. Let $`M`$ be a compact complex surface of general type. Assume one of the following conditions holds: 1. $`K_M`$ is not ample, or 2. $`c_1^2(M)`$ is odd, or 3. the signature $`\sigma (M)`$ is negative, or is zero and $`M`$ is not uniformised by the polydisk. Then for every Riemannian metric $`g`$ on $`M`$, all $`g`$-anti-self-dual harmonic $`2`$-forms have zeros and non-harmonic squares. In the first case, the argument is the same as in the proof of Proposition 13, because ampleness of $`K_M`$ only fails if there are rational curves of negative self-intersection. Note that we only need smoothly rather than holomorphically embedded spheres, so one can replace condition 1. by a weaker assumption. ## 7. General existence of non-formal metrics Having seen that only very few manifolds are geometrically formal, we now want to show that even these tend to also have non-formal metrics. The two-dimensional case of the following result was already proved in Section 2. ###### Theorem 16. A closed oriented manifold admits a non-formal Riemannian metric if and only if it is not a rational homology sphere. ###### Proof. It is clear that every metric on every rational homology sphere is formal because there are no nontrivial harmonic forms. Conversely, assume that $`M`$ is a manifold with a non-zero Betti number $`b_k(M)`$, for $`0<k<dim(M)`$. Let $`g`$ be a Riemannian metric which has positive curvature operator on an open set, say a ball $`BM`$, and assume it is formal. If $`\alpha `$ is a nontrivial $`g`$-harmonic $`k`$-form, then $$\alpha \alpha =\|\alpha \|^2dvol_g$$ shows that $`\alpha `$ has constant length. Therefore, the Bochner-Weitzenböck formula (3) $$\mathrm{\Delta }(\alpha )=^{}\alpha +(\alpha )$$ for $`k`$-forms allows us to compute: $$0=\mathrm{\Delta }(\frac{1}{2}\|\alpha \|^2)=g(^{}\alpha ,\alpha )\|\alpha \|^2=g((\alpha ),\alpha )\|\alpha \|^2.$$ Here the term $``$ is positive on $`B`$, because there the curvature operator is positive. Thus $`\alpha `$ vanishes identically on $`B`$. As $`\alpha `$ is harmonic, the unique continuation principle implies that $`\alpha `$ vanishes on all of $`M`$, contradicting the assumption that $`\alpha `$ is nontrivial. ∎ ###### Remark 17. The above proof shows that there is an open set of non-formal metrics in the space of all Riemannian metrics (with the $`C^{\mathrm{}}`$ topology, say) on any manifold which is not a rational homology sphere.
warning/0004/gr-qc0004053.html	ar5iv	text	# Response of massive bodies to gravitational waves ## 1 Introduction Gravitational waves were already considered by Einstein as the wave solutions of the linearized field equations of gravity. There is indirect evidence of their existence through systems of binary pulsars that loose energy in form of gravitational radiation , their direct experimental measurement presently is one of the most challenging tasks in gravitational physics. Very sensitive detectors operating at the quantum limit are needed to detect directly gravitational waves from cosmic events such as collapsing or colliding star systems. There are basically two different types of detectors: resonant mass antennas based on the resonant excitation of quadrupole-type modes of a appropriately chosen massive body, like the bar detectors conceived by Weber , and laser interferometric devices that detect the direction-dependent variation of the proper distance between the mirrors of a Michelson interferometer . Detectors of both types are presently under construction . Commonly a resonant mass antenna is described in Riemannian normal coordinates with respect to its center of mass, the proper frame of reference (PFR). The detector is analyzed in term of normal modes, idealized by a spring that couples two masses , the resonant energy input is calculated. The intention of the work presented in this article was twofold: First, to validate the results obtained from the normal mode model by microscopic considerations, second to give a complementary description of the detector in the reference system of the wave, in which the linearized solutions of the Einstein field equations are computed. Our model is based on the local properties of a detector in terms of the fundamental binding forces, electromagnetic and other, which we consider in both the PFR and the wave system. Analyzing the deviations from the tidal motion, we find that the energy input from gravitational waves on an electromagnetically coupled massive body is restricted to the surface of the body, whereas gravitational coupling leads to true bulk excitation of quadrupole modes. This result does not contradict the normal mode picture at all, rather it presents a complementary viewpoint that has eluded the normal mode analysis. The reason is that though the energy input into any single mode is nonlocal, in the special case of gravitational waves the superposition of all excited modes describes a localized excitation. Based on our observations we propose a new, microscopic type of detector. ## 2 Frames of reference There exist two frames of reference that can be used for the analysis of gravitational waves and gravitational wave detectors. The natural system to study the waves is a perturbed Minkowski system for which the linearized Einstein field equations are solved. In this reference system the plane wave is gauged, conventionally a transverse-traceless (TT) gauge is chosen . On the other hand, the use of the PFR system with Riemannian normal coordinates with respect to the center of mass is the natural system for the study of a detector. For our purpose to study of the detector on the microscopic level, both system have advantages and disadvantages. In the PFR system we have common argument that in the metric of the gravitational wave field leads to variations of the electromagnetic field of order $$\delta A/A\left(L^2/\lambda ^2\right)h^{TT}$$ (1) where $`L`$ is the distance from the origin, $`\lambda `$ and $`h^{TT}`$ the wave length and amplitude of the gravitational wave, respectively, whereas the tidal gravitational forces lead to displacements of the constituent particles of the body that in turn lead to variations of the electromagnetic field of order $$\delta A/Ah^{TT}$$ (2) so that the metric effect can be neglected. This means that in the PFR system the unperturbed solutions of the Maxwell equations can be used. On the other hand, in the PFR system we have to deal with the general problem of general relativity that local energy densities have no invariant meaning, so that we cannot easily control the exchange of energy between the wave and the detector. Here the TT system has the advantage that we can employ a Hamiltonian description of the constituent particles, and because the system is in principle a special relativistic one, we have the standard laws of energy-momentum conservation. Therefore the TT system proves useful for our considerations. Because considerations in both coordinate system contribute to the understanding of the detector response, we formulate our results in both systems. We first show that in the TT system no energy is transferred to a system of non-interacting particles; in the PFR system this corresponds to the well-known quadrupole type oscillations. Next we extend the discussion to electromagnetically coupled systems, where we can use the standard Coulomb force in the PFR system, whereas we have to solve the metric Maxwell equations in the TT system. Finally we discuss a gravitationally coupled system and show that there exists a decisive difference in the response of the detector which is due to the different nature of the fundamental forces. ## 3 Particle Motion A point-like test mass with electric charge $`e`$ can be described in the presence of electromagnetic fields and gravitation by the Hamiltonian $$H=c\sqrt{\left(p_\mu eA_\mu \right)g^{\mu \nu }\left(p_\nu eA_\nu \right)}$$ (3) where we use coordinates $`x^\mu ,\mu =0,1,2,3,4`$ for space-time with metric $`g_{\mu \nu }(x^\kappa )`$ and signature $`+`$, $`p_\mu `$ are the momentum coordinates in the cotangent space, and $`A_\mu (x^\kappa )`$ is the electromagnetic four-potential. We also use the notation $`(ct,x,y,z)`$ in an obvious manner. The evolution parameter will be denoted by $`\tau `$. The Hamiltonian is conserved, $`H/\tau =0`$, it represents the rest mass $`m=H/c^2`$of the particle. The canonical equations of motion are given by $$\dot{x}^\mu =\frac{H}{p_\mu },\dot{p}_\mu =\frac{H}{x^\mu }.$$ (4) The constancy of $`H`$ is equivalent to $`\dot{x}^\mu g_{\mu \nu }\dot{x}^\nu =c^2`$ for any trajectory, thus the evolution parameter is the proper time. A gravitational wave propagating in $`z`$-direction with $`+`$ and $`\times `$ polarization modes is described in TT gauge by the metric tensor $$g_{\mu \nu }=\left(\begin{array}{cccc}1\hfill & 0\hfill & 0\hfill & 0\hfill \\ 0\hfill & 1+f_+(ctz)\hfill & f_\times (ctz)\hfill & 0\hfill \\ 0\hfill & f_\times (ctz)\hfill & 1f_+(ctz)\hfill & 0\hfill \\ 0\hfill & 0\hfill & 0\hfill & 1\hfill \end{array}\right)=\eta _{\mu \nu }+h_{\mu \nu }$$ (5) where $`h_{\mu \nu }`$ is a small perturbation of the Minkowski metric $`\eta _{\mu \nu }`$ . This perturbation acts as classical, special relativistic field, so that the system can be treated within special relativity, except that the interpretation of the energy must be treated with care . In the absence of an electromagnetic field we obtain for this metric a Hamiltonian that leads to four conserved quantities: $`H,p_x,p_y`$, and $`p_0+p_3=E/c+p_z`$ where $`E`$ is the energy of the particle in the sense of special relativity. Note that $`p_3=m\dot{z}`$, so that the difference between the energy and the conventional $`z`$-momentum is conserved. This is natural, since the gravitational wave not only carries energy $`E_w`$, but also momentum $`P_w`$ with the relation $`E_w=c\left\|P_w\right\|`$ that holds for all massless objects in special relativity. The exchange $`\mathrm{\Delta }E`$ of energy between the wave and a test mass thus is always accompanied with an exchange $`\mathrm{\Delta }P`$ of momentum: $$\mathrm{\Delta }E=c\mathrm{\Delta }P.$$ (6) The existence of four conserved quantities now allows us to integrate the equations of motion completely: $`\begin{array}{cc}\dot{x}^0=\frac{1}{m}p_0\hfill & \dot{x}^2=\frac{1}{m}\left(g^{21}p_1+g^{22}p_2\right)\hfill \\ \dot{x}^1=\frac{1}{m}\left(g^{11}p_1+g^{12}p_2\right)\hfill & \dot{x}^3=\frac{1}{m}p_3\hfill \\ p_1=const\hfill & p_2=const\hfill \end{array}`$ (10) $`p_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(p_0+p_3\right)+{\displaystyle \frac{1}{2}}{\displaystyle \frac{m^2}{\left(p_0+p_3\right)}}\left(c^2{\displaystyle \frac{1}{m^2}}p_ag^{ab}p_b\right)`$ $`p_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(p_0+p_3\right){\displaystyle \frac{1}{2}}{\displaystyle \frac{m^2}{\left(p_0+p_3\right)}}\left(c^2{\displaystyle \frac{1}{m^2}}p_ag^{ab}p_b\right)`$ (11) where we use the indices $`a,b`$ for a summation over $`1,2`$ only. Let us consider a wave pulse. We denote the initial conditions before the arrival of the pulse by $`\overline{p}_\mu `$, and have $`p_0(\tau )`$ $`=`$ $`\overline{p}_0{\displaystyle \frac{1}{2}}{\displaystyle \frac{p_ah^{ab}p_b}{\overline{p}_0+\overline{p}_3}}`$ $`p_3(\tau )`$ $`=`$ $`\overline{p}_3+{\displaystyle \frac{1}{2}}{\displaystyle \frac{p_ah^{ab}p_b}{\overline{p}_0+\overline{p}_3}}`$ (12) where $`h^{ab}=g^{ab}\eta ^{ab}`$. Thus after the pulse, where the perturbation $`h`$ is zero again, the particle has the same four-momentum as before, and the only possible effect is a displacement of the straight trajectory after the pulse from the one before the pulse. If a particle is initially at rest, it stays at rest in this reference frame. Two particles that are at rest relative to each other, remain at rest relative to each other, though the proper distance between them changes with the wave amplitude. Thus free test particles do not take up energy from a gravitational wave. When we restrict this statement to a comparison of the energy before and after a wave pulse or wave train, we are moreover free from the ambiguity of the energy definition in general relativity. We now look at this result in the PFR system. In the PFR system particles do not stay at rest in the wave field, but move under tidal accelerations. In PFR coordinates $`(\widehat{t},\widehat{x},\widehat{y},\widehat{z})`$ the tidal accelerations produced by the wave are : $$\frac{d^2\widehat{x}^a}{d\widehat{t}^2}=\widehat{R}_{a0b0}\widehat{x}^b=\frac{1}{2}\frac{d^2h_{ab}}{d\widehat{t}^2}\widehat{x}^b$$ (13) For a particle that is initially at rest at $`\widehat{x}^{a(0)}`$ before a wave pulse arrives, the solution is simply given by $$\widehat{x}^a\left(\widehat{t}\right)=\widehat{x}^{a(0)}+\frac{1}{2}h_{ab}\left(\widehat{t}\right)\widehat{x}^{b(0)}.$$ (14) (Note that $`\widehat{R}_{a0b0}`$ and $`\widehat{R}{}_{}{}^{a}_{0b0}`$ differ only by terms of order $`O\left(h^2\right)`$). This is the well-known quadrupole-like tidal motion of a system of particles around the origin of the frame of reference. The coordinate distance of the particles varies according to the changing metric distance between them. We can ascribe a standard kinetic energy to this motion, but there will be additional contributions to the conserved energy from the metric. Our analysis in the TT system has shown that the energy of the wave does not vary in this case. The coordinate transformation from the TT to the PFR system is time-dependent, so that the energy conservation in the TT system translates into a more complicated conservation law in the PFR system. Because free particles do not effectively take up energy from a gravitational wave, we have to take the coupling between particles into account in order to describe the response of a detector. The coupling is basically of electromagnetic nature in small massive bodies, but may also be of gravitational nature in large bodies. Because the effects of the gravitational wave are so small, massive detectors must be cooled to zero temperature as near as possible. Thus the ground state of a body where the constituent atoms or ions are at rest relative to each other except for quantum effects will serve as an appropriate model. Rotational motion with respect to the TT frame of reference must be taken into account. Internal motion, thermal or otherwise, of the atoms leads to forces of order $`m\mathrm{\Delta }vh/t`$; we do not consider this kind of phonon-graviton interaction in this article, because it is the idea of the Weber detector to excite the quadrupole modes, not to enhance already excited modes. ## 4 Electromagnetic Field We first look at the Coulomb potential generated by a charge that resides in the field of the wave. In the PFR system we have the above cite argument that the metric perturbations of the electromagnetic potentials can be ignored. Thus, for the description of an electromagnetically bound solid body, we have to take only the Coulomb forces into account, and can ignore contributions from magnetic fields and electromagnetic radiation. In the PFR system the coordinate distance agrees with the metric distance up to order $`O(h)`$. Therefore we can state in an invariant manner, that the Coulomb potential depends only on the metric distance of the particles. We now show how this translates into the TT system. We assume that the wave field is slowly varying and the velocity of the charge relative to the source of the fields is so small that magnetic fields can be ignored. The corresponding equation to solve for the electromagnetic potential generated by a charge $`q_s`$ at rest at $`x=y=z=0`$ is $$_\mu \sqrt{g}g^{\mu \nu }_\nu A^0=\frac{q_s}{\epsilon _0}\delta ^3(x,y,z).$$ (15) Our considerations in the PFR system give rise to the following ansatz: $$A_s^0(\stackrel{}{r},ctz)=\frac{q_s}{4\pi \epsilon _0r},\stackrel{}{r}=\left(x^i\right)_{i=1,2,3},r^2=x^ig_{ij}\left(ctz\right)x^j.$$ (16) where the Euclidean coordinate distance from the source is replaced by the metric distance. We verify that equation (15) is satisfied to order $`O(h)`$. For simplicity only we consider a wave with the $`+`$-mode only where $$r^2=x^2(1f_+\left(ctz\right))+y^2\left(1+f_+\left(ctz\right)\right)+z^2.$$ (17) Using $`\sqrt{g}=1+O\left(h^2\right)`$, leaving out terms of order $`h^2`$ and higher, we obtain : $`_\mu \sqrt{g}g^{\mu \nu }_\nu {\displaystyle \frac{1}{r}}`$ $``$ $`\left(c^2_t^2{\displaystyle \frac{1}{1f_+}}_x^2{\displaystyle \frac{1}{1+f_+}}_y^2_z^2\right){\displaystyle \frac{1}{r}}`$ (18) $`=`$ $`c^1_t\left({\displaystyle \frac{x^2y^2}{2r^3}}f_+^{}\right)+_x{\displaystyle \frac{x}{r^3}}+_y{\displaystyle \frac{y}{r^3}}_z\left({\displaystyle \frac{z}{r^3}}{\displaystyle \frac{x^2y^2}{2r^3}}f_+^{}\right)`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{r}}{r^3}}+{\displaystyle \frac{3\left(x^2y^2\right)^2}{4r^5}}\left(f_+^{}\right)^2+{\displaystyle \frac{x^2y^2}{2r^3}}f_+^{\prime \prime }{\displaystyle \frac{x^2y^2}{2r^3}}f_+^{\prime \prime }+f_+^{}_z{\displaystyle \frac{y^2x^2}{2r^3}}`$ $`=`$ $`4\pi \delta \left(r\right)+{\displaystyle \frac{3\left(y^2x^2\right)^2}{4r^5}}\left(f_+^{}\right)^2+f_+^{}_z{\displaystyle \frac{y^2x^2}{2r^3}}`$ Now the first of these two remaining terms is of order $`O\left(h^2\right)`$ and can thus be ignored, the second is proportional to the spatial derivative of the wave field, which can be ignored for detectors that are small against the wave length . Thus the ansatz (16) solves equation (15) except for terms of order $`O\left(h^2\right)`$ or proportional to $`_zh`$. In the same sense the Lorentz gauge holds. The result thus agrees with that in the PFR system: The Coulomb potential depends only on the metric distance. In the following we assume that this principle may also by extended to other, phenomenological potentials, because the time scale set by the gravitational waves is by far larger than that of the induced changes of all other fundamental interactions. So in the TT system the Coulomb force between two charged particles varies in phase with the gravitational wave. It is not hard to see that the local energy density of the electromagnetic field, though time-dependent, is only relocated, so that the integrated energy density does not change to first order in $`h`$, implying that radiation effects are at most of second order in $`h`$, in agreement with the situation in the PFR system. On the other hand accelerated, moving charges produce magnetic fields and electromagnetic radiation, but both effects can be ignored for the analysis of weak mode excitation in a solid body. ## 5 Many particles The main point of this work is to show that a detailed microscopic, many-particle model of the detector exhibits two aspects that elude the spring model and the normal mode analysis. The first aspect is that the input of energy from the wave to the detector occurs only at the surface of the detector. This effect can already be modelled with a linear chain , as done in the next subsection. The second aspect is the dependence on the nature of the fundamental forces that stabilize a body: electromagnetic and gravitational binding forces lead to different mode structure of the responde, as shown subsequently. ### 5.1 Linear chain model In the PFR system $`N`$ classical particles with coordinates $`\stackrel{}{r}^{\left(s\right)}=(\widehat{x},\widehat{y},\widehat{z}),\alpha =1,\mathrm{},N`$ obey the equations of motion $$m_i\frac{d^2}{d\widehat{t}^2}\stackrel{}{r}^{\left(s\right)}=\underset{s^{}s}{}\stackrel{}{F}_{s^{}s}(\stackrel{}{r}^{\left(s^{}\right)},\stackrel{}{r}^{\left(s\right)})+\frac{m_i}{2}\frac{d^2h}{d\widehat{t}^2}\stackrel{}{r}^{\left(s\right)}$$ (19) where for simplicity the dot denotes the multiplication of the $`3\times 3`$ matrix $`h_{ij}`$ with the vector $`\stackrel{}{r}^{\left(s\right)}`$, and $`F_{s^{}s}`$ represents the fundamental forces between two particles. Let us consider the special model of a linear chain in $`\widehat{x}`$-direction with $`N=2n`$ equal particles numbered by $`s^{},s=n\mathrm{}n,`$ coupled by nearest-neighbor forces $$F_{s^{}s}(\widehat{x}^{\left(s^{}\right)},\widehat{x}^{\left(s\right)})=\omega _0^2\left(\widehat{x}^{\left(s^{}\right)}\widehat{x}^{\left(s\right)}l_0\right)\text{ for }\left\|s^{}s\right\|=1$$ (20) and a wave with $`+`$-mode only, so that the equations of motion are $$\frac{d^2}{d\widehat{t}^2}\widehat{x}^{\left(s\right)}=k\underset{\beta =\alpha \pm 1}{}\left(\widehat{x}^{\left(s\right)}\widehat{x}^{\left(s^{}\right)}\pm l_0\right)+\frac{1}{2}\frac{d^2f_+}{d\widehat{t}^2}\widehat{x}^{\left(s\right)}$$ (21) where $`k=\omega _0^2/m`$. Without the tidal acceleration, the equilibrium positions are $`\widehat{X}^{\left(s\right)}=sl_0`$, relative to the center of mass that is identical with the origin of the coordinate system. We are now interested in the deviation from the tidal motion (LABEL:prf2) because the tidal motion itself will be in general too small to be observed itself. We set $$\xi ^{\left(s\right)}=\widehat{x}^{\left(s\right)}\widehat{X}^{\left(s\right)}\left(1+\frac{f_+}{2}\right)$$ (22) leading to $`\ddot{\xi }^{\left(s\right)}`$ $`=`$ $`k{\displaystyle \underset{s^{}=s\pm 1}{}}\left(\xi ^{\left(s\right)}\xi ^{\left(s^{}\right)}\left(\widehat{X}^{\left(s\right)}\widehat{X}^{\left(s^{}\right)}\right)\left(1+{\displaystyle \frac{f_+}{2}}\right)\pm l_0\right)+{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2f_+}{d\widehat{t}^2}}\left(\widehat{x}^{\left(s\right)}\widehat{X}^{\left(s\right)}\right)`$ (23) $`=`$ $`k{\displaystyle \underset{s^{}=s\pm 1}{}}\left(\xi ^{\left(s\right)}\xi ^{\left(s^{}\right)}\pm l_0{\displaystyle \frac{f_+}{2}}\right)`$ where we have omitted the term $`\frac{1}{2}\frac{d^2f_+}{d\widehat{t}^2}\left(\widehat{x}^{\left(s\right)}\widehat{X}^{\left(s\right)}\right)`$ because we consider deviations from the tidal motion that are proportional to $`f_+`$ itself, as induced by the wave, and thus this term is of order $`O\left(h^2\right)`$. Now the following happens: the terms $`l_0\frac{f_+}{2}`$ in (23) cancel for all particles that have two neighbors, only for the particles at the ends of the chain they do not. The reason is that under the tidal motion the distances between pair of neighboring particles remains constant, thus the induced additional forces from the left and right cancel; the pattern of the tidal acceleration comes very close to a null mode. So we can rewrite $$\ddot{\xi }^{\left(s\right)}=k\underset{s^{}=s\pm 1}{}\left(\xi ^{\left(s\right)}\xi ^{\left(s\right)}\right)kl_0\frac{f_+}{2}\left(\delta _{s,n}\delta _{s,n}\right).$$ (24) These equations have a clear interpretation: The deviation from the tidal motion is driven by an effective force that applies to the ends of the chain only. Or, we can state that it suffices to apply additional forces $`kl_0f_+/2`$ to the ends of the chain in order to suppress any deviations from the tidal motion. In this case the tension of the chain varies uniformly along with the wave strength, but no work is done at all against the internal coupling forces. Regarding the energy we have to be careful to use standard expressions, because there are metric contributions, as we have seen above, but we can certainly state that apart from the tidal motion energy is transferred to the chain only locally at the ends. For resonance cross terms in the kinetic energy between the tidal motion and the deviations from it play no role, so that $`m\dot{\xi }^2/2`$ describes the kinetic energy of interest. Nevertheless this local picture does not contradict the normal-mode analysis in any way. We still have the possibility to decompose the perturbing force into normal modes and derive equations for the driving of these modes in the standard way. The fundamental mode of the chain is driven strongest, but all other symmetric modes are also excited. It is the superposition of all these modes that gives rise to the local excitation of the chain. Only refining the spring model to a linear chain model could exhibit this property. Because we see from (24) that the effective force depends on the microscopic property of the chain in form of the equilibrium distance $`l_0`$, it is necessary to analyze a detector on the basis of a more realistic microscopic model that yield information on the strength of the driving forces. It turns out that this is also necessary because the difference between gravitational and electromagnetic coupling can only then be shown. ### 5.2 Microscopic forces in the PFR system For a more realistic model we can use eq. (19) with the fundamental Coulomb, gravitational, and repulsive short-range forces inserted. We assume that all these forces depend only on the difference vectors $`\stackrel{}{r}^{\left(\beta \right)}\stackrel{}{r}^{\left(\alpha \right)}`$ and stable equilibrium positions $`\stackrel{}{R}^{\left(\alpha \right)}`$ exists. Again we look at the deviation from the tidal motion induced by these forces. The transformation $$\stackrel{}{r}^{\left(\alpha \right)}=\left(1+\frac{1}{2}h\right)\stackrel{}{\rho }^{\left(\alpha \right)}$$ (25) now leads, assuming small deviations from equilibrium, $$\stackrel{}{\rho }^{\left(\alpha \right)}=\stackrel{}{R}^{\left(\alpha \right)}+O\left(h\right),$$ (26) to $`{\displaystyle \frac{d^2}{d\widehat{t}^2}}\stackrel{}{\rho }^{\left(\alpha \right)}`$ $`=`$ $`{\displaystyle \underset{\beta \alpha }{}}\stackrel{}{F}_{\beta \alpha }\left(\left(1+{\displaystyle \frac{1}{2}}h\right)\left(\stackrel{}{\rho }^{\left(\beta \right)}\stackrel{}{\rho }^{\left(\alpha \right)}\right)\right)+O\left(h^2\right)`$ $`=`$ $`{\displaystyle \underset{\beta \alpha }{}}\left[\stackrel{}{F}_{\beta \alpha }\left(\stackrel{}{\rho }^{\left(\beta \right)}\stackrel{}{\rho }^{\left(\alpha \right)}\right)+{\displaystyle \frac{1}{2}}\left(h\right)\stackrel{}{F}_{\beta \alpha }\left(\stackrel{}{\rho }^{\left(\beta \right)}\stackrel{}{\rho }^{\left(\alpha \right)}\right)\right]+O\left(h^2\right)`$ We will see that this equation for the deviation is identical with that derived for the motion in the TT system. The reason is that the transformation (LABEL:prf11) basically agrees with the coordinate transformation from the PFR to the TT system. For the distance we have $$\stackrel{}{r}^i\delta _{ij}\stackrel{}{r}^j=\stackrel{}{\rho }^i\left(\delta _{ik}+\frac{1}{2}h_{ik}\right)\left(\delta _{kj}+\frac{1}{2}h_{kj}\right)\stackrel{}{\rho }^j=\stackrel{}{\rho }^ig_{ij}^{TT}\stackrel{}{\rho }^j+O\left(h^2\right)$$ and thus all potential forces acting on the deviation of the tidal motion in the PFR system are identical with that in the TT system. ### 5.3 Hamiltonian analysis in the TT system We make a potential approximation for the many-particle Hamiltonian, since in a fully relativistic approach we had to include necessarily the dynamics of the electromagnetic field in order to preserve energy-momentum conservation. Thus we write the total Hamiltonian for many particles with coordinates $`(x^{(s)\mu },p_\mu ^{(s)})`$, $`s=1,2,\mathrm{}`$ as $$H=\underset{s}{}H^{(s)}$$ (28) where $$H^{(s)}=\frac{1}{2m_1}p_i^{(s)}g^{ij}\left(ctz^{(s)}\right)p_j^{(s)}+\underset{ss}{}\frac{1}{2}V_{ss}(\stackrel{}{r}^{(s)}\stackrel{}{r}^{(s)},ctz^{(s)})$$ (29) is the contribution from a single particle. $`V_{ss}`$ is the total potential generated by the particle $`s^{}`$, acting on particle $`s`$. As it should be, to each particle only half the potential energy is attributed. In the electromagnetic case, the other half, as well as the infinite self-energy is subtracted with the contribution from the electromagnetic field energy . We assume that the potential is a central potential depending only on the metric distance: $$V_{ss}(\stackrel{}{r}^{(s)}\stackrel{}{r}^{(s)},ctz^{(s)})=V_{ss}^0(\sqrt{\left(\stackrel{}{r}^{(s)}\stackrel{}{r}^{(s)}\right)^ig_{ij}\left(ctz^{(s)}\right)\left(\stackrel{}{r}^{(s)}\stackrel{}{r}^{(s)}\right)^j}),$$ as derived for the Coulomb potential. The evolution parameter is the time $`t`$, common to all particles. The Hamiltonian (28) conserves the difference between the total energy and the center of mass $`z`$-momentum: $$\frac{d}{dt}\left(H+c\underset{s}{}p_3^{(s)}\right)=0.$$ (30) But since we have less conservation laws than coordinates, energy and momentum transfer from the wave to the particle system has become possible. The change of the total energy is calculated using the equations of motion $`{\displaystyle \frac{d}{dt}}p_z^{(s)}`$ $`=`$ $`{\displaystyle \frac{H}{z^{(s)}}}`$ $`=`$ $`{\displaystyle \frac{1}{2m_s}}p_i^{(s)}h^{ij}\left(ctz^{(s)}\right)p_j^{(s)}`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{z^{(s)}}}\left({\displaystyle \underset{ss}{}}V_{ss}(\stackrel{}{r}^{(s)}\stackrel{}{r}^{(s)},ctz^{(s)})+{\displaystyle \underset{s^{}s}{}}V_{ss}(\stackrel{}{r}^{(s)}\stackrel{}{r}^{(s)},ctz^{(s)})\right)`$ that lead us to $`{\displaystyle \frac{d}{dt}}{\displaystyle \underset{s}{}}p_z^{(s)}`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle \frac{1}{2m_s}}p_a^{(s)}h^{ab}\left(ctz^{(s)}\right)p_b^{(s)}`$ (32) $`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,s^{}s}{}}_2V_{ss}(\stackrel{}{r}^{(s)}\stackrel{}{r}^{(s^{})},ctz^{(s)})`$ where $`h^{ij}`$ denotes the derivative, $`_2V_{ss}`$ the partial derivative with respect to the second argument only. The derivatives of $`V_{ss}`$ with respect to the first argument cancel in the sum because of the dependence on $`\stackrel{}{r}^{(s)}\stackrel{}{r}^{(s^{})}`$ only. To first order in the perturbation we can approximate $$_2V_{ss}(\stackrel{}{r}^{(s)}\stackrel{}{r}^{(s^{})},ctz^{(s)})V_{ss}^0\left(r_{ss}\right)\frac{x_{ss}^ax_{ss}^b}{2r_{ss}}h_{ab}^{}\left(ctz^{(s)}\right)$$ (33) where we use $`x_{ss}^a=\left(\stackrel{}{r}^{(s)}\stackrel{}{r}^{(s)}\right)^a`$ for short, and $`r_{ss}`$ is the unperturbed Euclidean distance. Thus we have (with $`h^{ab}=h_{ab}`$ in lowest order) $$\frac{d}{dt}\underset{s}{}p_z^{(s)}=\underset{s}{}h_{ab}^{}\left(ctz^{(s)}\right)\left[\frac{1}{2m_s}p_a^{(s)}p_b^{(s)}\underset{s^{}s}{}V_{ss}^0\left(r_{ss}\right)\frac{x_{ss}^ax_{ss}^b}{2r_{ss}}\right].$$ (34) For a pulse of duration $`T`$ that interacts with the particle system the change of the center of mass $`z`$-momentum is given by $$\mathrm{\Delta }P_z=\underset{s}{}_{t_0}^{t_0+T}h_{ab}^{}\left(ctz^{(s)}\right)q_{ab}^{(s)}(t)𝑑t$$ (35) where $`q_{ab}^{(s)}`$ is the microscopic contribution from a single particle: $$q_{ab}^{(s)}(t)=\frac{1}{2m_s}p_a^{(s)}p_b^{(s)}\underset{ss}{}V_{ss}^0\left(r_{ss}\right)\frac{x_{ss}^ax_{ss}^b}{2r_{ss}}.$$ (36) This contribution is quadrupole-like, but it is weighted with derivative of the potential. Then the sum or the integral over all particles leads in the case of the Coulomb potential to alternating sums where contributions cancel, similar to the summation leading to the Madelung constant. On the other hand, if the potential is purely attractive, as for gravitation, this effect does not occur, giving rise to a completely different response, as well will see. For a small detector we can assume that the wave field is constant over the particle system represented by the center-of-mass coordinate $`z(t)`$, so we can change the integration variable to $`\tau =tz(t)/c`$ and obtain $$\mathrm{\Delta }P_z=_{\mathrm{}}^+\mathrm{}𝑑\tau h_{ab}^{}\left(c\tau \right)\underset{s}{}\frac{q_{ab}^{(s)}(t(\tau ))}{1\dot{z}(t(\tau ))/c}=_{\mathrm{}}^+\mathrm{}𝑑\tau h_{ab}^{}\left(c\tau \right)Q_{ab}(\tau )$$ (37) with $$Q_{ab}(\tau )=\underset{s}{}\frac{q_{ab}^{(s)}(t(\tau ))}{1\dot{z}(t(\tau ))/c}.$$ (38) For a wave pulse or wave train with duration $`T`$, such that $`h_{ab}\left(0\right)=h_{ab}\left(cT\right)=0`$ we use partial integration to write (37) as $$\mathrm{\Delta }E=c\mathrm{\Delta }P_z=_{\mathrm{}}^+\mathrm{}𝑑\tau h_{ab}\left(c\tau \right)Q_{ab}^{}(\tau )$$ (39) Thus the response of the detector to a gravitational wave pulse travelling in $`z`$-direction is described by the time-dependence of the microscopic function $`Q_{ab}`$. Alternatively, we may use Fourier transform to arrive at $$\mathrm{\Delta }E=_{\mathrm{}}^+\mathrm{}𝑑\nu \widehat{h}_{ab}^{}(\nu )i\nu \widehat{Q}_{ab}(\nu )$$ (40) with the spectral decompositions of the wave: $$\widehat{h}_{ab}(\nu )=h_{ab}(c\tau )e^{2\pi i\nu \tau }𝑑\tau ,$$ (41) and the particle system: $$\widehat{Q}_{ab}(\nu )=\underset{s}{}\frac{q_{ab}^{(s)}(t(\tau ))}{1\dot{z}(t(\tau ))/c}e^{2\pi i\nu \tau }𝑑\tau .$$ (42) Thus $`\widehat{Q}_{ab}`$ represents the effective cross section of the particle system on the microscopic level. ## 6 Force distribution in a massive body We first consider electromagnetically coupled bodies, as ion crystals and metals, where the gravitational forces between the constituents play no role. Ion crystals have the advantage that all charges can be considered to be point-like and located on lattice points. Thus only discrete sums have to be evaluated. As we have seen, the forces driving the deviations from the tidal motion in the PFR system are identical to the forces driving the whole motion in the TT system as long as we consider only motions of order $`O\left(h\right)`$. ### 6.1 <br>Ion crystal We first consider a gravitational wave propagating in $`z`$-direction incident on a lattice of ions. In order for the system to possess a stable ground state, we have to include not only the Coulomb potential, but also some (short-ranged) repulsive potential. The Coulomb force exerted by particle $`s`$ on particle $`s`$ is given, to first order in $`h`$, by $$F_{C,ss}^a=\frac{q_sq_s}{4\pi \epsilon _0}\left(\frac{x_{ss}^a}{r_{ss}^3}h_{ab}\frac{x_{ss}^b}{r_{ss}^3}+\frac{3}{2}h_{bc}\frac{x_{ss}^bx_{ss}^cx_{ss}^a}{r_{ss}^5}\right),$$ (43) for the $`x`$\- and $`y`$\- directions, the force in $`z`$-direction additionally involves the derivative of $`h`$, which is not of interest here. In the case of the Born-Meyer potential the repulsive potential is of exponential form $$V_{BM}(r_{ss})=A_{ss}e^{\beta r_{ss}}$$ (44) with couplings $`A_{ss}`$ that depend on the charges. Though it is necessary to include this potential, its precise form is not crucial for our further considerations. We obtain an additional force from this potential, up to first order in $`h`$ given by $`F_{BM,ss}^a`$ $`=`$ $`{\displaystyle \frac{V_{BM,}}{x^{a(s)}}},`$ (45) $`=`$ $`A_{ss}\beta {\displaystyle \frac{e^{\beta r_{ss}}}{r_{ss}}}x_{ss}^a+{\displaystyle \frac{1}{2}}A_{ss}\beta {\displaystyle \frac{e^{\beta r_{ss}}}{r_{ss}}}\left({\displaystyle \frac{1}{r_{ss}^2}}+\beta {\displaystyle \frac{1}{r_{ss}}}\right)h_{bc}x_{ss}^bx_{ss}^cx_{ss}^a`$ $`A_{ss}\beta {\displaystyle \frac{e^{\beta r_{ss}}}{r_{ss}}}h_{ab}x_{ss}^b.`$ The total force on ion $`s`$ is given by $`F_s^a`$ $`=`$ $`{\displaystyle \underset{ss}{}}\left({\displaystyle \frac{q_sq_s}{4\pi \epsilon _0}}{\displaystyle \frac{1}{r_{ss}^3}}+A_{ss^{}}\beta {\displaystyle \frac{e^{\beta r_{ss}}}{r_{ss}}}\right)x_{ss}^a`$ $`h_{ab}{\displaystyle \underset{ss}{}}\left({\displaystyle \frac{q_sq_s}{4\pi \epsilon _0}}{\displaystyle \frac{1}{r_{ss}^3}}+A_{ss}\beta {\displaystyle \frac{e^{\beta r_{ss}}}{r_{ss}}}\right)x_{ss}^b`$ $`+h_{bc}{\displaystyle \underset{ss}{}}\left({\displaystyle \frac{3q_sq_s}{8\pi \epsilon _0}}{\displaystyle \frac{1}{r_{ss}^5}}+{\displaystyle \frac{1}{2}}A_{ss}\beta {\displaystyle \frac{e^{\beta r_{ss}}}{r_{ss}}}\left({\displaystyle \frac{1}{r_{ss}^2}}+{\displaystyle \frac{\beta }{r_{ss}}}\right)\right)x_{ss}^bx_{ss}^cx_{ss}^a`$ Assuming that the ion chain is in its unperturbed equilibrium state, both the first and second terms vanish because the vanishing of the first term defines the equilibrium in absence of a gravitational wave, and the second term is just the first multiplied by the matrix $`h_{ab}`$. Thus the gravitational wave gives rise to a perturbation of the equilibrium state induce by the force $$\mathrm{\Delta }F_s^a=\left[h_{bc}\underset{ss}{}\left(\frac{3q_sq_s}{8\pi \epsilon _0}\frac{1}{r_{ss}^5}+\frac{1}{2}A_{ss}\beta \frac{e^{\beta r_{ss}}}{r_{ss}}\left(\frac{1}{r_{ss}^2}+\frac{\beta }{r_{ss}}\right)\right)x_{ss}^bx_{ss}^c\right]x_{ss}^a$$ (47) Once the system is driven out of the equilibrium state, we still can ignore the second term in (LABEL:27), because for deviations $`\mathrm{\Delta }x_{ss}^ah`$ from equilibrium this term is only of order $`h^2`$. The first term then describes phonon-graviton interaction. The sum over all other ion in (47) now depends on the dimension of the lattice. The sum over the short-ranged part converges even for an infinite lattice, so that its contribution is always limited. The sum over the Coulomb part is an alternating sum. In one dimension, this sum is always bounded by the first term that is not canceled by a contribution from a symmetric neighbor, thus the force has a maximum on the endpoints and decreases proportional to $`R^2`$ with the distance $`R`$ from the endpoints. In two dimensions we find that the force distribution exhibits the correct quadrupole structure . Therefore acoustic modes are excited. This result is obtained only when the short-range potential is included. The Coulomb forces alone gives rise to forces that additionally alternate in direction from ion to ion, so that we would have arrived at the wrong conclusion that optical modes are excited. The forces decay rapidly away from the boundary and do not follow the linear law of the tidal accelerations. Thus, the deviations from the tidal motion excite a broad spectrum of modes, of which few are resonantly driven. We now present a general argument that this pertains to the relevant case of three dimensions. ### 6.2 General case Since the temperature of the body must be low in order to achieve the desired sensitivity of a detector, we assume that the mass is a perfect crystal, either an ion crystal with discrete charges locate on some lattice, or a metal with ions on some lattice and the electron gas in between. The body is decomposed into a finite number of elementary cells $`V_i`$ with their charge centers at $`\stackrel{}{r}_i`$, that are (i) electrically neutral, $$_{V_i}d^Dr\rho (\stackrel{}{r})=0$$ (48) and (ii) do not possess an electric dipole moment, $$_{V_i}d^Dr\left(\stackrel{}{r}\stackrel{}{r}_i\right)\rho (\stackrel{}{r})=0.$$ (49) $`D`$ is the spatial dimension of the lattice. We now consider a charge element $`q^{}=\rho (\stackrel{}{r}^{})dV`$ located at $`\stackrel{}{r}^{}`$ and the perturbational force exerted on it by the elementary cell $`V`$ located at $`\stackrel{}{r}_i=0`$, according to (47). For $`\left\|\stackrel{}{r}^{}\right\|1/\beta `$ we can ignore the short-range potential and have $$F_s^a\left(\stackrel{}{r}^{}\right)=q^{}_Vd^Dr\frac{3\rho (\stackrel{}{r})}{8\pi \epsilon _0}\frac{\left(\stackrel{}{r}^{}\stackrel{}{r}\right)^a}{\left\|\stackrel{}{r}^{}\stackrel{}{r}\right\|^5}h_{bc}\left(\stackrel{}{r}^{}\stackrel{}{r}\right)^b\left(\stackrel{}{r}^{}\stackrel{}{r}\right)^c.$$ (50) We expand the integrand into powers of $`\stackrel{}{r}=(`$ $`x,y,z)`$ up to second order. Then the integrals of the first and second order vanish due to conditions (48) and (49), respectively. As an example, if $`h`$ is of $`+`$-polarization, the forces in two dimensions are given by $`F_{q,V}^x\left(\stackrel{}{r}^{}\right)`$ $`=`$ $`q^{}h_+\left[{\displaystyle \frac{x^{}\left(6x^29y^2\right)}{2r^7}}I_1+{\displaystyle \frac{x(3x^2+12y^2)}{2r^7}}I_2+{\displaystyle \frac{y^{}\left(12x^23y^2\right)}{r^7}}I_3\right]`$ $`F_{q,V}^y\left(\stackrel{}{r}^{}\right)`$ $`=`$ $`q^{}h_+\left[{\displaystyle \frac{y^{}\left(12x^23y^2\right)}{2r^7}}I_1+{\displaystyle \frac{y^{}\left(9x^2+6y^2\right)}{2r^7}}I_2+{\displaystyle \frac{x^{}\left(3x^2+12y^2\right)}{r^7}}I_3\right]`$ where the integrals $$I_1=_Vd^Drx^2\frac{3\rho (\stackrel{}{r})}{8\pi \epsilon _0},I_2=_Vd^Dry^2\frac{3\rho (\stackrel{}{r})}{8\pi \epsilon _0},I_3=_Vd^Drxy\frac{3\rho (\stackrel{}{r})}{8\pi \epsilon _0}$$ (53) describe the quadrupole moments of the electric charge distribution in the elementary cell. The generalization to three dimensions is straightforward, with the same qualitative properties: The force exerted by some elementary cell on a charge element $`q^{}`$ always decreases at least proportional to $`r^4`$ with the distance between the charge and the cell. If the quadrupole moments vanish, as for cubic lattices, the decrease is even faster by two powers of $`r^{}`$.This implies that for a given charge element $`q^{}`$ the sum over all elementary cells always converges in $`D=1,2,`$ or $`3`$ dimensions. Therefore the force on any ion in the body is bounded independently from the size of the body, $$\left\|\stackrel{}{F}_q\right\|_{\mathrm{max}}\underset{i}{}\left\|F_{q,V_i}^y\left(\stackrel{}{r}^{}\stackrel{}{r}_i\right)\right\|const\underset{i}{}\frac{1}{\left\|\stackrel{}{r}^{}\stackrel{}{r}_i\right\|^4}.$$ (54) Further, the reflection symmetry of the lattice implies that all forces from cells in a volume that possesses reflection symmetry around the charge element add up to zero. Hence also in the general case the charge element feels a force that depends on its distance $`R`$ to the surface of the body. A crude estimate gives $$\left\|\stackrel{}{F}_q\right\|_R\underset{i,\|\stackrel{}{r}\stackrel{}{r}_i\|>R}{}\left\|F_{q,V_i}^y\left(\stackrel{}{r}^{}\stackrel{}{r}_i\right)\right\|const\underset{\|\stackrel{}{r}\stackrel{}{r}_i\|>R}{}\frac{1}{\left\|\stackrel{}{r}^{}\stackrel{}{r}_i\right\|^4}R^{D4}.$$ (55) Depending on the dimension, we observe that the force decreases with the distance from the surface, at least with $`1/R`$ in three dimensions, if the quadrupole moments $`I_1,I_2,I_3`$ all vanish the decrease is even stronger. Scaling with the lattice constant $`a`$ leads us to $$\left\|\stackrel{}{F}_q\right\|_R\left\|\stackrel{}{F}_q\right\|_{\mathrm{max}}\left(\frac{a}{R}\right)^{4D}.$$ (56) Because the sum over the forces from the short range potential decreases even faster, the total perturbational forces, which are maximal on the surface, decrease to about $`10^3`$ of the surface value within 1000 atomic layer, which is about 1micrometer. When we integrate the forces (LABEL:32,LABEL:33) over an elementary cell in order to obtain the mean force on the cell, we again loose two powers due to (48) and (49), resulting in mean forces between two cells $`V_i`$ and $`V_j`$ that decrease at best proportional to $`\left\|\stackrel{}{r}_j\stackrel{}{r}_i\right\|^6`$. Thus we conclude that the bulk of the material remains, apart from tidal motion, unaffected by the forces induced by the gravitational wave. This result has its origin in the nature of the electromagnetic coupling with its charges of different signs, and the structure of solid bodies with reflection symmetry of the elementary cells. ### 6.3 Gravitationally coupled matter Clearly, the result of the preceding section was due to the conditions (48,49), but (48) does not hold for the attractive gravitational forces. The acceleration of a test mass is given by $$\stackrel{}{a}_G\left(\stackrel{}{r}^{}\right)=_Vd^Dr\frac{\rho _m(\stackrel{}{r})}{G}\frac{\left(\stackrel{}{r}^{}\stackrel{}{r}\right)}{\left\|\stackrel{}{r}^{}\stackrel{}{r}\right\|^5}h_{bc}\left(\stackrel{}{r}^{}\stackrel{}{r}\right)^b\left(\stackrel{}{r}^{}\stackrel{}{r}\right)^c,$$ (57) where $`\rho _m`$ is the mass density of the body. Assuming that the wavelength of the gravitational wave is large compared to the dimensions of a homogeneous massive body, we can take $`\rho _m`$ and $`h`$ to be constant and are able to evaluate the integral for simple geometries. For a $`+`$-polarized wave propagating in $`z`$-direction the maximal acceleration on the surface is calculated to $$\left\|\stackrel{}{a}_G\right\|_{\mathrm{max}}=\{\begin{array}{c}\rho _m\pi r\left(1\frac{1}{\sqrt{1+\frac{l^2}{r^2}}}\right)\text{ for a cylinder of radius }r\text{ and lenght }l\hfill \\ \rho _m\pi \frac{\mathrm{tan}^2\varphi }{\left(1+\mathrm{tan}^2\varphi \right)^{3/2}}l\text{ for the tip of a cone of opening angle }\varphi \text{ and height }l\hfill \\ \rho _m\pi \frac{8R}{15}\text{ at the surface a sphere of radius }R.\hfill \end{array}$$ (58) Naturally, there exists an angular dependence of the forces of quadrupole characteristic. Thus the forces grow linearly with the linear dimensions of the massive body. The maximal acceleration on the equator of a sphere is given by $$\left\|\stackrel{}{a}_G\right\|_{\mathrm{max}}^{TT}=\frac{2}{5}\left\|h\right\|g$$ (59) where $`g`$ is the surface acceleration of the mass. This force will truly excite the quadrupole modes in the bulk of a body and is able to do work against the gravitational and electromagnetic forces that keep the body together. Comparing this with the tidal acceleration in the PFR system, $$\left\|\stackrel{}{a}_G\right\|_{\mathrm{max}}^{PFR}=\frac{1}{2}\omega ^2R\left\|h\right\|,$$ (60) we see that (59) is several orders of magnitude smaller than (60). In the PFR system the forces drive primarily the tidal motion, only accelerations of order (59) drive the deviations. ## 7 <br>Detector Types ### 7.1 Rotational Detectors From the structure of the microscopic quantity $`Q_{ab}(\tau )`$ we can immediately identify the different types of detectors. If the body is rotating relative to the TT frame with some frequency $`\omega `$, then the momentum contribution $$\underset{s}{}\frac{1}{2m_s}p_a^{(s)}p_b^{(s)}\omega ^2\mathrm{cos}2\omega t$$ (61) dominates and gives rise to a response of order $`h`$. This type of detector was first suggested by Braginsky . We estimate the response of a rotating mass to the gravitational wave using our model. From (37) we obtain a momentum input to the center of mass of $$dP_z/dt=h_0\mathrm{cos}\left(\omega _0t+\phi \right)\mathrm{cos}2\omega t\frac{\omega ^2\omega _0l^2M}{24c}$$ (62) where $`l`$ is the length of the bar, $`M`$ its mass, and $`h_0,`$ $`\omega _0,`$ and $`\phi `$ are the amplitude, frequency and phase difference of the gravitational wave, respectively. This is for the case where the wave vector is perpendicular to the plane of rotation. In the ideal case $`\omega \omega _0/2,\phi =0`$ or $`\pi `$, the mean energy input is given by $$\dot{E}=\pm h_0\frac{\omega _0^3l^2M}{192}.$$ (63) Note that the change can be of either sign, thus the gravitational wave cannot only be absorbed, but can also stimulate the emission of gravitational waves from the system. For reasonable values (bar of $`1m`$, $`1100kg`$ mass, $`\omega /2\pi 101000`$ Hertz, $`h_010^{20}`$) the attainable energy input ranges in about $$\dot{E}10^{13}\mathrm{}10^9W.$$ (64) Though this is quite large compared to the resonant detector, it seems questionable whether this change in the rotational energy of the bar can be measured. Certainly the acceleration of the rotating bar in the direction of the wave, given by (62), is too small to be detectable. This gives additional justification to the negligence of the acceleration of the center of mass in the transformation to the PFR system . ### 7.2 Resonant Detectors For a non-rotating mass, the lowest order contribution stems from the time derivatives of the positions and momenta induced by the wave and thus leads to a response of order $`h^2`$. This is the Weber detector. Our considerations have shown that the energy input into such a detector occurs primarily at its surface. This implies that when resonance is discussed, we have to take the time into account that is needed to transport form the surface to the interior of the body. For a material with high Q-factor the velocity of transport is given by the velocity of sound. Thus it takes a time $$T_v=\frac{L}{2v_s}$$ (65) where $`L`$ is the diameter of the body and $`v_s`$ the velocity of sound before energy reaches the center of the body. This time limits the onset of resonance. For example, for GRAIL \[grail\] with diameter $`L=3m`$, $`v_s4000m/s`$ we have $`T_v3/8`$ milliseconds. Thus this time scale will play a role for the detection of millisecond pulses. In order to arrive at the resonant amplitude, pulses must be considerably longer than $`T_v`$. How many oscillations are needed to arrive at the maximal resonant amplitude, can also be estimated in the normal-mode model. In general, if a weakly damped oscillator characterized by a frequency $`\omega _0`$ is excited by a force $`F=\epsilon \omega ^2sin\left(\omega t\right)`$ near resonance, $`\omega \omega _0`$, then if the oscillator is initially at rest, the subsequent maxima of the amplitude of the oscillator will at the beginning follow a linear law $$\left\|A_{max}\right\|\frac{\epsilon \omega ^2}{2\omega _0}t$$ (66) until the resonant amplitude of order $`C=\epsilon C_0`$ is reached after a transient time $$T_{trans}\frac{2C_0}{\omega }\frac{2C_0}{\omega _0}.$$ (67) Thus e.g. it takes for a signal of frequency $`\omega =1000Hz`$ about $`2s`$ until it is enhanced by a factor $`C_0=1000`$. Thus $`2000`$ oscillations are needed before the maximal resonant level is reached. In general the effective cross section of a resonant mass detector is proportional to $`\left(C_0\right)^2`$, but for pulses shorter than $`T_{trans}`$ an additional factor $`\left(T/T_{trans}\right)^2`$ must be taken into account. This is relevant to the detection of millisecond pulses with a detector like GRAIL that operates at a fundamental frequency of $`650Hz`$ . A further point to consider certainly is the impurity of the material. When grains of material stick together, we also have to consider internal boundary that behave like surfaces. But because this results in random effects, the impurity will hardly improve the behavior of the detector. ### 7.3 Microscopic Detector We have shown that in the TT system the metric effects on the electromagnetic coupling cancel in the bulk of a massive body, in the PFR system the deviations from the tidal motion are driven only at the surface. This raises the question whether this type of surface effect could be used in some other way for detection of gravitational waves. Similar to expression (59) for the gravitationally coupled matter, we can estimate the acceleration induce by the gravitational wave on the surface of an ion lattice by $$\left\|\stackrel{}{a}_{EM}\right\|_{\mathrm{max}}^{TT}\left\|\frac{Z_+Z_{}}{m_\pm d^2\pi \epsilon _0}h\right\|$$ (68) where $`Z_+,Z_{}`$ and $`m_+,m_{}`$ are the charges and the masses of the ions, respectively, and $`d`$ the lattice constant. A geometric factor of order $`1`$ has to be included in addition. This factor will depend on the structure of the lattice and the precise behavior of the repulsive potential. The numerical value of (68) is of the order of $`h10^{16}ms^2`$ (for KCl) which is several orders of magnitude larger than (60). This suggests that the piezoelectric effect might be used for experimental verification, in a similar way like a pressure sensor works. For $`h`$ in the order of $`10^{22}`$ accelerations are of order $`10^6ms^2`$ and possibly are within the reach of sensitive detectors. If we consider a piezo crystal under strain our analysis has to be revised. Because we then have a dipole moment in each elementary cell, the leading order of the forces will be proportional to $`r^3`$ as compared to $`r^4`$ in (LABEL:32) and (LABEL:33). Due to the uniform direction of the dipoles and the broken refection symmetry the contributions in the integration will no longer cancel, so we expect that the forces grow logarithmically with the dimension of the crystal and the excitation occurs truly throughout the bulk. We hope to be able present a detailed analysis of this case, together with possible detector design, soon elsewhere. ## 8 Conclusions We analyzed the resonant mass gravitational wave detector form a microscopic point of view, using the wave guide (TT) frame of reference along with the conventional PFR system. In the TT system the variation of the Coulomb field of the constituent charges of the body gives the dominant contribution, the resulting forces agrees with those driving the deviation from the tidal motion in the PFR system. For an electromagnetically coupled body with reflection symmetry, this force distribution is such that the forces are restricted to a small surface layer, with no bulk force. Hence the relevant energy input occurs at the surface only. This effect has its origin in the fact that the pattern of the tidal forces comes close to a null mode of the system, not doing work against the coupling forces, as we saw from a linear chain model. This result is not in contradiction to the standard normal mode analysis, rather it reflects a local property of the response compared to global nature of normal modes. For the gravitationally coupled body, we observe a linear force law and bulk excitation, but the part that causes resonant response is orders of magnitude smaller than expected from the tidal forces in the PFR system. Regarding resonant mass detectors, the local nature of the driving force makes it necessary to consider the time scale on which the energy is transported into the interior. This time scale is set by the velocity of sound. As a result, it might take too long a time for a resonant detector to be excited measurably by short millisecond pulses as are emitted by collapsing stellar objects. Finally, we presented ideas for a new type of microscopic detector that employs the piezo effect. The force pattern on a crystal is basically that of a quadrupole-like distributed pressure change, the surface force is independent of the size of the crystal. For a crystal under strain the force even acts throughout the bulk. An analysis how the desired sensitivity can be achieved was outside the scope of the work presented here.

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