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# Mathematical Foundations of Holonomic Quantum Computer
## 1 Introduction
Quantum Computer is a very attractive and challenging object for New Science.
After the breakthrough by P. Shor there has been remarkable progress in Quantum Computer or Computation (QC briefly). This discovery had a great influence on scientists. This drived not only theoreticians to finding other quantum algorithms, but also experimentalists to building quantum computers. See and in outline. is also very useful.
On the other hand, Gauge Theories are widely recognized as the basis in quantum field theories. Therefore it is very natural to intend to include gauge theories in QC $`\mathrm{}`$ a construction of “gauge theoretical” quantum computation or of “geometric” quantum computation in our terminology. The merit of geometric method of QC is strong against the influence from the environment. See .
Zanardi and Rasetti in and proposed such an idea using non-abelian Berry phase (quantum holonomy). See also and as another geometric models. In their model a Hamiltonian (including some parameters) must be degenerated because an adiabatic connection is introduced using this degeneracy . In other words, a quantum computational bundle on some parameter space (see ) is introduced due to this degeneracy.
They gave a few simple but interesting examples to explain their idea. We believe that these examples will become important in the near future. But their works (, and ) are a bit coarse in the mathematical point of view. Therefore in this paper we give a mathematical reinforcement to them. See and also as a further generalization.
It is not easy to predict the future of geometric quantum computation. However it is an arena worth challenging for mathematical physicists.
## 2 Mathematical Foundation of Quantum Holonomy
We start with mathematical preliminaries. Let $``$ be a separable Hilbert space over $`𝐂`$. For $`m𝐍`$, we set
$$St_m\left(\right)\left\{V=(v_1,\mathrm{},v_m)\times \mathrm{}\times \right|V^{}V=1_m\},$$
(1)
where $`1_m`$ is a unit matrix in $`M(m,𝐂)`$. This is called a (universal) Stiefel manifold. Note that the unitary group $`U(m)`$ acts on $`St_m\left(\right)`$ from the right:
$$St_m\left(\right)\times U\left(m\right)St_m\left(\right):(V,a)Va.$$
(2)
Next we define a (universal) Grassmann manifold
$$Gr_m\left(\right)\left\{XM\left(\right)\right|X^2=X,X^{}=X\mathrm{and}\mathrm{tr}X=m\},$$
(3)
where $`M()`$ denotes a space of all bounded linear operators on $``$. Then we have a projection
$$\pi :St_m\left(\right)Gr_m\left(\right),\pi \left(V\right)VV^{},$$
(4)
compatible with the action (2) ($`\pi \left(Va\right)=Va(Va)^{}=Vaa^{}V^{}=VV^{}=\pi \left(V\right)`$).
Now the set
$$\{U\left(m\right),St_m\left(\right),\pi ,Gr_m\left(\right)\},$$
(5)
is called a (universal) principal $`U(m)`$ bundle, see and . We set
$$E_m\left(\right)\left\{(X,v)Gr_m\left(\right)\times \right|Xv=v\}.$$
(6)
Then we have also a projection
$$\pi :E_m\left(\right)Gr_m\left(\right),\pi \left((X,v)\right)X.$$
(7)
The set
$$\{𝐂^m,E_m\left(\right),\pi ,Gr_m\left(\right)\},$$
(8)
is called a (universal) $`m`$-th vector bundle. This vector bundle is one associated with the principal $`U(m)`$ bundle (5) .
Next let $`M`$ be a finite or infinite dimensional differentiable manifold and the map $`P:MGr_m\left(\right)`$ be given (called a projector). Using this $`P`$ we can make the bundles (5) and (8) pullback over $`M`$ :
$`\{U\left(m\right),\stackrel{~}{St},\pi _{\stackrel{~}{St}},M\}P^{}\{U\left(m\right),St_m\left(\right),\pi ,Gr_m\left(\right)\},`$ (9)
$`\{𝐂^m,\stackrel{~}{E},\pi _{\stackrel{~}{E}},M\}P^{}\{𝐂^m,E_m\left(\right),\pi ,Gr_m\left(\right)\},`$ (10)
see . (10) is of course a vector bundle associated with (9).
Let $``$ be a parameter space and we denote by $`\lambda `$ its element. Let $`\lambda _\mathrm{𝟎}`$ be a fixed reference point of $``$. Let $`H_\lambda `$ be a family of Hamiltonians parametrized by $``$ which act on a Fock space $``$. We set $`H_0`$ = $`H_{\lambda _\mathrm{𝟎}}`$ for simplicity and assume that this has a $`m`$-fold degenerate vacuum :
$$H_0v_j=\mathrm{𝟎},j=1m.$$
(11)
These $`v_j`$’s form a $`m`$-dimensional vector space. We may assume that $`v_i|v_j=\delta _{ij}`$. Then $`(v_1,\mathrm{},v_m)St_m\left(\right)`$ and
$$F_0\left\{\underset{j=1}{\overset{m}{}}x_jv_j\right|x_j𝐂\}𝐂^m.$$
Namely, $`F_0`$ is a vector space associated with o.n.basis $`(v_1,\mathrm{},v_m)`$.
Next we assume for simplicity that a family of unitary operators parametrized by $``$
$$W:U(),W(\lambda _\mathrm{𝟎})=\mathrm{id}.$$
(12)
is given and $`H_\lambda `$ above is given by the following isospectral family
$$H_\lambda W(\lambda )H_0W(\lambda )^1.$$
(13)
In this case there is no level crossing of eigenvalues. Making use of $`W(\lambda )`$ we can define a projector
$$P:Gr_m\left(\right),P(\lambda )W(\lambda )\left(\underset{j=1}{\overset{m}{}}v_jv_j^{}\right)W(\lambda )^1$$
(14)
and have the pullback bundles over $``$
$$\{U\left(m\right),\stackrel{~}{St},\pi _{\stackrel{~}{St}},\},\{𝐂^m,\stackrel{~}{E},\pi _{\stackrel{~}{E}},\}.$$
(15)
For the later we set
$$|vac=(v_1,\mathrm{},v_m).$$
(16)
In this case a canonical connection form $`𝒜`$ of $`\{U\left(m\right),\stackrel{~}{St},\pi _{\stackrel{~}{St}},\}`$ is given by
$$𝒜=vac|W(\lambda )^1dW(\lambda )|vac,$$
(17)
where $`d`$ is a differential form on $``$, and its curvature form by
$$d𝒜+𝒜𝒜,$$
(18)
see and .
Let $`\gamma `$ be a loop in $``$ at $`\lambda _\mathrm{𝟎}`$., $`\gamma :[0,1],\gamma (0)=\gamma (1)`$. For this $`\gamma `$ a holonomy operator $`\mathrm{\Gamma }_𝒜`$ is defined :
$$\mathrm{\Gamma }_𝒜(\gamma )=𝒫exp\left\{_\gamma 𝒜\right\}U\left(m\right),$$
(19)
where $`𝒫`$ means path-ordered. This acts on the fiber $`F_0`$ at $`\lambda _\mathrm{𝟎}`$ of the vector bundle $`\{𝐂^m,\stackrel{~}{E},\pi _{\stackrel{~}{E}},M\}`$ as follows : $`\text{x}\mathrm{\Gamma }_𝒜(\gamma )\text{x}`$. The holonomy group $`Hol(𝒜)`$ is in general subgroup of $`U\left(m\right)`$ . In the case of $`Hol(𝒜)=U\left(m\right)`$, $`𝒜`$ is called irreducible, see .
In the Holonomic Quantum Computer we take
$`\mathrm{Encoding}\mathrm{of}\mathrm{Information}\text{x}F_0,`$
$`\mathrm{Processing}\mathrm{of}\mathrm{Information}\mathrm{\Gamma }_𝒜(\gamma ):\text{x}\mathrm{\Gamma }_𝒜(\gamma )\text{x}.`$ (20)
## 3 Holonomic Quantum Computation
We apply the results of last section to Quantum Optics and discuss about (optical) Holonomic Quantum Computation proposed by and .
Let $`a(a^{})`$ be the annihilation (creation) operator of the harmonic oscillator. If we set $`Na^{}a`$ (: number operator), then
$$[N,a^{}]=a^{},[N,a]=a,[a,a^{}]=1.$$
(21)
Let $``$ be a Fock space generated by $`a`$ and $`a^{}`$, and $`\{|n|n𝐍\{0\}\}`$ be its basis. The actions of $`a`$ and $`a^{}`$ on $``$ are given by
$$a|n=\sqrt{n}|n1,a^{}|n=\sqrt{n+1}|n+1,$$
(22)
where $`|0`$ is a vacuum ($`a|0=0`$).
Next we consider the system of two-harmonic oscillators. If we set
$$a_1=a1,a_1^{}=a^{}1;a_2=1a,a_2^{}=1a^{},$$
(23)
then it is easy to see
$$[a_i,a_j]=[a_i^{},a_j^{}]=0,[a_i,a_j^{}]=\delta _{ij},i,j=1,2.$$
(24)
We also denote by $`N_i=a_i^{}a_i`$ number operators.
Now since we want to consider coherent states based on Lie algebras $`su(2)`$ and $`su(1,1)`$, we make use of Schwinger’s boson method, see , . Namely if we set
$`[\mathrm{C}]`$ $`su(2):J_+=a_1^{}a_2,J_{}=a_2^{}a_1,J_3={\displaystyle \frac{1}{2}}(a_1^{}a_1a_2^{}a_2),`$ (25)
$`[\mathrm{NC}]`$ $`su(1,1):K_+=a_1^{}a_2^{},K_{}=a_2a_1,K_3={\displaystyle \frac{1}{2}}(a_1^{}a_1+a_2^{}a_2+1),`$ (26)
then we have
$`[\mathrm{C}]`$ $`su(2):[J_3,J_+]=J_+,[J_3,J_{}]=J_{},[J_+,J_{}]=2J_3,`$ (27)
$`[\mathrm{NC}]`$ $`su(1,1):[K_3,K_+]=K_+,[K_3,K_{}]=K_{},[K_+,K_{}]=2K_3.`$ (28)
In the following we treat unitary coherent operators based on Lie algebras $`su(2)`$ and $`su(1,1)`$.
Definition We set
$`[\mathrm{C}]`$ $`U(\xi )=e^{\xi a_1^{}a_2\overline{\xi }a_2^{}a_1}\mathrm{for}\xi 𝐂,`$ (29)
$`[\mathrm{NC}]`$ $`V(\zeta )=e^{\zeta a_1^{}a_2^{}\overline{\zeta }a_2a_1}\mathrm{for}\zeta 𝐂.`$ (30)
For the details of $`U(\xi )`$ and $`V(\zeta )`$ see and . For the latter convenience let us list well-known disentangling formulas.
Lemma 1 We have
$`[\mathrm{C}]U(\xi )`$ $`=`$ $`e^{\eta a_1^{}a_2}e^{\mathrm{log}\left(1+|\eta |^2\right)\frac{1}{2}\left(a_1^{}a_1a_2^{}a_2\right)}e^{\overline{\eta }a_2^{}a_1},where\eta ={\displaystyle \frac{\xi \mathrm{tan}|\xi |}{|\xi |}},`$ (31)
$`[\mathrm{NC}]V(\zeta )`$ $`=`$ $`e^{\kappa a_1^{}a_2^{}}e^{\mathrm{log}\left(1|\kappa |^2\right)\frac{1}{2}\left(a_1^{}a_1+a_2^{}a_2+1\right)}e^{\overline{\kappa }a_2a_1},where\kappa ={\displaystyle \frac{\zeta \mathrm{tanh}|\zeta |}{|\zeta |}}.`$ (32)
As for a genelalization of these formulas see .
Let $`H_0`$ be a Hamiltonian with nonlinear interaction produced by a Kerr medium., that is $`H_0=\mathrm{}\mathrm{X}N(N1)`$, where X is a certain constant, see . The eigenvectors of $`H_0`$ corresponding to $`0`$ is $`\{|0,|1\}`$, so its eigenspace is $`\mathrm{Vect}\{|0,|1\}𝐂^2`$. We correspond to $`0|0,1|1`$ for a generator of Boolean algebra $`\{0,1\}`$. The space $`\mathrm{Vect}\{|0,|1\}`$ is called 1-qubit (quantum bit) space, see or . Since we are considering the system of two particles, the Hamiltonian that we treat in the following is
$$H_0=\mathrm{}\mathrm{X}N_1(N_11)+\mathrm{}\mathrm{X}N_2(N_21).$$
(33)
The eigenspace of $`0`$ of this Hamiltonian becomes therefore
$$F_0=\mathrm{Vect}\{|0,|1\}\mathrm{Vect}\{|0,|1\}𝐂^2𝐂^2.$$
(34)
We denote the basis of $`F_0`$ as $`\{|0,0,|0,1,|1,0,|1,1\}`$ and set $`|vac=(|0,0,|0,1,|1,0,|1,1)`$.
Next we consider the following isospectral family of $`H_0`$ above :
$`H_{(\xi ,\zeta )}`$ $`=`$ $`W(\xi ,\zeta )H_0W(\xi ,\zeta )^1,`$ (35)
$`W(\xi ,\zeta )`$ $`=`$ $`U(\xi )V(\zeta )U(),W(0,0)=\mathrm{id}.`$ (36)
For this system let us calculate a connection form (17) in the last section. For that we set
$$A_\xi =vac|W(\xi ,\zeta )^1\frac{}{\xi }W(\xi ,\zeta )|vac,A_\zeta =vac|W(\xi ,\zeta )^1\frac{}{\zeta }W(\xi ,\zeta )|vac.$$
(37)
Here remaking
$`W(\xi ,\zeta )^1{\displaystyle \frac{}{\xi }}W(\xi ,\zeta )`$ $`=`$ $`V(\xi )^1\left\{U(\xi )^1{\displaystyle \frac{}{\xi }}U(\xi )\right\}V(\xi ),`$
$`W(\xi ,\zeta )^1{\displaystyle \frac{}{\zeta }}W(\xi ,\zeta )`$ $`=`$ $`V(\xi )^1{\displaystyle \frac{}{\zeta }}V(\xi )`$
and using Lemma 1,
Lemma 2 we have
$`W^1{\displaystyle \frac{}{\xi }}W`$
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right)\left\{\mathrm{cosh}(2|\zeta |)a_1^{}a_2+{\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\left(a_1^{}\right)^2+{\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\left(a_2\right)^2\right\}`$
$`+`$ $`{\displaystyle \frac{\overline{\xi }}{2|\xi |^2}}\left(1\mathrm{cos}(2|\xi |)\right){\displaystyle \frac{1}{2}}\left(a_1^{}a_1a_2^{}a_2\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\xi }^2}{2|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right)\left\{\mathrm{cosh}(2|\zeta |)a_2^{}a_1+{\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\left(a_1\right)^2+{\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\left(a_2^{}\right)^2\right\},`$
$`W^1{\displaystyle \frac{}{\zeta }}W`$
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\right)a_1^{}a_2^{}+{\displaystyle \frac{\overline{\zeta }}{2|\zeta |^2}}\left(1+\mathrm{cosh}(2|\zeta |)\right){\displaystyle \frac{1}{2}}\left(a_1^{}a_1+a_2^{}a_2+1\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta }^2}{2|\zeta |^2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta |)}{2|\xi |}}\right)a_1a_2.`$ (39)
From this lemma it is easy to calculate $`A_\xi `$ and $`A_\zeta `$. Before stating the result let us prepare some notations.
$$\widehat{E}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right),\widehat{F}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 0\end{array}\right),\widehat{H}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& \frac{1}{2}& 0& 0\\ 0& 0& \frac{1}{2}& 0\\ 0& 0& 0& 0\end{array}\right).$$
(40)
$$\widehat{A}=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right),\widehat{C}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 1& 0& 0& 0\end{array}\right),\widehat{B}=\left(\begin{array}{cccc}\frac{1}{2}& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& \frac{3}{2}\end{array}\right).$$
(41)
Proposition 3 We have
$`A_\xi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right)\mathrm{cosh}(2|\zeta |)\widehat{F}{\displaystyle \frac{\overline{\xi }}{2|\xi |^2}}\left(1\mathrm{cos}(2|\xi |)\right)\widehat{H}`$ (42)
$`+{\displaystyle \frac{\overline{\xi }^2}{2|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right)\mathrm{cosh}(2|\zeta |)\widehat{E},`$
$`A_\zeta `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\right)\widehat{C}+{\displaystyle \frac{\overline{\zeta }}{2|\zeta |^2}}\left(1+\mathrm{cosh}(2|\zeta |)\right)\widehat{B}`$ (43)
$`+{\displaystyle \frac{\overline{\zeta }^2}{2|\zeta |^2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\right)\widehat{A}.`$
Since the connection form $`𝒜`$ is anti-hermitian ($`𝒜^{}=𝒜`$), it can be written as
$$𝒜=A_\xi d\xi +A_\zeta d\zeta A_\xi ^{}d\overline{\xi }A_\zeta ^{}d\overline{\zeta },$$
(44)
so that it’s curvature form $`=d𝒜+𝒜𝒜`$ becomes
$``$ $`=`$ $`\left(_\xi A_\zeta _\zeta A_\xi +[A_\xi ,A_\zeta ]\right)d\xi d\zeta `$ (45)
$`\left(_\xi A_\xi ^{}+_{\overline{\xi }}A_\xi +[A_\xi ,A_\xi ^{}]\right)d\xi d\overline{\xi }`$
$`\left(_\xi A_\zeta ^{}+_{\overline{\zeta }}A_\xi +[A_\xi ,A_\zeta ^{}]\right)d\xi d\overline{\zeta }`$
$`\left(_\zeta A_\xi ^{}+_{\overline{\xi }}A_\zeta +[A_\zeta ,A_\xi ^{}]\right)d\zeta d\overline{\xi }`$
$`\left(_\zeta A_\zeta ^{}+_{\overline{\zeta }}A_\zeta +[A_\zeta ,A_\zeta ^{}]\right)d\zeta d\overline{\zeta }`$
$`\left(_{\overline{\xi }}A_\zeta ^{}_{\overline{\zeta }}A_\xi ^{}+[A_\zeta ^{},A_\xi ^{}]\right)d\overline{\xi }d\overline{\zeta }.`$
Now we state our main result.
Theorem 4
$`=`$
$`\left\{\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{F}+{\displaystyle \frac{\overline{\xi }^2}{|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{E}\right\}d\xi d\zeta `$
$`\{{\displaystyle \frac{\xi }{|\xi |^2}}(1+\mathrm{cos}(2|\xi |))\mathrm{cosh}(2|\zeta |)\widehat{F}{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{|\xi |}}(1+\mathrm{cosh}^2(2|\zeta |))\widehat{H}`$
$`+{\displaystyle \frac{\overline{\xi }}{|\xi |^2}}(1+\mathrm{cos}(2|\xi |))\mathrm{cosh}(2|\zeta |)\widehat{E}\}d\xi d\overline{\xi }`$
$`\left\{\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{F}+{\displaystyle \frac{\overline{\xi }^2}{|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{E}\right\}d\xi d\overline{\zeta }`$
$`\left\{\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{E}+{\displaystyle \frac{\xi ^2}{|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{F}\right\}d\zeta d\overline{\xi }`$
$`{\displaystyle \frac{\mathrm{sinh}(2|\zeta |)}{|\zeta |}}\left(2\widehat{B}\text{1}_4\right)d\zeta d\overline{\zeta }`$
$`+\left\{\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{E}+{\displaystyle \frac{\xi ^2}{|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{F}\right\}d\overline{\xi }d\overline{\zeta }.`$ (46)
From this and the theorem of Ambrose–Singer (see ) it is easy to see that
Corollary 5
$$Hol(𝒜)=SU(2)\times U(1)U(4).$$
(47)
Therefore $`𝒜`$ is not irreducible.
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# 1 Introduction
## 1 Introduction
In recent years, modeling traffic flow has been the subject of comprehensive studies by statistical physicists . Needless to say many general phenomena in vehicular traffic can be explained in general terms with these models. Distinct traffic states have been identified and some of these models have found empirical applications in real traffic . In these investigations, various theoretical approaches namely microscopic car-following models , hydro-dynamical coarse-grained macroscopic models , and gas-kinetic models have been developed in order to find a better quantitative as well as qualitative understanding toward vehicular traffic phenomena. Recently as an alternative microscopic description, Probabilistic Cellular Automata (PCA) have come into play (for an overview see Refs. ). This approach to theoretical description of traffic flow is one of the most effective and well-established ones and there is a relatively rich amount of results both numeric and analytic in the literature .
In PCA models, space (road), time and velocities of vehicles are assumed to take discrete values. This realization of traffic flow provides PCA as an ideal tool for the computer simulation. One of the prototype PCA models is the so-called Nagel-Schreckenberg (Na-Sch) model which describes a single-lane traffic flow. Although the initial observations of the Na-Sch model were numerical, shortly thereafter analytical technique were also proposed . Analytical treatments to CA are difficult in general. This is mainly due to the discreteness and the use of parallel (synchronous) updating procedures which produce the largest correlation among the vehicles with regard to other updating schemes.
Soon after its introduction, the Na-Sch model was extended to account for more realistic situations such as multi-lane traffic flow , bi-directional roads and urban traffic . In multi-lane traffic, fast cars are capable of passing the slow ones by using the fast-lane. The possibility of lane-changing allows for these models to exhibit non-trivial and interesting properties which are exclusive to multi-lane traffic flow. Despite the quite large approximative methods applied to single-lane Na-Sch based models, there are few analytical approaches to multi-lane traffic flow . One main reason is the large number of rules in PCA modeling multi-lane traffic. In reality, a driver attempting to overtake the car ahead (in a uni-directional road) has to take the following criteria into consideration:
1) There must be enough forward space in the passing-lane
2) There must be enough backward space in passing-lane so that no accident could occur between two simultaneously passing cars.
Moreover, in bi-directional roads additional criteria are necessary for a successful passing (for details see ). The main purpose of the present paper is to introduce an analytical approach to study a uni-directional two-lane road. The approach we use is to some extent similar to PCA, however basic differences are distinguishable. The major distinction is concerned with the type of updating scheme. In contrast to PCA which are realized in parallel update, our models are based on time-continuous random sequential update. The mechanism of modeling the two-lane traffic we use, is based on the stochastic reaction-diffusion processes, however the rules have roots in the Na-Sch rules. This paper is organized as follows: In section two we define the first model ( model I ) and interpret the rules in terms of those in Na-Sch model. Section three starts with the Hamiltonian description of the related master equation and continues with mean field rate equations and their solutions. The results of the numerical simulation of the model I ends this section. Next we introduce the second model (model II) in section four which is formulated in symmetric as well as asymmetric versions and follow the same steps performed in section three to obtain the fundamental diagrams of the both versions. The paper ends with some concluding remarks in section five.
## 2 Definitions of the Models
In the first model, a uni-directional two-lane road is approximated by a set of two parallel one dimensional chains, each with $`N`$ sites. The periodic boundary condition applies to both. Cars are considered as particles which occupy sites of the chains. Two type of cars exist in the system: slow cars which are denoted by $`A`$ and fast cars denoted by $`B`$. Also $`\mathrm{\Phi }`$ represents an empty site. Each site of the chains is either empty, occupied by a slow or by a fast car.
Fast cars can pass the slow ones with certain probabilities while approaching them. The bottom lane is the home-lane and cars are only allowed to use the top lane for passing. Once the passing process is achieved, they should return to the home-lane. This realization of a two-lane road is regarded as ”asymmetric ” type. Nonetheless ”symmetric ” type could also be implemented where passing from the right is allowed as well.
In model I, we restrict ourselves to ”asymmetric ” type. The state of the system is characterized by two sets of occupation numbers $`(\xi _1,\xi _2,\mathrm{},\xi _N)`$ and $`(\sigma _1,\sigma _2,\mathrm{},\sigma _N)`$ for the home and passing-lane respectively. $`\xi _i,\sigma _i=0,1,2`$ where zero refers to an empty site whereas one and two refer to a site being occupied by a slow or a fast car respectively.
To investigate the characteristics of this model, a simplification has been considered. If simultaneous two-car occupation of parallel sites of the chains is forbidden, one can describe configurations with a single set of occupation numbers $`\{\tau _i\}`$ where $`\tau _i=0,1,2`$.
Inspired by the $`v_{max}=2`$ version of the Na-Sch model , we propose the following set of stochastic processes which evolve according to a random sequential updating scheme:
$$A\mathrm{\Phi }\mathrm{\Phi }A(\mathrm{with}\mathrm{rate}h)$$
(1)
$$B\mathrm{\Phi }\mathrm{\Phi }B(\mathrm{with}\mathrm{rate}p)$$
(2)
$$A\mathrm{\Phi }\mathrm{\Phi }B(\mathrm{with}\mathrm{rate}q)$$
(3)
$$B\mathrm{\Phi }\mathrm{\Phi }A(\mathrm{with}\mathrm{rate}r)$$
(4)
$$BAAA(\mathrm{with}\mathrm{rate}\lambda )$$
(5)
$$BA\mathrm{\Phi }\mathrm{\Phi }AB(\mathrm{with}\mathrm{rate}s)$$
(6)
In order to illustrate the above definitions, let us express their interpretations :
The first and the second of the above rules correspond to the free moving of slow and fast cars respectively. The third one expresses the accelerated movement of slow cars. This step corresponds to the so-called acceleration step in the Na-Sch model. The fourth rule simulates the behaviour of a driver randomly reducing his/her speed as a result of environmental effects, road conditions etc. This step corresponds to the so-called ”random breaking” step in Na-Sch model. Finally the last two processes simulate the behaviour of the fast-car drivers when approaching a slow car. Either they pass the slow car using the passing-lane or they prefer to move behind it which give rises to their speed reduction.
We recall that in Na-Sch model, the forward movement of each car is highly affected by the car ahead. Here for simplicity we have considered the two-site interactions and only use three-site interaction for the passing process. In this particular case, it is crucial that the site ahead of the slow car should be empty. Despite the partial explanation of microscopic rules necessary for the description of a traffic flow in a two-lane road, the present model ignores the effect of oncoming fast cars (in the passing-lane) on the fast car (in the home-lane). In reality, a fast car attempts to overtake provided that there is enough back- space behind him in the passing-lane i.e. there is no passing car close to him in the passing-lane . In model (I) passing occurs locally and irrespective of the state of passing-lane behind the fast car in the home-lane.
## 3 Master Equation and Mean-Field Rate Equations
The processes (1) to (6) could be regarded as a two-species one-dimensional reaction-diffusion stochastic process. This is an example of hard-core driven lattice gas far from equilibrium which has proven to be excellent systems for theoretical investigations of low dimensional systems out of thermal equilibrium. A large variety of phenomena had already been described by driven lattice gases ( for an overview see and the references therein). Using the rates given by (1-6), one can rewrite the corresponding master equation as a Schrödinger-like equation in imaginary time.
$$\frac{}{t}|p(t)=|p(t)$$
(7)
The explicit form of $``$ could be written down via the rate equations. Let $`n_{k,A}`$ $`(n_{k,B})`$ denotes the probability that at time $`t`$, the site $`N=k`$ of the chain is occupied by a slow (fast) car. The Hamiltonian formulation of master equation allows for evaluating the average quantities in a well-established manner. It could be easily verified that the following rate equations hold for the average occupation probabilities.
$$\frac{d}{dt}n_{k,A}=hn_{k1,A}e_k+rn_{k1,B}e_k+\lambda n_{k,B}n_{k+1,A}hn_{k,A}e_{k+1}qn_{k,A}e_{k+1}$$
(8)
In the above equation $`e_k`$ stands for $`1n_{k,A}n_{k,B}`$. Similarly for $`n_{k,B}`$ we have:
$$\frac{d}{dt}n_{k,B}=qn_{k1,A}e_k+pn_{k1,B}e_k+sn_{k2,B}n_{k1,A}e_krn_{k,B}e_{k+1}pn_{k,B}e_{k+1}$$
(9)
$$\lambda n_{k,B}n_{k+1,A}sn_{k,B}n_{k+1,A}e_{k+2}$$
Apparently the total number of neither slow nor fast cars are conserved according to the dynamics and therefore the right hand side of eqs (8,9) cannot be written as a difference of two currents. However, the total number of cars i.e. the sum of slow and fast cars is a conserved quantity and the time rate of changing $`n_{A,k}+n_{B,k}`$ is equal to a difference between oncoming and outgoing currents. Summing up eqs (8) and (9) yields the following discrete form of the continuity equation:
$$\frac{d}{dt}\left[n_{k,A}+n_{k,B}\right]=J_k^{in}J_k^{out}$$
(10)
In which the explicit form of $`J_k^{out}`$ is given below.
$$J_k^{out}=hn_{k,A}e_{k+1}+rn_{k,B}e_{k+1}+qn_{k,A}e_{k+1}+pn_{k,B}e_{k+1}+sn_{k,B}n_{k+1,A}e_{k+2}$$
(11)
Equations (8), (9) and (11) are valid for arbitrary time $`t`$, however our particular interest is focused in the longtime behaviour of the system where stationarity is established. In the steady state regime, one and two-points correlators in (8,9) will be time-independent. Equation (10) implies that in steady state the current would be site-independent as expected.
So far, our result have been exact and no approximation has been implemented. At this stage and in order to solve equation (8-11) we resort to a mean-field approximation and replace the two point correlators with the product of one-point correlators. Moreover, since the closed boundary condition has been been applied, it can be anticipated that the steady values of $`n_{k,A}_s`$ and $`n_{k,B}_s`$ be site-independent and therefore we omit the site-dependence subscripts from equations (8-11). Denoting the steady values of $`n_A_s`$ and $`n_B_s`$ by $`n_A`$ and $`n_B`$ respectively, the steady current $`J`$ turns out to be
$$J=(hn_A+rn_B+qn_A+pn_B+sn_An_B)(1n)$$
(12)
In the above expression, the total density of the cars has been taken to be $`n`$
$$n_A+n_B=n$$
(13)
our final aim is to write $`J`$ in terms of total density $`n`$ and the rates . This is performed if one writes $`n_A`$ as a function $`n`$ and the rates. By applying the mean-field approximation to the equation (9) in its steady state form and using (13), one obtains the following equation
$$r(nn_A)(1n)+\lambda (nn_A)n_A=qn_A(1n)$$
(14)
which simply yields the solutions:
$$n_A=\frac{1}{2\lambda }\left[n\lambda (1n)(q+r)\pm \left([n\lambda (1n)(q+r)]^2+4rn(1n)\lambda \right)^{\frac{1}{2}}\right]$$
(15)
the solution with the minus sign is unphysical $`(n_A<0)`$ so the unique solution is the one with the positive sign. We Remark that within the mean-filed approach, one also can solve the time-dependent version of the equations (8,9). In this case, the equation for $`n_A`$ turns out to be:
$$\frac{d}{dt}n_A=rn(1n)[(q+r)(1n)n\lambda ]n_A\lambda (n_A)^2$$
(16)
which simply give rises to the following solution:
$$n_A(t)=\frac{n_AC_1e^{C_2(C_3t)}}{1e^{C_2(C_3t)}}$$
(17)
In which $`C_1,C_2`$ and $`C_3`$ are constants depending on the rates. In the long-time limit, the mean concentration of slow cars exponentially relaxes toward the steady value $`n_A`$. Replacing the above $`n_A`$ into the equation (13), one now has the total current $`J`$ as a function of $`n`$ and the rates. In order to have better insights into the problem, extended computer simulation were carried out. Here we present the result of numerical investigations of model I. In these computer simulations, the system size is typically 2400. With no loss of generality, we re-scale the time so that the rate of hopping a fast car is set to one. The speed of slow cars is supposed to be 70 percent of the speed of the fast cars which is realized by taking $`h=0.7`$ . The values of $`q`$ and $`\lambda `$ are set 1 and 0.7 respectively. One sub-update step consists of a random selection of a site, say $`N=i`$ and developing the state of the link $`(i,i+1)`$ according to the dynamics. One update step contains $`L`$ sub-updates . The typical number of updates developed in order that the system reaches to stationarity is 400000 and the averaging has been performed over 500000 updating steps. The initial state of the system was prepared randomly i.e. each site is occupied with the probability $`n`$. figures below show the result of numerical simulations.
## 4 Model II
### 4.1 Asymmetric regulation
The second model we consider, has less resemblance to the Na-Sch model. Here there is no specification of fast and slow cars and only one kind of particle exists in the chain, nevertheless the distinction between fast and slow cars is realized by their appearance in the passing and home-lanes. In this periodic double-chain model the following processes occur in a random sequential updating scheme :
$``$
$`\mathrm{with}\mathrm{rate}h`$
$``$
$``$
$`\mathrm{with}\mathrm{rate}a`$
$``$
$``$
$`\mathrm{with}\mathrm{rate}g`$
$``$
$``$
$`\mathrm{with}\mathrm{rate}b`$
$``$
As depicted, the ”asymmetric ” regulation has been adopted so that the top-lane can only be used for passing. According to the above rules, once a successful passing has taken place, the passing car should return to its home-lane unless the next site in the home-lane is already occupied. In this circumstance, it can continue to pass the second slow car ( multi-passing ). Each site of the double chain takes four different states but according to the above dynamics only three of them appears in the course of time. The forbidden state is the one in which the passing-lane site is full and its parallel home-lane site is empty. Regarding this fact, we characterize the three allowed states by $`\mathrm{\Phi },A`$ and $`B`$. $`\mathrm{\Phi }`$ represents the situation where both parallel sites are empty, $`A`$ represents the case of an occupied site in the home-lane and empty parallel site in passing-lane and finally $`B`$ refers to the case of both parallel sites being occupied.
This notation yields the following reaction-diffusion processes:
$$A\mathrm{\Phi }\mathrm{\Phi }A(\mathrm{with}\mathrm{rate}h)$$
(18)
$$AA\mathrm{\Phi }B(\mathrm{with}\mathrm{rate}a)$$
(19)
$$B\mathrm{\Phi }AA(\mathrm{with}\mathrm{rate}g)$$
(20)
$$BAAB(\mathrm{with}\mathrm{rate}b)$$
(21)
It is worth mentioning that the above model for a two-lane road is simultaneously being considered within the approach of Deterministic Cellular Automata (DCA) .
### 4.2 Master equation and mean-field approach
Similar to the steps performed in model I, one can write the following form of discrete continuity equation.
$$\frac{d}{dt}\left[n_{k,A}+2n_{k,B}\right]=J_k^{in}J_k^{out}$$
(22)
in which
$$J_k^{out}=hn_{k,A}e_{k+1}+bn_{k,B}n_{k+1,A}+gn_{k,B}e_{k+1}+an_{k,A}n_{k+1,A}$$
(23)
The above expression for $`J_k`$ has a clear interpretation in terms of rules (18-21). In steady state, the time dependences in the equation disappear and the current will be site-independent. Next we apply the mean-field approximation through which all the two-point correlators are replaced by the product of one-point correlators. This leads to the following equation for $`J`$:
$$J=hn_A(1n)+b(n\frac{n_A}{2})n_A+g(n\frac{n_A}{2})(1n)+an_A^2$$
(24)
where the relation $`\frac{n_A}{2}+n_B=n`$ has been used.
In order to obtain $`J`$ in terms of total density $`n`$ and the rates, we must write $`n_A`$ as a function of $`n`$ and the rates. This is done by solving the following equation with its left hand side set to zero.
$$\frac{d}{dt}n_A=2gn_B(1n)2an_A^2$$
(25)
The unique physical solution of the above equation is:
$$n_A=\frac{1}{4a}\left[\{(g^2(1n)^2+16an(1n)g)\}^{\frac{1}{2}}g(1n)\right]$$
(26)
putting (26) in the eq. (24), the current $`J`$ is now obtained in terms of $`n`$ and the rates. The result of computer simulations are shown in the following set of figures. Here the rate $`b,g`$ and $`h`$ are chosen to be $`1.0,1.0`$ and $`0.7`$ respectively while $`a`$ is varied. We recall that ”$`a`$” measures the tendency of fast cars to pass the slow ones.The simulation specifications are the same as those in model I.
### 4.3 Symmetric Regulation
Here we allow the fast cars to pass rightward as well. In this case, both the top and bottom lanes become identical and fast cars can pass the slow ones irrespective of their home-lane. In this symmetric two-lane model, each particle hops one site ahead in its home-lane provided that the next site is empty. Otherwise it tries to pass the car ahead. This attempt is successful if there is an empty site ahead on the opposite lane. The following rules illustrates the model definition.
$``$
$`\mathrm{with}\mathrm{rate}h`$
$``$
$``$
$`\mathrm{with}\mathrm{rate}h`$
$``$
$``$
$`\mathrm{with}\mathrm{rate}g`$
$``$
$``$
$`\mathrm{with}\mathrm{rate}g`$
$``$
The astrix symbol indicates that the process in the opposite lane occurs independently of the configuration of the sites filled with astrix. If we denote the state of two parallel sites in which the bottom site is empty and the top one is occupied by $`B`$, the state of simultaneous occupation of parallel sites by $`C`$ and adopting the notations $`\mathrm{\Phi }`$ and A as the same in the asymmetric version of the model, then it could easily be verified that the forms of the discrete continuity equation and the current are as follows:
$$\frac{d}{dt}\left(a_k+b_k+2c_k\right)=J_{k1}J_k$$
(27)
and
$`J_{k,k+1}=h\left(a_ke_{k+1}+a_kb_{k+1}+2c_ke_{k+1}+c_kb_{k+1}+b_ke_{k+1}+b_ka_{k+1}+c_ka_{k+1}\right)`$
$$+g\left(b_kb_{k+1}+a_ka_{k+1}\right)$$
(28)
Where $`a_k,b_k`$ and $`c_k`$ refer to the probabilities that at time $`t`$, the site $`N=k`$ of the double-chain has one car in bottom lane, one car in top lane and double-occupancy in both lanes respectively. In steady state, the system is both time and site independent. Denoting the steady values of $`a_k,b_k`$ and $`c_k`$ by $`a,b`$ and $`c`$, one has the relation:
$$\frac{a+b}{2}+c=n$$
(29)
Moreover, the symmetry between the lanes implies that $`a=b`$. The steady value $`a`$ is easily found to be obtained from the following equation:
$$(g+h)a^2=hc(1n)$$
(30)
Solving the steady-state equation for $`a`$, one finds:
$$a=\frac{(\left[h^2(1n)^2+4hn(1n)(g+h)\right])^{\frac{1}{2}}h(1n)}{2(g+h)}$$
(31)
Also equation (31) leads to the following equation for $`J`$.
$$J=2[hn(1n)+h\{a^2+2a(na)\}+ga^2)]$$
(32)
Where by putting the eq. (31) into it, one reaches to expression for $`J`$ in terms of $`n,g`$ and $`h`$. We remark that the factor two reflects the number of lanes. The result of computer simulations are shown in the following set of figures. The value of $`h`$ is set to one and $`g`$ is varied.
## 5 Concluding Remarks
We have introduced a two-species reaction-diffusion model for description of a uni-directional two-lane road. The type of update we have used is random-sequential which sounds more appropriate for analytical treatments. In the first model, the result of numeric simulations are very close to those in mean-field approach which indicates that the effects of correlations are suppressed. However in the second model, there are remarkable differences between analytical and numeric results. in model I, the current-density diagram is slightly affected by changing the passing rate and the passing process has its most effect in the intermediate densities. This could be anticipated since in the low and high densities the number of passing considerably reduces. The space-time diagrams of the model I reveal the discriminating effect of passing.
In model II (both symmetric and asymmetric), the maximum of $`J`$ occurs in different values of $`n`$ in simulation and analytical approach. Mean-field predicts a shift toward higher densities while in simulation a slight shift toward left is observed. We note that in the PCA based models the maximum of $`J`$ corresponds to a considerable left-shifted value of the density . In symmetric version of the model II, we observe an increment of the current with regard to the asymmetric version. In contrast to the asymmetric version, the maximum of $`J`$ in mean-field approach is higher than its value obtained through simulation. Although the current diagram (10) appears asymmetrically with respect to the density, the lane-changing diagram (13) is symmetric to a high accuracy.
Acknowledgements:
I would like to express my gratitude to V. Karimipour and G. Schütz, for their fruitful comments. I acknowledge the assistance given by A. Schadschneider, N.Hamadani, V.Shahrezaei, R. Sorfleet and in particular I highly appreciate R. Gerami for his valuable helps in the computer simulations.
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# How to uncross some modular metrics Supported by the Fields Institute for Research in Mathematical Sciences and by grant 97-01-00115 from the Russian Foundation of Basic Research.
## 1 Introduction
We deal with a variant of the multifacility location problem. In its setting, there are a finite metric space $`(T,\mu )`$, a finite set $`X`$, and a nonnegative function $`c`$ on the pairs of elements of $`TX`$. (The elements of $`T`$ are thought of as the points where the existing facilities are located, the elements of $`X`$ as new facilities, and $`c(x,y)`$ as a measure of communication between $`x`$ to $`y`$.) The objective is to place each new facility at a point of $`T`$ minimizing the sum of values $`c(x,y)\mu (x^{},y^{})`$, where $`x,y`$ range over the pairs of facilities and $`x^{},y^{}`$ are the points of $`T`$ where $`x,y`$ are placed. For a survey on location problems, see, e.g., .
This problem can be reformulated in terms of metric extensions. We start with some terminology and notation. A semimetric on a set $`S`$ is a function $`d:S\times S\text{R}_+`$ that establishes distances on the pairs of elements (points) of $`S`$ satisfying $`d(x,x)=0`$, $`d(x,y)=d(y,x)`$ and $`d(x,y)+d(y,z)d(x,z)`$, for all $`x,y,zS`$. We use notation $`xy`$ for an unordered pair $`\{x,y\}`$ and usually write $`d(xy)`$ instead of $`d(x,y)`$. The set of pairs $`xy`$ with $`xy`$ is denoted by $`E_S`$. When $`d(xy)>0`$ for all $`xyE_S`$, $`d`$ is a metric. We do not distinguish between the (semi)metric $`d`$ and the (semi)metric space $`(S,d)`$ and usually deal with only finite (semi)metric spaces. A special case is the path metric $`d^G`$ of a connected graph $`G`$, where $`d^G(xy)`$ is the minimum number of edges of a path in $`G`$ connecting nodes $`x`$ and $`y`$.
A semimetric $`m`$ on a set $`VS`$ is said to be an extension of $`d`$ if the restriction (submetric) of $`m`$ to $`S`$ is just $`d`$. Such an $`m`$ is called a 0-extension if the distance $`m(x,S)`$ from each point $`xV`$ to $`S`$ is zero, i.e., $`m(xs)=0`$ for some $`sS`$. Clearly each 0-extension $`m`$ is uniquely determined by the 0-distance sets $`X_s=\{xV:m(xs)=0\}`$, $`sS`$, and these sets give a partition of $`V`$ when $`d`$ is a metric.
The above problem is equivalent to the minimum 0-extension problem : Given a metric $`\mu `$ on a set $`T`$, a superset $`VT`$, and a function $`c:E_V\text{Z}_+`$,
1. Find a 0-extension $`m`$ of $`\mu `$ to $`V`$ with $`cm:=(c(e)m(e):eE_V)`$ minimum.
In this paper we extend earlier results on the complexity of ((1.1)) for fixed metrics $`\mu `$.
Two classes of metrics $`\mu `$ have been found for which ((1.1)) is solvable in polynomial time. One class consists of the metrics for which ((1.1)) becomes as easy as its linear programming relaxation. More precisely, let $`\tau =\tau (V,c,\mu )`$ denote the minimum $`cm`$ in ((1.1)), and let $`\tau ^{}=\tau ^{}(V,c,\mu )`$ denote the minimum $`cm`$ in the problem:
1. Find an extension $`m`$ of $`\mu `$ to $`V`$ with $`cm`$ minimum.
Then $`\tau \tau ^{}`$. A metric $`\mu `$ is called minimizable if $`\tau (V,c,\mu )=\tau ^{}(V,c,\mu )`$ holds for any $`V`$ and $`c`$. Since ((1.2)) is a linear program whose constraint matrix size is polynomial in $`|V|`$, ((1.2)) is solvable in strongly polynomial time. This easily implies that for every minimizable metric $`\mu `$, on optimal 0-extension in ((1.1)) can be found in strongly polynomial time as well. The following theorem characterizes the set of minimizable path metrics.
###### Theorem 1.1
For a graph $`H`$, the metric $`d^H`$ is minimizable if and only if $`H`$ is hereditary modular and orientable.
Recall that a metric $`\mu `$ on $`T`$ is modular if every three points $`s_0,s_1,s_2T`$ have a median, a node $`zT`$ satisfying $`\mu (s_iz)+\mu (zs_j)=\mu (s_is_j)`$ for all $`0i<j2`$. A graph $`H`$ is called modular if $`d^H`$ is modular, and hereditary modular if every isometric subgraph of $`H`$ is modular, where a subgraph (or circuit) $`H^{}=(T^{},U^{})`$ of $`H`$ is isometric if $`d^H^{}(st)=d^H(st)`$ for all $`s,tT^{}`$. Every modular graph is bipartite. A graph is called orientable if its edges can be oriented so that for any 4-circuit $`C=(v_0,e_1,v_1,\mathrm{},e_4,v_4=v_0)`$ and $`i=1,2`$, the edge $`e_i`$ is oriented from $`v_{i1}`$ to $`v_i`$ if and only if the opposite edge $`e_{i+2}`$ is oriented from $`v_{i+2}`$ to $`v_{i+1}`$; see Fig. 1(a). For example, every bipartite graph with at most five nodes is hereditary modular and orientable. The simplest hereditary modular but not orientable graph is the graph $`K_{3,3}^{}`$ obtained from $`K_{3,3}`$ by deleting one edge; see Fig. 1(b). Using terminology in , we refer to an orientable hereditary modular graph as a frame.
Theorem 1.1 is extended to general metrics using the notion of underlying graph of $`\mu `$. This is the least graph $`H(\mu )=(T,U(\mu ))`$ which enables us to restore $`\mu `$ if we know the distances of its edges. Formally, nodes $`x,yT`$ are adjacent in $`H(\mu )`$ if and only if no other node $`zT`$ lies between $`x`$ and $`y`$, i.e., satisfies $`\mu (xz)+\mu (zy)=\mu (xy)`$. This graph is modular if $`\mu `$ is modular .
###### Theorem 1.2
A metric $`\mu `$ is minimizable if and only if $`\mu `$ is modular and $`H(\mu )`$ is a frame.
Another tractable case involves median metrics, the metrics $`\mu `$ with precisely one median for each triplet of points. Chepoi showed that ((1.1)) with any median metric $`\mu `$ is solvable in strongly polynomial time. A simple alternative method, based on cut uncrossing techniques, is suggested in . Note that a minimizable metric need not be a median one, and vice versa. For example, $`d^{K_{2,3}}`$ is minimizable but not median, while the path metric of the (skeleton of the) cube is median but not minimizable (the cube is not hereditary modular as it contains an isometric 6-circuit).
In this paper we show the polynomial solvability for a class of modular metrics which includes the median ones as a special case. It uses the notion of orbit graphs that we now introduce. Given a modular graph $`H=(T,U)`$, two edges are called mates if they are opposite in some 4-circuit; when dealing with graphs with possible parallel edges, we refer to such edges as mates as well. Edges $`e,e^{}`$ of $`H`$ are called projective if there is a sequence $`e=e_0,e_1,\mathrm{},e_k=e^{}`$ of edges such that any two consecutive $`e_i,e_{i+1}`$ are mates; such a sequence is called projective too. A maximal set $`Q`$ of mutually projective edges is called an orbit. Define the orbit graph $`H_Q`$ to be $`H//(UQ)`$, where for a graph $`H^{}`$ and a subset $`Z`$ of its edges, $`H^{}//Z`$ denotes the graph obtained by contracting $`Z`$ (i.e., forming $`H^{}/Z`$) and then identifying parallel edges appeared.
The main result in this paper is the following.
###### Theorem 1.3
Let $`\mu `$ be a modular metric with underlying graph $`H=(T,U)`$, and let for each orbit $`Q`$ of $`H`$,
(i) the orbit graph $`H_Q`$ is a frame, and
(ii) $`H_Q`$ is isomorphic to some subgraph of the graph $`(T,Q)`$.
Then ((1.1)) can be solved in strongly polynomial time.
We shall explain later that each orbit graph of a frame is a frame, and each orbit graph of a median graph is $`K_2`$, which is a trivial case of frames. Since condition (ii) in Theorem 1.3 obviously holds when $`H_Q`$ is $`K_2`$, this theorem generalizes the above result for median metrics. On the other hand, the set of metrics $`\mu `$ in this theorem does not contain some minimizable metrics since there are frames $`H`$ for which (ii) is not valid. One can show that (ii) holds when each orbit graph is either $`K_2`$ or $`K_{2,r}`$ for $`r3`$, the simplest cases of frames with one orbit. Figure 2 illustrates the graph $`H`$ with three orbits, drawn by thin, dashed and bold lines, whose orbit graphs are $`H_1K_{2,3}`$, $`H_2K_2`$ and $`H_3K_2`$.
The proof of Theorem 1.3 will involve a number of reductions. One of them is to show that this theorem can be derived from Theorem 1.1 and Theorem 1.4 below that claims the existence of a retraction for certain graphs. Here a retraction of a bipartite graph $`K=(V(K),E(K))`$ onto its subgraph $`K^{}=(V(K^{}),E(K^{}))`$ is meant to be a mapping $`\gamma :V(K)V(K^{})`$ which is identical on $`V(K^{})`$ (i.e., $`\gamma (v)=v`$ for all $`vV(K^{})`$) and brings each edge of $`K`$ to an edge of $`K^{}`$ (i.e., $`\gamma (u)\gamma (v)E(K^{})`$ for all $`uvE(K)`$). Suppose $`K`$ is the Cartesian product $`H_1\times \mathrm{}\times H_k`$ of graphs $`H_i=(T_i,U_i)`$, $`i=1,\mathrm{},k`$, i.e., $`V(K)=T_1\times \mathrm{}\times T_k`$ and nodes $`(s_1,\mathrm{},s_k)`$ and $`(t_1,\mathrm{},t_k)`$ of $`K`$ are adjacent if and only if there is $`i\{1,\mathrm{},k\}`$ such that $`s_it_iU_i`$ and $`s_j=t_j`$ for $`ji`$. For a subgraph $`K^{}`$ of $`K`$ and $`i\{1,\mathrm{},k\}`$, an $`i`$-layer of $`K^{}`$ is a maximal subgraph of $`K^{}`$ induced by nodes $`(t_1,\mathrm{},t_k)`$ with $`t_1,\mathrm{},t_{i1},t_{i+1},\mathrm{},t_k`$ fixed.
###### Theorem 1.4
Let $`K`$ be the Cartesian product of frames $`H_i=(T_i,U_i)`$, $`i=1,\mathrm{},k`$. Let $`K^{}`$ be an isometric subgraph of $`K`$ such that $`K^{}`$ is modular and for each $`i=1,\mathrm{},k`$, some of the $`i`$-layers of $`K^{}`$ is isomorphic to $`H_i`$. Then there exists a retraction of $`K`$ onto $`K^{}`$.
(Note that $`K^{}`$ is not an absolute retract in general, i.e., $`K^{}`$ need not admit retraction of any bipartite graph which contains $`K^{}`$ as an isometric subgraph. Necessary and sufficient conditions on a bipartite graph to be an absolute retract are given in .) In our case, the role of graphs $`H_i`$ and $`K^{}`$ in Theorem 1.4 will play the graphs $`H_Q`$ and $`H`$ in Theorem 1.3, using the important observation that $`H`$ has a canonical isometric embedding in the Cartesian product $`K`$ of its orbit graphs. It turns out that Theorem 1.4 can be rather easily reduced to its special case with $`k=2`$; moreover, such a reduction takes place for arbitrary modular graphs $`H_1,\mathrm{},H_k`$. To show the existence of a retraction for this special case, with $`H_1,H_2`$ frames, is the core of the whole proof of Theorem 1.3. Such a retraction is just behind our “metric uncrossing operation”, an analogue of the cut uncrossing operation for 0-extensions of the corresponding orbit metrics (when both $`H_1,H_2`$ are $`K_2`$, the retraction is evident and it induces the uncrossing of two cuts, as we explain later).
Next we deal with intractable cases. When $`\mu =d^{K_p}`$, ((1.1)) turns into the minimum multiterminal (or multiway) cut problem, which is strongly NP-hard already for $`p=3`$ . That result has been generalized to a larger set of path metrics.
###### Theorem 1.5
For a fixed graph $`H`$, problem ((1.1)) with $`\mu =d^H`$ is strongly NP-hard if $`H`$ is non-modular or non-orientable.
We extend this theorem as follows.
###### Theorem 1.6
For a fixed rational-valued metric $`\mu `$, ((1.1)) is strongly NP-hard if $`\mu `$ is non-modular or if the underlying graph $`H(\mu )`$ is non-orientable.
The structure of this paper is as follows. Section 2 demonstrates some basic properties of modular metrics and graphs and their orbit graphs. Section 3 describes our approach to proving Theorem 1.3 and is aimed to explain why this theorem reduces to Theorem 1.4 with $`k=2`$. The desired retraction is constructed in Section 4, using combinatorial arguments and relying on some result concerning the tight spans of minimizable path metrics from . The construction also relies on a key lemma proved in Section 5. The proof of Theorem 1.6 is given in Section 6.
By technical reasons, in problems ((1.1)) and ((1.2)) we will sometimes admit $`\mu (st)=0`$ for distinct $`s,t`$ and may speak about minimizable semimetrics rather than metrics; this does not affect the problem area in essense. The sets of extensions and 0-extensions of a (semi)metric $`\mu `$ to a set $`V`$ are denoted by $`\mathrm{Ext}(\mu ,V)`$ and $`\mathrm{Ext}^0(\mu ,V)`$, respectively.
## 2 Modular metrics, modular graphs, and orbits
By a $`u`$$`v`$ path on a set $`V`$ we mean any sequence $`P=(x_0,x_1,\mathrm{},x_k)`$ of elements of $`V`$ with $`x_0=u`$ and $`x_k=v`$. For a semimetric $`m`$ on $`V`$, the $`m`$-length $`m(P)`$ of $`P`$ is $`m(x_0x_1)+\mathrm{}+m(x_{k1}x_k)`$, and $`P`$ is called shortest w.r.t. $`m`$, or $`m`$-shortest, if $`m(P)=m(uv)`$. If each pair $`e_i=x_{i1}x_i`$ is an edge of a graph $`G=(V,E)`$, then $`P=(x_0,e_1,x_1,\mathrm{},e_k,x_k)`$ is a path in $`G`$, and we say that $`P`$ is $`G`$-shortest if its length $`|P|:=k`$ is equal to $`d^G(uv)`$. When it is not confusing, we abbreviate $`P=x_0x_1\mathrm{}x_k`$. Given nonnegative lengths $`\mathrm{}(e)`$ of the edges $`eE`$, we denote by $`d^{G,\mathrm{}}(xy)`$ the minimum length $`\mathrm{}(P)=(\mathrm{}(x_{i1}x_i):i=1,\mathrm{},k)`$ of a path $`P=x_0x_1\mathrm{}x_k`$ connecting nodes $`x`$ and $`y`$ in $`G`$ (the path (semi)metric for $`G,\mathrm{}`$). From the definition of the underlying graph $`H(\mu )`$ of a metric $`\mu `$ it follows that $`\mu =d^{H(\mu ),\mathrm{}}`$ for the restriction $`\mathrm{}`$ of $`\mu `$ to the edges of $`H(\mu )`$.
Bandelt showed useful relations between modular graphs and metrics. They can be stated in terms of orbits as follows (cf. ).
1. For an orbit $`Q`$ of a modular graph $`H=(T,U)`$ and nodes $`u,vT`$, if $`P`$ is a shortest $`u`$$`v`$ path and $`P^{}`$ is a $`u`$$`v`$ path in $`H`$, then $`|PQ||P^{}Q|`$; in particular, $`|PQ|=|P^{}Q|`$ if both $`P,P^{}`$ are shortest.
1. For a modular metric $`\mu `$, the graph $`H(\mu )`$ is modular and $`\mu `$ is orbit-invariant, i.e., it is constant on the edges of each orbit of $`H(\mu )`$.
1. For a modular graph $`H=(T,U)`$ and an orbit-invariant function $`\mathrm{}:U\text{R}_+`$, the semimetric $`\mu =d^{H,\mathrm{}}`$ is modular, $`\mu (e)=\mathrm{}(e)`$ for all $`eU`$, and every $`H`$-shortest path is $`\mu `$-shortest; moreover, if, in addition, $`\mathrm{}`$ is positive, then $`H=H(\mu )`$, and the metrics $`d^H`$ and $`\mu `$ have the same sets of shortest paths.
Note that $`\mu `$ need not be modular when $`H(\mu )`$ is modular. (Properties ((2.2)) and ((2.3)) are easily derived from ((2.1)). The latter can be seen as follows (a sketch). Let $`w`$ be the node of $`P`$ following $`u`$. One may assume $`P^{}`$ is simple and all intermediate nodes $`x`$ of $`P^{}`$ are different from $`w`$. Since $`P`$ is shortest and $`H`$ is bipartite, some node $`x`$ of $`P^{}`$ satisfies $`d^H(wx)1=d^H(wy)=d^H(wz)`$, where $`y,z`$ are the neighbours of $`x`$ in $`P^{}`$. Take a median $`x^{}`$ of $`y,z,w`$. Then $`x^{}y`$ and $`x^{}z`$ are edges of $`H`$ projective to $`xz`$ and $`xy`$, respectively. Therefore, the path $`P^{\prime \prime }`$ obtained from $`P^{}`$ by replacing $`x`$ by $`x^{}`$ obeys $`|P^{\prime \prime }Q|=|P^{}Q|`$, and we can apply induction since the sum of distances from $`w`$ to the nodes of $`P^{\prime \prime }`$ is less than the corresponding sum for $`P^{}`$, in view of $`d^H(wx^{})=d^H(wx)2`$.)
By ((2.3)), every modular graph is the underlying graph for the class of modular metrics determined by positive orbit-invariant functions on its edges, and all these metrics have the same sets of shortest paths. This fact will often allow us to work with modular graphs rather than modular metrics.
Consider a modular graph $`H=(T,U)`$, and let $`Q_1,\mathrm{},Q_k`$ be the orbits of $`H`$. Let $`\chi ^S`$ denote the incidence vector of a subset $`SU`$, i.e. $`\chi ^S(e)=1`$ for $`eS`$, and 0 for $`eUS`$. Any modular metric $`\mu `$ with $`H(\mu )=H`$ is representable as
$$\mu =h_1\mu _1+\mathrm{}+h_k\mu _k,$$
(2.4)
where $`\mu _i=d^{H,\mathrm{}_i}`$ for $`\mathrm{}_i=\chi ^{Q_i}`$ and $`h_i=\mu (e)`$ for $`eQ_i`$ ($`h_i`$ is well-defined by ((2.2))). Indeed, for any $`s,tT`$, a shortest $`s`$$`t`$ path $`P`$ in $`H`$ is shortest for each of $`\mu ,\mu _1,\mathrm{},\mu _k`$, and $`\mu _i`$ coincides with $`\mathrm{}_i`$ on $`U`$, by ((2.3)). Therefore,
$$\mu (st)=\mu (P)=h_1\mathrm{}_1(P)+\mathrm{}+h_k\mathrm{}_k(P)=h_1\mu _1(st)+\mathrm{}+h_k\mu _k(st),$$
as required. When all $`h_i`$’s are ones, (2.4) is specified as
$$d^H=\mu _1+\mathrm{}+\mu _k.$$
(2.5)
Some properties of $`H`$ preserve under contraction of orbits. Let $`H^{}=(T^{},U^{})`$ be the graph $`H/Q_1`$. We identify the edges in $`UQ_1`$ with their images in $`H^{}`$ and denote by $`\phi (x)`$ (resp. $`\phi (P)`$) the image in $`H^{}`$ of a node $`x`$ (resp. a path $`P`$) of $`H`$. By ((2.3)) applied to the orbit-invariant function $`\mathrm{}=\chi ^{UQ_1}`$,
1. if $`P`$ is a shortest path of $`H`$, then $`\phi (P)`$ is a shortest path of $`H^{}`$.
Therefore, if $`v`$ is a median of nodes $`x,y,z`$ in $`H`$, then $`\phi (v)`$ is a median of $`\phi (x),\phi (y),\phi (z)`$ in $`H^{}`$. This implies that $`H^{}`$ is modular.
###### Statement 2.1
$`Q_2,\mathrm{},Q_k`$ are the orbits of $`H^{}`$.
Proof. Obviously, mates $`e,e^{}UQ_1`$ of $`H`$ remain mates in $`H^{}`$, i.e., they are either opposite in a 4-circuit or parallel. This implies that each set $`Q_i`$ ($`i>1`$) is entirely included in some orbit of $`H^{}`$. To see the reverse inclusion, consider a 4-circuit $`C=(v_0,e_1,v_1,\mathrm{},e_4,v_4=v_0)`$ of $`H^{}`$, and let $`L_j`$ denote the path $`(v_j,e_{j+1},v_{j+1},e_{j+2},v_{j+2})`$ for $`j=0,\mathrm{},3`$ (taking indices modulo 4). Each $`L_j`$ is a shortest path since $`H^{}`$ is bipartite (as being modular). Choose $`x_0\phi ^1(v_0)`$ and $`x_2\phi ^1(v_2)`$, and let $`P_0`$ and $`P_2`$ be two $`x_0`$$`x_2`$ paths of $`H`$ whose images in $`H^{}`$ are $`L_0`$ and the reverse to $`L_2`$, respectively. Let $`P`$ be a shortest $`x_0`$$`x_2`$ path in $`H`$. Then $`|\phi (P)|=|L_0|=|L_2|=2`$. This together with ((2.1)) (applied to $`P`$ and $`P^{}=P_0,P_2`$) implies $`|PQ_i|=|L_0Q_i|=|L_2Q_i|`$ for $`i=2,\mathrm{},k`$. Similarly, $`|L_1Q_i|=|L_3Q_i|`$ for each $`i`$. These equalities are possible only if each pair of mates in $`C`$ belongs to the same set $`Q_i`$. Similar arguments are applied to parallel edges of $`H^{}`$ (if any).
Repeatedly applying this statement to orbits of $`H`$, we obtain the following.
###### Corollary 2.2
For any $`I\{1,\mathrm{},k\}`$, the graph $`H/(_{iI}Q_i)`$ is modular and its orbits are the sets $`Q_j`$ for $`j\{1,\mathrm{},k\}I`$. In particular, each orbit graph $`H_Q`$ of a modular graph $`H=(T,U)`$ is modular and has only one orbit, which is obtained by identifying parallel edges in $`H/(UQ)`$.
Next we explain that each orbit graph of $`H(\mu )`$ is $`K_2`$ when $`\mu `$ is a median metric; this follows from properties of median graphs revealed by Mulder and Schrijver . Since $`\mu `$ and $`H(\mu )`$ have the same sets of shortest paths (by ((2.2)) and ((2.3))), a point $`v`$ is a median of points $`x,y,z`$ for $`\mu `$ if and only if $`v`$ is a median of this triplet for $`d^{H(\mu )}`$. So $`d^{H(\mu )}`$ is a median metric, which means that $`H(\mu )`$ is a median graph. It is shown in that
1. $`H=(T,U)`$ is a median graph if and only if $`d^H=\mu _1+\mathrm{}+\mu _k`$, where each $`\mu _i`$ is the cut metric corresponding to a bi-partition $`\{A_i,TA_i\}`$ of $`T`$ (i.e., $`\mu _i(st)=1`$ if $`|\{s,t\}A_i|=1`$, and 0 otherwise), and the family $`=\{A_1,\mathrm{},A_k,TA_1,\mathrm{},TA_k\}`$ satisfies the Helly property (i.e., any subfamily $`^{}`$ has a nonempty intersection provided that each two members of $`^{}`$ meet).
Let $`Q_i`$ be the set of edges of $`H`$ connecting $`A_i`$ and $`TA_i`$; clearly $`Q_1,\mathrm{},Q_k`$ are pairwise disjoint. These sets are precisely the orbits of $`H`$. Indeed, in view of $`d^H=\mu _1+\mathrm{}+\mu _k`$, a shortest path of $`H`$ is $`\mu _i`$-shortest for each $`i`$. This easily implies that: (i) the subgraphs of $`H`$ induced by $`A_i`$ and by $`TA_i`$ are connected, and (ii) $`Q_i`$ is a matching. ( shows the sharper property that $`H`$ is median if and only if $`H`$ is modular and has a cutset edge colouring.) Since $`Q_i`$ is simultaneously a cut and a matching, if $`e,e^{}`$ are mates in $`H`$ and $`eQ_i`$, then $`e^{}Q_i`$. So each orbit $`Q`$ of $`H`$ is included in some $`Q_i`$. Suppose $`QQ_i`$. Then the subgraph $`(T,UQ)`$ is connected, by (i) above, whence the semimetric $`\mu ^{}=d^{H,\mathrm{}}`$ for $`\mathrm{}=\chi ^Q`$ is identically zero. This is impossible because $`\mu ^{}`$ coincides with $`\mathrm{}`$ on $`U`$, by ((2.3)). Thus, $`Q_i`$ is an orbit. Now (i) implies that $`H/(UQ_i)`$ is a tuple of parallel edges, and we conclude that each orbit graph of $`H`$ is $`K_2`$.
As mentioned in the Introduction, our approach to solving problem ((1.1)) with a metric figured in Theorem 1.3 generalizes the cut uncrossing method for median metrics $`\mu `$. We now briefly describe that method, referring the reader for details to \[11, Sec. 5\].
Given a median metric $`\mu `$ on $`T`$, a set $`VT`$ and a function $`c:E_V\text{Z}_+`$, represent $`\mu `$ as in (2.4), where each $`\mu _i`$ is the cut metric corresponding to a bi-partition $`\{A_i,TA_i\}`$ of $`T`$ as in ((2.7)). For $`i=1,\mathrm{},k`$, find a bi-partition $`\{X_i,\overline{X}_i\}`$ of $`V`$ such that $`X_iT=A_i`$ and $`(c(xy):xX_i\ni ̸y)`$ is minimum (a minimum cut separating $`A_i`$ and $`TA_i`$). Let $`𝒳=\{X_1,\mathrm{},X_k,\overline{X}_1,\mathrm{},\overline{X}_k\}`$, and let $`m=h_1m_1+\mathrm{}+h_km_k`$, where $`m_i`$ is the cut metric on $`V`$ corresponding to $`\{X_i,\overline{X}_i\}`$. Choose a pair $`Y,Z𝒳`$ such that $`YZT=\mathrm{}`$ but $`Y,Z`$ meet, and make “uncrossing” by replacing $`Y,Z`$ in $`𝒳`$ by $`Y^{}=YZ`$ and $`Z^{}=ZY`$ (taking into account that $`\{Y^{},\overline{Y}^{}\}`$ induces a minimum cut separating $`YT`$ and $`\overline{Y}T`$, and $`\{Z^{},\overline{Z}^{}\}`$ induces a minimum cut separating $`ZT`$ and $`\overline{Z}T`$). Iterate until the current family $`𝒳^{}`$ has no such pair $`Y,Z`$, i.e., $`YZT=\mathrm{}`$ implies $`YZ`$. Using the Helly property for $``$ in ((2.7)), one can see that the corresponding metric $`m^{}=h_1m_1^{}+\mathrm{}+h_km_k^{}`$ is a 0-extension. Moreover, the fact that each $`m_i^{}`$ is induced by a minimum cut implies that $`m^{}`$ is optimal. One shows that the number of iterations does not exceed $`|T|^2|V|`$ (in fact, one can arrange a process consisting of only $`O(k^2)`$ uncrossing operations).
It turns out that the Helly property for median graphs exhibited in ((2.7)) is extended to general modular graphs. More precisely, for a modular graph $`H=(T,U)`$ with orbits $`Q_1,\mathrm{},Q_k`$, let $`H_i=(T_i,U_i)`$ stand for $`H_{Q_i}`$, and define $`\pi _i=\{A_i(t):tT_i\}`$ to be the partition of $`T`$ where each member $`A_i(t)`$ is the node set of the component of $`(T,UQ_i)`$ whose contraction creates the node $`t`$ of $`H_i`$. Each $`A_i(t)`$ is just the corresponding maximal 0-distance set of the metric $`\mu _i=d^{H,\mathrm{}_i}`$ as in (2.5). We assert that
1. the family $`\pi =\pi (H)`$ of subsets of $`T`$ occurring in $`\pi _1,\mathrm{},\pi _k`$ has the Helly property.
Indeed, each set $`A\pi `$ is convex, i.e., for any $`x,yA`$, each node on a shortest $`x`$$`y`$ path $`P`$ of $`H`$ belongs to $`A`$. To see this, assume $`A\pi _i`$. Then $`\mu _i(xy)=0`$, and therefore, $`\mathrm{}_i(P)=0`$ (by ((2.3))). So all nodes of $`P`$ are in $`A`$, as required. Now the result follows from the simple fact that the family $`\overline{\pi }`$ of convex node sets of an arbitrary modular graph has the Helly property. (This is shown by induction on $`n`$, considering a collection $`\pi ^{}=\{A^1,\mathrm{},A^n\}`$ of $`n3`$ members of $`\overline{\pi }`$ such that any $`n1`$ of them meet. For $`i=1,2,3`$, choose an element $`x_i`$ contained in all sets in $`\pi ^{}`$ except possibly $`A^i`$. Let $`z`$ be a median of $`x_1,x_2,x_3`$. For each $`A^j\pi ^{}`$, at least two of $`x_1,x_2,x_3`$ belong to $`A^j`$, hence $`zA^j`$ by the convexity. Thus, the members of $`\pi ^{}`$ have a common element.)
In conclusion of this section we show the hereditary property for orbit graphs of frames.
###### Statement 2.3
Let $`H=(T,U)`$ be a frame, and let $`Z`$ be the union of some orbits of $`H`$. Then $`H/Z`$ is a frame. In particular, each orbit graph of $`H`$ is a frame.
Proof. One can try to prove directly that the graph $`H/Z=:H^{}=(T^{},U^{})`$ is hereditary modular and orientable. We, however, can use Theorem 1.2 and standard compactness arguments to show that $`d^H^{}`$ is minimizable. Then $`H^{}`$ is a frame by Theorem 1.1.
More precisely, define the semimetric $`\mu `$ on $`T`$ to be $`d^{H,\mathrm{}}`$ for $`\mathrm{}=\chi ^{UZ}`$. Consider $`V^{}T^{}`$ and $`c^{}:E_V^{}\text{Z}_+`$. We have to show that $`\tau (V^{},c^{},d^H^{})=\tau ^{}(V^{},c^{},d^H^{})`$. Let $`V=V^{}T`$ (assuming $`V^{}T=T^{}`$) and define $`c(e)=c^{}(e)`$ for $`eE_V^{}`$, and $`c(e)=0`$ for $`eE_VE_V^{}`$. Clearly $`\tau (V,c,\mu )=\tau (V^{},c^{},d^H^{})`$ and $`\tau ^{}(V,c,\mu )=\tau ^{}(V^{},c^{},d^H^{})`$. So it is suffices to prove $`\tau (V,c,\mu )=\tau ^{}(V,c,\mu )`$.
To see the latter, consider the infinite sequence $`d_1,d_2,\mathrm{}`$ of approximations for $`\mu `$, where $`d_i`$ is $`d^{H,\rho _i}`$ with $`\rho _i(e)=1`$ for $`eUZ`$, and $`\rho _i(e)=1/i`$ for $`eZ`$. Since $`H`$ is modular and $`\rho _i`$ is positive and orbit-invariant, $`H=H(d_i)`$ for each $`i`$ by ((2.3)). So $`d_i`$ is minimizable (by Theorem 1.2), whence $`\tau (V,c,d_i)=\tau ^{}(V,c,d_i)`$. When $`i`$ grows, $`\tau (V,c,d_i)`$ tends to $`\tau (V,c,\mu )`$ (since the number of partitions of $`V`$ is finite). Also $`\tau ^{}(V,c,d_i)`$ tends to $`\tau ^{}(V,c,\mu )`$, because of the obvious fact that for any $`m\mathrm{Ext}(\mu ,V)`$, there exists $`m^{}\mathrm{Ext}(d_i,V)`$ such that $`|m^{}(e)m(e)||V|/i`$ for each $`eE_V`$. Thus, $`\tau (V,c,\mu )=\tau ^{}(V,c,\mu )`$, as required.
## 3 Reduction to the case of two orbits, and uncrossing method
In this section we describe our approach to proving Theorem 1.3. A majority of arguments below are applicable to general modular metrics, and unless explicitly said otherwise, we assume that $`\mu `$ is an arbitrary modular metric on a set $`T`$.
Let $`H=(T,U)`$ be the underlying graph $`H(\mu )`$ of $`\mu `$ with orbits $`Q_1,\mathrm{},Q_k`$. As before, for $`i=1,\mathrm{},k`$, $`H_i=(T_i,U_i)`$ stands for $`H_{Q_i}`$, $`\mathrm{}_i`$ for $`\chi ^{Q_i}`$, $`\mu _i`$ for $`d^{H,\mathrm{}_i}`$, and $`\pi _i=\{A_i(t):tT_i\}`$ for the corresponding partition of $`T`$ defined in the previous section. We formally identify each $`tT_i`$ with some element of $`A_i(t)`$, which allows us to speak of $`\mu _i`$ as a 0-extension of $`d^{H_i}`$ to $`T`$.
For the given $`\mu `$, consider an instance of the minimum 0-extension problem with $`VT`$ and $`c:E_V\text{Z}_+`$. By (2.4), any 0-extension $`m`$ of $`\mu `$ to $`V`$ is representable as
$$m=h_1m_1+\mathrm{}+h_km_k,$$
(3.1)
where each $`m_i`$ is the 0-extension of $`\mu _i`$ to $`V`$, defined by
1. $`m_i(xy)=\mu _i(st)`$ for $`x,yV`$ and $`s,tT`$ with $`m(xs)=m(yt)=0`$.
Then $`cm=c(h_1m_1)+\mathrm{}+c(h_km_k)`$ and $`cm_i\tau (V,c,\mu _i)`$ for each $`i`$. Taking as $`m`$ an optimal 0-extension for $`V,c,\mu `$, we conclude that
$$\tau (V,c,\mu )h_1\tau (V,c,\mu _1)+\mathrm{}+h_k\tau (V,c,\mu _k).$$
(3.3)
In particular, this is valid for $`h_1=\mathrm{}=h_k=1`$ and $`\mu =d^H`$. We say that $`H`$ is orbit-additive if
$$\tau (V,c,d^H)=\tau (V,c,\mu _1)+\mathrm{}+\tau (V,c,\mu _k)$$
(3.4)
holds for any $`V`$ and $`c`$. Such an $`H`$ has a sharper property.
###### Statement 3.1
Let $`H`$ be orbit-additive. Then for any numbers $`h_1,\mathrm{},h_k0`$, the semimetric $`\mu =d^{H,\mathrm{}}`$ with $`\mathrm{}=h_1\mathrm{}_1+\mathrm{}+h_k\mathrm{}_k`$ satisfies
$$\tau (V,c,\mu )=h_1\tau (V,c,\mu _1)+\mathrm{}+h_k\tau (V,c,\mu _k).$$
(3.5)
Moreover, if $`m`$ is an optimal 0-extension for $`V,c,d^H`$ and $`m_1,\mathrm{},m_k`$ are defined as in ((3.2)), then $`m^{}=h_1m_1+\mathrm{}+h_km_k`$ is an optimal 0-extension for $`V,c,\mu `$.
Proof. Since $`\tau (V,c,d^H)=cm=cm_1+\mathrm{}+cm_k`$, (3.4) implies $`cm_i=\tau (V,c,\mu _i)`$ for each $`i`$. Clearly $`m^{}\mathrm{Ext}^0(\mu ,V)`$. Therefore, $`\tau (V,c,\mu )cm^{}=h_1\tau (V,c,\mu _1)+\mathrm{}+h_k\tau (V,c,\mu _k)`$, yielding $`\tau (V,c,\mu )=cm^{}`$ and (3.5), in view of (3.3).
Because of (3.5), problem ((1.1)) for a metric $`\mu `$ whose underlying graph $`H`$ is orbit-additive becomes as easy as that for the path metrics of orbit graphs of $`H`$. Indeed, to compute $`\tau (V,c,\mu )`$ is reduced to finding the numbers $`\tau (V,c,\mu _i)`$. Moreover, once there is a subroutine to compute $`\tau (V^{},c^{},\mu )`$ for arbitrary $`V^{},c^{}`$, we can find an optimal 0-extension for the given $`\mu ,V,c`$ by applying this subroutine $`O(|T||V|)`$ times (similarly to the case of minimizable metrics $`\mu `$, mentioned in the Introduction).
In turn, $`\tau (V,c,\mu _i)`$ is equal to $`\tau (V_i,c_i,d^{H_i})`$, where $`V_i`$ and $`c_i`$ arise by shrinking the sets $`A_i(t)`$ in the partition $`\pi _i`$ of $`T`$ to the nodes $`tT_i`$. Formally, $`V_i=(VT)T_i`$, $`c_i(xy)=c(xy)`$ for $`x,yVT`$, $`c_i(xt)=c(\{x\},A_i(t))`$ for $`xVT,tT_i`$, and $`c_i(st)=c(A_i(s),A_i(t))`$ for $`s,tT_i`$, where $`c(A,B)`$ denotes $`(c(xy):xA,yB)`$ for $`A,BV`$.
In light of the above discussion, Theorem 1.3 would follow from Theorem 1.1 and the property that if $`H`$ is as in the hypotheses of Theorem 1.3, then
1. $`H`$ is orbit-additive.
Remark 1. The property of being orbit-additive is immediate in two cases of modular graphs $`H`$. Given $`V,c`$, let $`m_i`$ be an optimal 0-extension for $`V,c,\mu _i`$, and let $`m=m_1+\mathrm{}+m_k`$. By (2.5), $`m\mathrm{Ext}(d^H,V)`$. (i) If $`H`$ is a frame, then (3.4) holds since $`\tau (V,c,d^H)=\tau ^{}(V,c,d^H)cm=\tau (V,c,\mu _1)+\mathrm{}+\tau (V,c,\mu _k)\tau (V,c,d^H)`$. (ii) If $`H`$ is isomorphic to the Cartesian product of $`H_1,\mathrm{},H_k`$, then $`m`$ is already a 0-extension of $`d^H`$, yielding (3.4); cf. .
We further explain that ((3.6)) would follow from the existence of a retraction onto $`H`$ of the Cartesian product $`K=K(H)`$ of the orbit graphs $`H_1,\mathrm{},H_k`$ of $`H`$ (see the Introduction for needed definitions). We will use notation $`z_i`$ for $`i`$th coordinate (component) of a point $`zV(K)`$. Since each $`H_i`$ is bipartite, so is $`K`$. For $`vT`$, define
1. $`\varphi (v)`$ to be $`zV(K)`$ such that $`vA_i(z_i)`$ for $`i=1,\mathrm{},k`$.
###### Statement 3.2
For any $`u,vT`$, $`d^H(uv)=d^K(\varphi (u)\varphi (v))`$.
Proof. Let $`\varphi (u)=s`$ and $`\varphi (v)=t`$. We have $`d^K(st)=d^{H_1}(s_1t_1)+\mathrm{}+d^{H_k}(s_kt_k)`$. Consider a shortest $`u`$$`v`$ path $`P`$ in $`H`$, and for $`i=1,\mathrm{},k`$, let $`P_i`$ be the image of $`P`$ in $`H_i`$. Then $`|P|=|P_1|+\mathrm{}+|P_k|`$, and each $`P_i`$ is a shortest path, by ((2.6)). By ((3.7)), $`uA_i(s_i)`$ and $`vA_i(t_i)`$, so $`s_i,t_i`$ are the ends of $`P_i`$ and $`|P_i|=d^{H_i}(s_it_i)`$. Therefore, $`|P|=d^K(st)`$.
Thus, $`\varphi `$ induces an isometric embedding of $`H`$ into $`K`$, called the canonical embedding of $`H`$. We extend $`\varphi `$ to the edges of $`H`$ and, when no confusion can arise, identify $`H`$ with the subgraph $`\varphi (H)`$ of $`K`$. In particular, $`\varphi `$ is injective; in other words,
1. for $`zV(K)`$, the subset $`A_1(z_1)\mathrm{}A_k(z_k)`$ of $`T`$ consists of a single element (namely, $`\varphi ^1(z)`$) if $`z\varphi (T)`$, and is empty otherwise.
An elementary property of a retraction of a (bipartite) graph $`G=(V,E)`$ onto its subgraph $`G^{}=(V^{},E^{})`$ is that $`\gamma `$ turns every path of $`G`$ into a path of $`G^{}`$. This implies that $`d^G(xy)d^G^{}(\gamma (x)\gamma (y))`$ is a nonnegative even integer for any $`x,yV`$. Therefore, $`\gamma `$ is non-expansive (does not increase the distances) and preserves the distance parity.
###### Statement 3.3
A modular graph $`H`$ is orbit-additive if there exists a retraction of $`K=K(H)`$ onto $`H`$.
Proof. Given $`V,c`$, for each $`i=1,\mathrm{},k,`$ take an optimal 0-extension $`m_i`$ for $`V,c,\mu _i`$, and form the extension $`m=m_1+\mathrm{}+m_k`$ of $`d^H`$ to $`V`$. Assuming there exists a retraction $`\gamma `$ of $`K`$ onto $`H`$, we construct a 0-extension $`m^{}`$ of $`d^H`$ to $`V`$ such that $`m^{}m`$. This will imply (3.4) since $`\tau (V,c,\mu )cm^{}cm`$ and $`cm=\tau (V,c,\mu _1)+\mathrm{}+\tau (V,c,\mu _k)`$. For $`zV(K)`$, define
$`X_i(z_i)`$ $`=`$ $`\{xV:m_i(xv)=0\text{some}vA_i(z_i)\},i=1,\mathrm{},k;`$
$`X_z`$ $`=`$ $`X_1(z_1)\mathrm{}X_k(z_k).`$ (3.9)
The mapping $`\omega :VV(K)`$, defined by $`\omega (x)=z`$ for $`xX_z`$, isometrically embeds $`(V,m)`$ in $`(V(K),d^K)`$. Indeed, for $`xX_z`$ and $`yX_z^{}`$, we have
$$m(xy)=m_1(xy)+\mathrm{}+m_k(xy)=d^{H_1}(z_1z_1^{})+\mathrm{}+d^{H_k}(z_kz_k^{})=d^K(zz^{}).$$
Also $`\omega (v)=v`$ for each $`vT`$ (cf. ((3.7))), i.e., $`\omega `$ is identical on the node set of the graph $`H`$ embedded in $`K`$ by $`\varphi `$. The sets $`X_z`$ give a partition of $`V`$, and if it happens that for each nonempty set $`X_z`$, the set $`A_1(z_1)\mathrm{}A_k(z_k)`$ is nonempty too (thus consisting of a single node, by ((3.8))), then $`m`$ is already a 0-extension. In general, define the semimetric $`m^{}`$ on $`V`$ by
$$m^{}(xy)=d^H(\gamma (\omega (x))\gamma (\omega (y)))\text{for}x,yV.$$
Then $`m^{}`$ is a 0-extension of $`d^H`$ (corresponding to the partition $`\{\omega ^1\gamma ^1(t):tT\}`$). Now the fact that $`\gamma `$ is non-expansive while $`\omega `$ is isometric implies $`m^{}m`$, as required.
One can see that for each orbit $`Q_i`$, the components of the graph $`(T,Q_i)`$ are just the $`i`$-layers of $`H`$ (canonically embedded in $`K`$ by $`\varphi `$). Thus, condition (ii) in Theorem 1.3 says that each orbit graph $`H_i`$ is isomorphic to some of the $`i`$-layers of $`H`$, and now summing up the above reasonings, we conclude that Theorem 1.3 is implied by Theorem 1.4.
So it remains to prove Theorem 1.4. For convenience we denote $`K^{}`$ by $`H=(T,U)`$. Note that now any graph $`H_i`$ may have more than one orbit, but this is not important for us. First of all we explain that it suffices to consider the case $`k=2`$ (in the reduction below we only use the fact that each $`H_i`$ is modular rather than $`H_i`$ is a frame).
Let $`1i<jk`$ and $`K_{ij}=H_i\times H_j`$. Define $`H_{ij}=(T_{ij},U_{ij})`$ to be the projection of $`H`$ to $`K_{ij}`$, i.e., $`T_{ij}=\{(z_i,z_j):zT\}`$ and $`U_{ij}=\{(z_i,z_j)(z_i^{},z_j^{}):zz^{}U,z_p=z_p^{}`$ for $`pi,j\}`$. (When $`H`$ is as in Theorem 1.3, $`H_{ij}`$ is isomorphic to the “two-orbit graph” $`H//(UQ_iQ_j)`$.) Suppose a retraction $`\gamma _{ij}`$ of $`K_{ij}`$ onto $`H_{ij}`$ exists for each pair $`i,j`$. Define the mapping $`\psi _{ij}:V(K)V(K)`$ by $`\psi _{ij}(z)=z^{}`$, where $`(z_i^{},z_j^{})=\gamma _{ij}(z_i,z_j)`$ and $`z_p^{}=z_p`$ for $`pi,j`$. Clearly $`\psi _{ij}`$ is identical on $`T`$ and brings every edge of $`K`$ to an edge. Then the desired retraction $`\gamma `$ of $`K`$ onto $`H`$ is devised by successively applying transformations $`\psi _{ij}`$, as follows.
At the first step, set $`W_1:=V(K)`$ and choose a pair $`i,j`$ such that there is a point $`zW_1`$ with $`(z_i,z_j)T_{ij}`$. Set $`\alpha _1:=\psi _{ij}`$ and reduce $`W_1`$ to $`W_2:=\alpha _1(W_1)`$. Note that $`\alpha `$ decreases at least one distance, namely, for $`u=\alpha _1(z)`$, we have $`\alpha _1(u)=u`$, so $`d^K(zu)>d^K(\alpha _1(z)\alpha _1(u))=0`$. Similarly, at each step $`q`$, we choose $`i^{},j^{}`$ with $`(v_i^{},v_j^{})T_{i^{}j^{}}`$ for some $`vW_q`$, set $`\alpha _q:=\psi _{i^{}j^{}}`$ and reduce $`W_q`$ to $`W_{q+1}:=\alpha _q(W_q)`$, and so on. Since each transformation is non-expansive and brings some pair of points of the current set $`W`$ to closer points, the process is finite. It terminates when, after $`N`$ steps, for any $`zW_{N+1}`$, each pair $`(z_i,z_j)`$ is already in $`T_{ij}`$. Let $`\gamma =\alpha _N\alpha _{N1}\mathrm{}\alpha _1`$. Then $`\gamma `$ is identical on $`T`$, brings every edge to an edge and maps $`V(K)`$ to $`W_{N+1}`$. To conclude that $`\gamma `$ is a retraction of $`K`$ onto $`H`$, we have to show that $`W_{N+1}=T`$.
###### Statement 3.4
Let $`z`$ be a point in $`V(K)`$ such that $`(z_i,z_j)T_{ij}`$ for all $`0i<jk`$. Then $`z`$ is in $`H`$.
Proof. For each $`p=1,\mathrm{},k`$, the set $`B_p:=\{tT:t_p=z_p\}`$ is convex in $`H`$ (but not necessarily in $`K`$!). Indeed, if $`u,vB_p`$ and $`P`$ is a shortest $`u`$$`v`$ path in $`H`$, then $`P`$ is shortest in $`K`$ (since $`H`$ is isometric). Therefore, $`u_p=v_p=z_p`$ implies $`w_p=z_p`$ for any node $`w`$ on $`P`$, whence $`wB_p`$.
We know that the family of convex sets of a modular graph has the Helly property. The inclusion $`(z_i,z_j)T_{ij}`$ means that the sets $`B_i`$ and $`B_j`$ meet. Therefore, $`B_1,\mathrm{},B_k`$ have a common element $`z^{}T`$. Clearly $`z^{}=z`$.
Thus, it suffices to prove Theorem 1.4 for $`k=2`$. The desired retraction will be constructed in the next section.
Remark 2. The above arguments prompt a method to solve ((1.1)) with $`\mu `$ as in Theorem 1.3 in which each particular problem concerning $`\mu _i`$ is solved only once (so the method looks more efficient than that described after the proof of Statement 3.1). More precisely, given $`V,c`$, find an optimal 0-extension $`m_i`$ for each $`i=1,\mathrm{},k`$. This gives the family $`𝒳`$ of sets $`X_i(z_i)`$ as in (3.9), and we can select, in polynomial time, the set $`𝒱`$ consisting of all points $`zK(V)`$ with $`X_z\mathrm{}`$. Starting with $`𝒱_1=𝒱`$, at each, $`q`$th, iteration, we examine the current set $`𝒱_q`$ to find $`z𝒱`$ with $`(z_i,z_j)T_{ij}`$ for some $`i,j`$. If such a $`z`$ exists and is chosen, we set $`\alpha _q:=\psi _{ij}`$, reduce $`𝒱_q`$ to $`𝒱_{q+1}:=\alpha _q(𝒱_q)`$ (which changes $`𝒳`$) and continue the process. Otherwise $`𝒱_q=T`$, by Statement 3.4, and the partition $`\{Y_t:tT\}`$ of $`V`$ into the corresponding 0-distance sets induces an optimal 0-extension for $`V,c,d^H`$ (and therefore, for $`V,c,\mu `$, by Statement 3.1), where $`Y_t`$ is the union of sets $`X_z`$ for $`z𝒱`$ such that $`\alpha _{q1}\mathrm{}\alpha _1(z)=t`$. Since each transformation moves some point of the current set $`𝒱`$ closer to $`T`$, the number of iterations is $`O(|T|^2|V|)`$.
Remark 3. The above transformation of $`𝒳`$ induced by the retraction $`\gamma _{ij}`$ can be thought of as an analogue of the cut uncrossing operation for median metrics (reviewed in Section 2), thus justifying the term “uncrossing” used in a more general context in this paper. Recall that each orbit graph of a median graph $`H`$ is $`K_2`$, and therefore, each “two-orbit graph” $`H_{ij}`$ is isomorphic either to $`K_2\times K_2`$ or to the path $`P=xyz`$ of length two, as drawn in Fig. 3. When $`H_{ij}P`$, the (unique) retraction $`\gamma =\gamma _{ij}`$ brings the point $`(x,z)`$ of $`H_i\times H_j`$ not in $`H_{ij}`$ to $`y`$. This retraction is just behind the uncrossing operation on the corresponding cuts in that method.
## 4 Retraction
In this and next sections we prove Theorem 1.4 with $`k=2`$, using notation, conventions and results from Sections 2 and 3. One may assume $`KH`$. We will essentially use the condition in the theorem that $`H`$ includes a subgraph (“row-layer”) of the form $`H_1\times s_2`$ and a subgraph (“column-layer”) of the form $`s_1\times H_2`$ for some $`s_1T_1`$ and $`s_2T_2`$, i.e.,
1. for any $`uT_1`$ and $`vT_2`$, $`(u,s_2)T`$ and $`(s_1,v)T`$.
We fix such $`s_1,s_2`$ and call the node $`s=(s_1,s_2)`$ the origin of $`K`$.
In the proof below we everywhere admit that $`H_1,H_2`$ are arbitrary modular graphs until ((4.9)) where the assumption that $`H_1,H_2`$ are frames is essential. We abbreviate $`d^K,d^{H_1},d^{H_2}`$ as $`d,d_1,d_2`$, respectively. The interval $`\{vV(K):d(xv)+d(vy)=d(xy)\}`$ of nodes (points) $`x,y`$ of $`K`$ is denoted by $`I(x,y)=I(y,x)`$. We denote by $`J(x)`$ and $`r(x)`$ the interval $`I(x,s)`$ and the distance $`d(xs)`$, called the principal interval and the rank of $`x`$, respectively. $`M(x,y,z)`$ denotes the set of medians of points $`x,y,zV(K)`$. For $`i=1,2`$, $`I_i(x_i,y_i)`$, $`J_i(x_i)`$, $`r_i(x_i)`$, and $`M_i(x_i,y_i,z_i)`$ stand for the corresponding objects concerning the graph $`H_i`$. Then $`I(x,y)=I_1(x_1,y_1)\times I_2(x_2,y_2)`$, $`J(x)=J_1(x_1)\times J_2(x_2)`$, $`r(x)=r_1(x_1)+r_2(x_2)`$, and $`M(x,y,z)=M_1(x_1,y_1,z_1)\times M_2(x_2,y_2,z_2)`$ (as being immediate consequences from the equality $`d(uv)=d_1(u_1v_1)+d_2(u_2v_2)`$ for any $`u,vV(K)`$). The latter correspondence between medians in $`K,H_1,H_2`$ implies the following elementary property, which will be often used later on:
1. for $`x,y,zV(K)`$ and $`i\{1,2\}`$, if $`z_iI_i(x_i,y_i)`$, then $`z_i=v_i`$ for each median $`vM(x,y,z)`$; in particular, $`x_i=z_i`$ implies $`v_i=x_i`$.
The modularity of $`H`$ implies that
1. for each $`uT_1`$, the set $`Z(u):=\{vT_2:(u,v)T\}`$ is convex in $`H_2`$, and similarly for each $`vT_2`$, the set $`\{uT_1:(u,v)T\}`$ is convex in $`H_1`$
(cf. the proof of Statement 3.4). Indeed, for $`v,wZ(u)`$ and $`v^{}I_2(v,w)`$, consider the nodes $`x=(u,v)`$, $`y=(u,w)`$ and $`z=(s_1,v^{})`$ of $`H`$ (where $`z`$ is in $`T`$ by ((4.1))). These nodes have a median $`q`$ in $`H`$. Then $`q_1=u`$ and $`q_2=v^{}`$ (cf. ((4.2))). Hence, $`v^{}Z(u)`$. It follows from ((4.3)) that
$$J(t)T\text{for all}tT.$$
(4.4)
(However, the whole set $`T`$ is not convex in $`K`$ unless $`H=K`$.)
The mapping (retraction) $`\gamma `$ that we wish to construct will be some kind of reflection of points in $`V(K)T`$ with respect to their closest points in $`H`$. Consider a point $`xV(K)`$. Define the excess $`\mathrm{\Delta }^x`$ to be the distance $`d(x,T)`$ from $`x`$ to $`T`$, i.e., $`\mathrm{\Delta }^x=\mathrm{min}\{d(xt):tT\}`$, and define $`N(x)`$ to be the set of points $`tT`$ with $`d(xt)=\mathrm{\Delta }^x`$. In particular, $`\mathrm{\Delta }^xr_i(x_i)`$ for $`i=1,2`$ (since $`(x_1,s_2),(s_1,x_2)T`$), and $`\mathrm{\Delta }^x=0`$ if and only if $`xT`$.
###### Statement 4.1
$`N(x)J(x)`$.
Proof. Let $`tN(x)`$. The points $`x^{}=(x_1,s_2)`$, $`x^{\prime \prime }=(s_1,x_2)`$ and $`t`$ are in $`T`$, so they have a median $`q`$ in $`T`$ as well. Then $`q_1M_1(x_1,s_1,t_1)`$ and $`q_2M_2(s_2,x_2,t_2)`$. Therefore, $`q_1`$ belongs to both $`J_1(x_1)`$ and $`I_1(x_1,t_1)`$, and $`q_2`$ belongs to both $`J_2(x_2)`$ and $`I_2(x_2,t_2)`$, which means that $`qJ(x)`$ and $`qI(x,t)`$. Now $`d(xq)\mathrm{\Delta }^x=d(xt)`$ implies $`q=t`$.
By this statement, the rank $`r(t)`$ is equal to the same number $`r(x)\mathrm{\Delta }^x`$ for all $`tN(x)`$. Note that for any $`x,yV(K)`$, $`|\mathrm{\Delta }^x\mathrm{\Delta }^y|=|d(x,T)d(y,T)|d(xy)`$. Therefore,
$$|\mathrm{\Delta }^x\mathrm{\Delta }^y|1\text{for each edge}xyE(K).$$
(4.5)
We partition $`E(K)`$ into the sets $`E_1=\{xy:x_2=y_2\}`$ and $`E_2=\{xy:x_1=y_1\}`$, and for $`i=1,2`$, define
$$E_i^==\{xyE_i:\mathrm{\Delta }^x=\mathrm{\Delta }^y\}\text{and}E_i^{}=E_iE_i^=.$$
(4.6)
The desired retraction is devised by use of certain 0-extensions of metrics $`d_1`$ and $`d_2`$. First we introduce the auxiliary graphs $`G_1=(𝒱_1,_1)`$ and $`G_2=(𝒱_2,_2)`$, as follows. For $`i=1,2`$, let $`𝒜_i`$ be the set of pairs $`tt_i=\{t,t_i\}`$ for $`tT`$, and $`_i`$ the set of pairs $`xs_i=\{x,s_i\}`$ for $`xV(K)`$. Then $`G_i`$ is the (disjoint) union of the graphs $`H_i`$ and $`K`$ to which the pairs from $`𝒜_i_i`$ are added as edges, i.e.,
$$𝒱_i=T_iV(K)\text{and}_i=U_iE(K)𝒜_i_i.$$
The edges $`e`$ of $`G`$ are endowed with the lengths $`\delta _i(e)`$ defined by
$`\delta _i(e)`$ $`=`$ $`1\text{for}eU_iE_i^=E_{3i}^{},`$
$`=`$ $`0\text{for}eE_i^{}E_{3i}^=𝒜_i,`$
$`=`$ $`r_i(x_i)\mathrm{\Delta }^x\text{for}e=xs_i_i.`$
We say that a semimetric $`m`$ on a set $`V`$ is cyclically even if $`m(xy)+m(yz)+m(zx)`$ is an even integer for all $`x,y,zV`$ (equivalently: the $`m`$-length of any cycle on $`V`$ is even). All values of such an $`m`$ are integers since $`m(xy)+m(yx)+m(xx)=2m(xy)2\text{Z}`$.
###### Lemma 4.2
For $`i=1,2`$, define $`m_i=d^{G_i,\delta _i}`$. Then: (i) $`m_i`$ is an extension of $`d_i`$ to $`𝒱_i`$, and (ii) $`m_i`$ is cyclically even and coincides with $`\delta _i`$ on $`_i`$.
This lemma (the keystone in our arguments) will be proved later, and now we explain how it help us to construct the desired mapping $`\gamma `$. We apply some results from and .
More precisely, for a metric $`\mu ^{}`$ on a set $`T^{}`$, an extension $`m^{}`$ of $`\mu ^{}`$ to $`VT^{}`$ is called tight if there exists no $`m^{\prime \prime }(\mu ^{},V)\{m^{}\}`$ such that $`m^{\prime \prime }m^{}`$; equivalently: $`m^{}`$ has no loose pair $`x,y`$, i.e, for any $`x,yV`$, the path $`(u,x,y,v)`$ on $`V`$ is $`m^{}`$-shortest for at least one pair $`u,vT^{}`$. It is shown in \[13, Sec.5\] that for any cyclically even metric $`\mu ^{}`$,
1. if $`m\mathrm{Ext}(\mu ^{},V)`$ is cyclically even, then there exists $`m^{}\mathrm{Ext}(\mu ^{},V)`$ such that $`m^{}`$ is cyclically even and tight, $`m^{}(e)m(e)`$ for all $`eE_V`$, and $`m^{}(e)=m(e)`$ whenever $`m(e)1`$.
(Such an $`m^{}`$ is constructed by the following process. If there is no loose pair $`x,yV`$ with $`m(xy)2`$, then one easily shows that there is no loose pair at all, i.e., $`m`$ is already tight. Otherwise choose such a pair $`x,y`$, and let $`m^{}:=d^{K_V,\mathrm{}}`$, where $`\mathrm{}(xy):=m(xy)2`$ and $`\mathrm{}(e):=m(e)`$ for $`eE_V\{xy\}`$. Then $`m^{}`$ is a cyclically even extension of $`\mu ^{}`$. Update $`m:=m^{}`$ and iterate.)
Next, the proof of the “if” part of Theorem 1.1 in relies of an explicit construction of the so-called tight span of a frame, which in turn is based on the following result (Claim 5 in Section 4 there):
1. if $`H^{}=(T^{},U^{})`$ is a frame and $`m`$ is a tight extension of $`d^H^{}`$ to $`VT^{}`$, then each point $`xV`$ satisfies at least one of the following:
* $`m(xt)=0`$ for some node $`tT^{}`$;
* $`m(ux)+m(xv)=1`$ for some edge $`uvU^{}`$;
* $`m(v_0x)+m(xv_2)=m(v_1x)+m(xv_3)=2`$ for some 4-circuit $`C=v_0v_1v_2v_3v_0`$ of $`H^{}`$.
Using ((4.8)) and ((4.9)), we argue as follows. For $`i=1,2`$, let $`m_i`$ be as in Lemma 4.2, and let $`m_i^{}m_i`$ be a cyclically even tight extension of $`d_i`$ as in ((4.8)). Then
$$m_i^{}(e)=\delta _i(e)\text{for}e_i_i,\text{and}m_i^{}(e)\delta _i(e)\text{for}e_i.$$
(4.10)
Moreover, in view of ((4.9)), for each $`x𝒱_i`$, there exists $`tT_i`$ with $`m_i^{}(tx)=0`$. This is immediate in cases (i) and (ii) of ((4.9)). And if we are in case (iii) (with $`m=m_i^{}`$) and if $`C=v_0v_1v_2v_3v_0`$ is the corresponding 4-circuit for $`x`$, then $`\alpha _j:=m_i^{}(v_jx)>0`$ for $`j=0,1,2,3`$ would imply $`\alpha _j=1`$ for each $`j`$. Then $`m_i^{}(v_0v_1)+\alpha _0+\alpha _1=1+1+1=3`$, contrary to the fact that $`m_i^{}`$ is cyclically even. Thus, $`m_i^{}`$ is a 0-extension of $`d_i`$ to $`𝒱_i`$.
Now for $`xV(K)`$, define $`\gamma (x)`$ to be the point $`(\gamma _1(x),\gamma _2(x))`$, where $`\gamma _i(x)`$ is the node $`vT_i`$ with $`m_i^{}(xv)=0`$.
###### Statement 4.3
$`\gamma `$ is the retraction of $`K`$ onto $`H`$.
Proof. For each $`tT`$, $`m_i^{}(tt_i)=0`$ (since $`\delta _i`$ is zero on $`𝒜_i`$, by (4)), so $`\gamma `$ is identical on $`T`$.
To see $`\gamma (V(K))T`$, consider $`xV(K)`$, and let $`x^{}=\gamma (x)`$ and $`tN(x)`$. Let $`P=z^0z^1\mathrm{}z^k`$ ($`k=\mathrm{\Delta }^x`$) be a shortest $`t`$$`x`$ path in $`K`$. Then for $`j=0,\mathrm{},k1`$, one has $`tN(z^j)`$ and $`\mathrm{\Delta }_j:=\mathrm{\Delta }^{z^j}=j`$, whence $`\mathrm{\Delta }_{j+1}\mathrm{\Delta }_j=1`$ and $`z^jz^{j+1}E_1^{}E_2^{}`$, cf. (4.6). This implies $`\delta _1(P)=d_2(t_2x_2)`$ and $`\delta _2=d_1(t_1x_1)`$, by the definition of $`\delta _i`$ on $`E(K)`$. Therefore,
$$d_1(x_1^{}t_1)=m_1^{}(xt)\delta _1(P)=d_2(t_2x_2)=\mathrm{\Delta }^xd_1(t_1x_1).$$
(4.11)
Since $`\delta _1(s_1x)=r_1(x_1)\mathrm{\Delta }^x`$ (by (4)) and $`r_1(x_1)=r_1(t_1)+d_1(t_1x_1)`$ (by Statement 4.1),
$$d_1(s_1x_1^{})=m_1^{}(s_1x)\delta _1(s_1x)=r_1(x_1)\mathrm{\Delta }^x=r_1(t_1)+d_1(t_1x_1)\mathrm{\Delta }^x.$$
(4.12)
Comparing (4.11) and (4.12), we obtain $`d_1(s_1x_1^{})+d_1(x_1^{}t_1)r_1(t_1)`$, whence $`x_1^{}J_1(t_1)`$. Similarly, $`x_2^{}J_2(t_2)`$. So $`x^{}J(t)`$, yielding $`x^{}T`$, by (4.4).
Finally, consider an edge $`e=xyE(K)`$, and let $`x^{}=\gamma (x)`$ and $`y^{}=\gamma (y)`$. We have $`\delta _1(e)+\delta _2(e)=1`$, by (4). Also $`m_i^{}(e)=\delta _i(e)`$, $`i=1,2`$, by (4.10). Hence,
$$d(x^{}y^{})=d_1(x_1^{}y_1^{})+d_2(x_2^{}y_2^{})=m_1^{}(e)+m_2^{}(e)=\delta _1(e)+\delta _2(e)=1,$$
i.e., $`x^{}y^{}`$ is an edge of $`K`$, as required.
It remains to prove Lemma 4.2.
## 5 Proof of Lemma 4.2
We may prove this lemma for $`i=1`$. First we explain that $`\delta _1`$ is cyclically even, i.e., the $`\delta _1`$-length of any cycle in $`G_1`$ is even.
For any 4-circuit $`C=x^0x^1x^2x^3x^0`$ in $`K`$, an edge of $`C`$ belongs to $`E_1`$ if and only if the opposite edge does. Also, letting $`\eta _j:=\mathrm{\Delta }^{x^{j+1}}\mathrm{\Delta }^{x^j}`$, the numbers $`\eta _0,\eta _2`$ have the same parity if and only if $`\eta _1,\eta _3`$ do so. From these properties and the definition of $`\delta _i`$ one can deduce that the $`\delta _1`$-length of $`C`$ is even. Then $`\delta _1`$ is cyclically even within $`K`$, because $`K`$ is modular and, therefore, the 4-circuits form a basis in the space of cycles of $`K`$ over $`\text{Z}_2`$. (Indeed, any cycle of length $`q6`$ in a modular graph can be represented as the modulo two sum of three cycles with length less than $`q`$ each.) Next, using the fact that $`\delta _1`$ takes value one on $`U_1(E_1U)`$ and zero on $`(E_2U)𝒜_1`$, one can see that the $`\delta _1`$-length of any cycle with all edges in $`U_1U𝒜_1`$ is even. Finally, for an edge $`e=xs_1_1`$, choose $`tN(x)`$ and a shortest $`t`$$`x`$ path $`L`$ in $`K`$. Then $`\delta _1(L)=d_2(t_2x_2)`$. Concatenating $`L`$ with the edge $`e`$, the edge $`t_1t`$ in $`𝒜_1`$ and a shortest $`s_1`$$`t_1`$ path $`R`$ in $`H_1`$, we obtain a cycle whose $`\delta _1`$-length is equal to
$$\delta _1(L)+\delta _1(e)+\delta _1(R)+\delta _1(t_1t)=d_2(t_2x_2)+(r_1(x_1)\mathrm{\Delta }^x)+r_1(t_1)+0=2r_1(t_1).$$
Summing up the above observations, one can conclude that $`\delta _1`$ is cyclically even within the entire set $`_1`$. Then $`m_1`$ is cyclically even as well.
Next we prove that $`m_1`$ is an extension of $`d_1`$. The main part of this proof is to show the following property:
1. for any path $`P=x^0x^1\mathrm{}x^k`$ in $`K`$ with $`x^0T`$, there exists a path $`L=z^0z^1\mathrm{}z^\alpha `$ with $`z^0=x^0`$ and $`z^\alpha =x^k`$ and a number $`0\beta \alpha `$ such that $`z^0,\mathrm{},z^\beta T`$, that $`r(z^\beta )<r(z^{\beta +1})<\mathrm{}<r(z^\alpha )`$, and that $`\delta _1(L)\delta _1(P)`$.
The proof of ((5.1)) includes Claims 1–3 below. Recall that any edge $`xyE(K)`$ satisfies $`|r(x)r(y)|=1`$ (since $`K`$ is bipartite), and if $`xT`$ and $`r(x)>r(y)`$, then $`yT`$ (by (4.4)). In particular, $`L`$ as in ((5.1)) entirely lies in $`H`$ if $`x^kT`$. To show ((5.1)), it suffices to consider the case when $`P`$ is simple, $`k2`$, and all intermediate nodes of $`P`$ are not in $`T`$ (for if $`x^iT`$ for some $`0<i<k`$, we can split $`P`$ into two paths $`P^{}=x^0\mathrm{}x^i`$ and $`P^{\prime \prime }=x^i\mathrm{}x^k`$ and prove ((5.1)) for each of $`P^{},P^{\prime \prime }`$ independently). For $`i=0,\mathrm{},k`$, let $`r(i):=r(x^i)`$. An intermediate node $`x^i`$ of $`P`$ is called a peak if $`r(i)>r(i1)=r(i+1)`$. The set of peaks is denoted by $`F=F(P)`$. We prove ((5.1)) by induction on
$$\omega (P)=(4^{r(i)}:x^iF(P)).$$
If $`F=\mathrm{}`$, then $`r(0)<r(1)<\mathrm{}<r(k)`$ (as $`r(0)>r(1)`$ would imply $`x^1T`$), i.e., $`P`$ is just the desired path $`L`$. So assume $`F\mathrm{}`$. Let $`x^p`$ be the first peak in $`P`$, and let $`x,y,z`$ stand for $`x^{p1},x^p,x^{p+1}`$, respectively. Choose a median $`y^{}`$ for $`x,z,s`$ in $`K`$. Since $`r(x)=r(z)`$ and $`d(xz)=2`$, both $`xy^{},y^{}z`$ are edges of $`K`$ and $`r(y^{})<r(x)<r(y)`$. Replace $`y`$ by $`y^{}`$ in $`P`$, forming the path $`P^{}=x^0\mathrm{}x^{p1}y^{}x^{p+1}\mathrm{}x^k`$; we say that $`P^{}`$ is obtained by cutting off the peak $`y`$. Since $`4^{r(p)}>24^{r(p)1}=4^{r(p1)}+4^{r(p+1)}`$, we have $`\omega (P^{})<\omega (P)`$. Also $`\delta _1(P)\delta _1(P^{})`$ is equal to
$$\rho :=\rho (x,y,z,y^{}):=\delta _1(xy)+\delta _1(yz)\delta _1(xy^{})\delta _1(y^{}z).$$
Therefore, if $`\rho 0`$ occurs, we can immediately apply induction. Let $`\overline{\mathrm{\Delta }}:=\mathrm{\Delta }^y`$.
Claim 1 A median $`y^{}`$ for $`x,z,s`$ can be chosen so that $`\rho (x,y,z,y^{})<0`$ is possible only if both edges $`e=xy,e^{}=yz`$ are in $`E_2`$, $`\mathrm{\Delta }^x=\mathrm{\Delta }^z=\overline{\mathrm{\Delta }}`$, and $`\mathrm{\Delta }^y^{}=\overline{\mathrm{\Delta }}1`$.
Proof. Since the $`\delta _1`$-length of the 4-circuit $`C=xyzy^{}x`$ is even, $`\rho <0`$ implies
$$\delta _1(e)=\delta _1(e^{})=0\text{and}\delta _1(xy^{})=\delta _1(y^{}z)=1.$$
(5.2)
This is impossible when $`eE_1`$ and $`e^{}E_2`$ (or $`eE_2`$ and $`e^{}E_1`$). Indeed, in this case we would have $`\mathrm{\Delta }^x=\overline{\mathrm{\Delta }}1`$ and $`\mathrm{\Delta }^z=\overline{\mathrm{\Delta }}`$, by (4). Then $`xy^{}E_2`$ and $`\delta _1(xy^{})=1`$ imply $`\mathrm{\Delta }^y^{}=\mathrm{\Delta }^x1=\overline{\mathrm{\Delta }}2`$, while $`y^{}zE_1`$ and $`\delta _1(y^{}z)=1`$ imply $`\mathrm{\Delta }^y^{}=\mathrm{\Delta }^z=\overline{\mathrm{\Delta }}`$; a contradiction.
If $`e,e^{}E_2`$, then $`xy^{},y^{}zE_2`$. So (5.2) yields $`\mathrm{\Delta }^x=\mathrm{\Delta }^z=\overline{\mathrm{\Delta }}`$ and $`\mathrm{\Delta }^y^{}=\mathrm{\Delta }^x1=\overline{\mathrm{\Delta }}1`$, as required.
Now, suppose $`e,e^{}E_1`$ and $`\delta _1(e)=\delta _1(e^{})=0`$. Choose $`uN(x)`$ and $`vN(z)`$. We have $`\mathrm{\Delta }^x=\mathrm{\Delta }^z=\overline{\mathrm{\Delta }}1`$, whence $`u,vN(y)`$. Choose in $`T`$ a median $`q`$ for $`u,v,(y_1,s_2)`$ and a median $`w`$ for $`u,v,(s_1,y_2)`$. We assert that $`q,wN(y)`$. Indeed,
$$q_1M_1(u_1,v_1,y_1),w_1M_1(u_1,v_1,s_1),q_2M_1(u_2,v_2,s_2),w_2M_1(u_2,v_2,y_2).$$
In particular, $`q_1,w_1I_1(u_1,v_1)`$. Also $`u_1,v_1I_1(q_1,w_1)`$ (in view of $`u_1,v_1I_1(y_1,s_1)`$, by Statement 4.1). These relations imply $`d_1(u_1q_1)=d_1(v_1w_1):=a`$. Similarly, $`d_2(u_2q_2)=d_2(v_2w_2):=a^{}`$. Then $`d(yu)=\overline{\mathrm{\Delta }}d(yq)=d(yu)a+a^{}`$ and $`d(yv)d(yw)=d(yv)+aa^{}`$. This is possible only if $`a=a^{}`$, yielding $`d(yq)=d(yw)=\overline{\mathrm{\Delta }}`$, as required.
Assume $`y^{}`$ is chosen to be a median for $`x,z,w`$. Then $`y^{}`$ is a median for $`x,z,s`$ as well, taking into account that $`x_2=z_2`$ and the paths $`(x_1,u_1,w_1,s_1)`$ and $`(z_1,v_1,w_1,s_1)`$ on $`T_1`$ are $`d_1`$-shortest. Now $`d(y^{}w)=d(xw)1`$ implies $`\mathrm{\Delta }^y^{}<\mathrm{\Delta }^x`$. Hence, $`\delta _1(xy^{})=\delta _1(y^{}z)=0`$ and $`\rho =0`$.
Arguing as in the above proof, one can see that for any $`x^{}V(K)`$, there are elements $`t,t^{}N(x^{})`$ such that $`r_1(t_1)r_1(t_1^{})`$ (and $`r_2(t_2)r_2(t^{})`$) and $`N(x^{})I(t,t^{})`$. We denote $`t`$ by $`t(x^{})`$ and refer to it as the minimal element of $`N(x^{})`$ (with respect to the rank in $`H_1`$).
Remark 4. For $`i=1,2`$ and $`f,gN(x^{})`$, denote $`f_i_ig_i`$ if $`f_iJ_i(g_i)`$. Then $`_i`$ is the partial order on $`N_i=\{w_i:wN(x^{})\}`$ with unique minimal and maximal elements. Moreover, the correspondence $`w_1w_2`$ establishes the isomorphism between $`(N_1,_1)`$ and $`(N_2,_2^1)`$ (where $`^1`$ is the reverse to $``$). One can show that if none of $`H_1,H_2`$ containes $`K_{3,3}^{}`$ as an induced subgraph (see Fig. 1b), then $`(N_i,_i)`$ is a modular lattice, i.e., (i) any $`u,vN_i`$ have unique lower and upper bounds, denoted by $`uv`$ and $`uv`$, respectively; (ii) for each $`uN_i`$, all maximal chains to $`u`$ from the minimal element have the same length $`\rho (u)`$, and (iii) each pair $`u,v`$ satisfies the modular equality $`\rho (u)+\rho (v)=\rho (uv)+\rho (uv)`$. We, however, do not need these properties in further arguments.
In light of Claim 1, we may assume that $`\rho <0`$ and $`e,e^{}E_2`$. Consider the minimal element $`t(y)=(t_1(y),t_2(y))`$ in $`N(y)`$. Suppose $`t_1(y)y_1`$. Then there is a node $`w`$ of $`K`$ adjacent to $`y`$ such that $`w_1I_1(y_1,t_1(y))`$ and $`w_2=y_2`$. We have $`ywE_1`$, $`r(w)=r(y)1`$ and $`t(y)N(w)`$. Then $`\mathrm{\Delta }^w<\overline{\mathrm{\Delta }}`$ and $`\delta _1(yw)=0`$. Transform $`P`$ into the (non-simple) path $`P^{}=x^0\mathrm{}x^{p2}xywyzx^{p+2}\mathrm{}x^k`$ and then cut off both copies of $`y`$ (which are peaks of $`P^{}`$). This results in a path $`P^{\prime \prime }`$ of the form $`x^0\mathrm{}x^{p2}xy^{}wy^{\prime \prime }zx^{p+2}\mathrm{}x^k`$; clearly $`x,w,z`$ are peaks of $`P^{\prime \prime }`$. Since $`ywE_1`$, $`y^{}`$ and $`y^{\prime \prime }`$ can be chosen so that $`\rho (x,y,w,y^{})0`$ and $`\rho (w,y,z,y^{\prime \prime })0`$, by Claim 1. Therefore, $`\delta _1(P^{\prime \prime })\delta _1(P^{})=\delta _1(P)`$. Also $`4^{r(y)}>34^{r(y)1}=4^{r(x)}+4^{r(w)}+4^{r(z)}`$, yielding $`\omega (P^{\prime \prime })<\omega (P)`$. So we can apply induction.
It remains to consider the case when $`t_1(y)=y_1`$. Then $`t(y)`$ is the unique element of $`N(y)`$. We will use the following property.
Claim 2. Let $`\overline{x}\overline{y}E_2`$ satisfy $`r(\overline{x})<r(\overline{y})`$, let $`N(\overline{y})`$ consist of a single element $`u`$, and let $`u_1=\overline{y}_1`$. Then $`N(\overline{x})`$ consists of a single element $`v`$, and $`v_1=\overline{y}_1`$. Moreover, $`u=v`$ if $`\mathrm{\Delta }^{\overline{x}}<\mathrm{\Delta }^{\overline{y}}`$, and $`u`$ and $`v`$ are adjacent if $`\mathrm{\Delta }^{\overline{x}}=\mathrm{\Delta }^{\overline{y}}`$.
Proof. If $`\mathrm{\Delta }^{\overline{x}}<\mathrm{\Delta }^{\overline{y}}`$, then $`N(\overline{x})N(\overline{y})`$, whence $`N(\overline{x})=\{u\}`$. So assume $`\mathrm{\Delta }^{\overline{x}}=\mathrm{\Delta }^{\overline{y}}`$, and let $`vN(\overline{x})`$. Choose $`qM(u,v,(\overline{y}_1,s_2))T`$ and $`wM(u,v,(s_1,\overline{y}_2))T`$. We have $`q_2,w_2I_2(u_2,v_2)`$ and $`u_2I_2(q_2,w_2)`$ (in view of $`u_2I_2(\overline{y}_2,s_2)`$). Note that the path $`(\overline{y}_2,\overline{x}_2,v_2,s_2)`$ on $`T_2`$ is $`d_2`$-shortest (since $`r(\overline{x})<r(\overline{y})`$ and $`\overline{x}_1=\overline{y}_1`$ imply $`\overline{x}_2I_2(\overline{y}_2,s_2)`$). This yields $`v_2I_2(q_2,w_2)`$, and we can conclude that $`d_2(u_2w_2)=d_2(v_2q_2)=:a^{}`$.
Next, $`q_1M_1(u_1,v_1,\overline{y}_1)`$ and $`u_1=\overline{y}_1`$ imply $`q_1=\overline{y}_1`$, while $`w_1M_1(u_1,v_1,s_1)`$, $`v_1I_1(\overline{x}_1,s_1)`$ and $`\overline{x}_1=\overline{y}_1=\overline{u}_1`$ imply $`w_1=v_1`$. Let $`a:=d_1(\overline{y}_1v_1)`$. Then $`d(\overline{x}v)d(\overline{x}q)=d(\overline{x}v)a+a^{}`$ and $`d(\overline{y}u)d(\overline{y}w)=d(\overline{y}u)+aa^{}`$, whence $`a=a^{}`$, $`qN(\overline{x})`$ and $`wN(\overline{y})`$. Since $`|N(\overline{y})|=1`$, we have $`w=u`$. This implies $`a=0`$ and $`q=v`$, yielding $`v_1=q_1=\overline{y}_1`$. So $`v_1=\overline{y}_1`$, regardless of the choice of $`v`$ in $`N(\overline{x})`$. This is possible only if $`N(\overline{x})`$ consists of a single element (for if $`v,v^{}N(\overline{x})`$ and $`vv^{}`$, then a median $`f`$ for $`v,v^{},(s_1,\overline{x}_2)`$ in $`T`$ satisfies $`f_1=\overline{y}_1`$ and $`d_2(\overline{x}_2f_2)<d_2(\overline{x}_2v_2)`$, whence $`d(\overline{x}f)<\mathrm{\Delta }^{\overline{x}}`$).
Finally, to see that $`u_2`$ and $`v_2`$ are adjacent, take in $`T`$ a median $`h`$ for $`u,v,(s_1,\overline{x}_2)`$. Then $`d(\overline{x}h)d(\overline{x}v)`$, $`h_2I_2(\overline{x}_2,v_2)`$ and $`h_1=\overline{y}_1`$, implying $`h=v`$. So $`v_2I_2(\overline{x}_2,u_2)`$. Also $`d_2(\overline{y}_2u_2)=\mathrm{\Delta }^{\overline{y}}=\mathrm{\Delta }^{\overline{x}}=d_2(\overline{x}_2v_2)`$ and $`u_2I_2(\overline{y}_2,v_2)`$ (since $`u_2I_2(\overline{y}_2,s_2)`$ and $`v_2=q_2I_2(u_2,s_2)`$). Now $`d_2(\overline{x}_2\overline{y}_2)=1`$ implies $`d_2(u_2v_2)=1`$, as required.
For $`i=0,\mathrm{},p`$, define $`P_i`$ to be the subpath $`x^i\mathrm{}x^p`$ of $`P`$. Let $`P_j`$ be the maximal subpath with all edges in $`E_2`$ (i.e., $`j`$ is minimum subject to $`x_1^j=\mathrm{}=x_1^p`$). Since $`r(j)<r(j+1)<\mathrm{}<r(p)`$, we can repeatedly apply Claim 1 to the edges of $`P_j`$, starting with $`x^{p1}x^p`$, and conclude that $`N(x^i)`$ is a singleton $`\{u^i\}`$ with $`u_1^i=y_1`$ for each $`i=j,\mathrm{},p`$. Also $`u^i=u^{i+1}`$ if $`\mathrm{\Delta }_i<\mathrm{\Delta }_{i+1}`$, and $`u^iu^{i+1}E_2`$ if $`\mathrm{\Delta }_i=\mathrm{\Delta }_{i+1}`$, where $`\mathrm{\Delta }_q`$ stands for $`\mathrm{\Delta }^{x^q}`$. Consider two possible cases.
Case 1: $`j1`$. By the maximality of $`P_j`$, $`x^{j1}x^jE_1`$. Let $`b:=x_1^{j1}`$. For $`i=j,\mathrm{},p`$, define $`z^i`$ and $`v^i`$ to be the points with $`z_1^i=v_1^i=b`$, $`z_2^i=x_2^i`$ and $`v_2^i=u_2^i`$, i.e., $`z^i`$ and $`v^i`$ are obtained by shifting the points $`x^i`$ and $`u^i`$, respectively, along the edge $`y_1b`$ of $`H_1`$. In particular, $`z^j=x^{j1}`$. Denote $`\mathrm{\Delta }^{z^i}`$ by $`\mathrm{\Delta }_i^{}`$.
Claim 3. $`\mathrm{\Delta }_i^{}=\mathrm{\Delta }_i`$ and $`v^iN(z^i)`$ for each $`i=j,\mathrm{},p`$.
Proof. Since $`r(j1)<r(j)`$ and $`x_2^{j1}=x_2^j`$, $`r_1(b)<r_1(x^j)`$. Therefore, $`u^iT`$ implies $`v^iT`$, and we have $`\mathrm{\Delta }_i^{}d_2(z^iv^i)=d_2(x^iu^i)=\mathrm{\Delta }_i`$. Suppose $`\mathrm{\Delta }_i^{}<\mathrm{\Delta }_i`$. Then $`N(z^i)N(x^i)`$, whence $`N(z^i)=\{u^i\}`$. But $`d(z^iu^i)=d_1(by_1)+d_2(x_2^iu_2^i)=1+d(z^iv^i)`$; a contradiction. Thus, $`\mathrm{\Delta }_i^{}=\mathrm{\Delta }_i`$ and $`v^iN(z^i)`$.
Consider the $`x^{j1}`$$`y`$ paths $`P_{j1}`$ and $`R=z^j\mathrm{}z^px^p`$ in $`K`$. From Claim 3 it follows that $`\delta _1(z^iz^{i+1})=\delta _1(x^ix^{i+1})`$ for $`i=j,\mathrm{},p1`$, and that $`\delta _1(x^{j1}x^j)=\delta _1(z^px^p)`$. Therefore, $`\delta _1(P_{j1})=\delta _1(R)`$. Replace in $`P`$ the part $`P_{j1}`$ by $`R`$, forming the path $`P^{}=x^0\mathrm{}x^{j1}z^j\mathrm{}z^px^p\mathrm{}x^k`$. Clearly $`y=x^p`$ is the first peak of $`P^{}`$. Cut off $`y`$ in $`P^{}`$ by replacing $`y`$ by a median $`y^{\prime \prime }`$ for $`z^p,z,s`$; let $`P^{\prime \prime }`$ be the resulting path. Since $`z^pyE_1`$ and $`yzE_2`$, one has $`\rho (z^p,y,z,y^{\prime \prime })0`$, by Claim 1. Therefore, $`\delta _1(P^{\prime \prime })\delta _1(P^{})=\delta _1(P)`$, and ((5.1)) follows by induction because $`z^p`$ and $`z`$ are the first and second peaks of $`P^{\prime \prime }`$ and $`4^{r(y)}>4^{r(z^p)}+4^{r(z)}`$.
Case 2: $`j=0`$. Then $`x^0=u^0`$. By Claim 2 applied to the edge $`zy`$, $`N(z)`$ is a singleton $`\{\widehat{u}\}`$ with $`\widehat{u}_1=y_1`$. As before, let $`y^{}M(x,z,s)T`$; then $`y_1^{}=y_1`$ and $`N(y^{})`$ is a singleton $`\{v\}`$ (by Claim 2 applied to the edge $`y^{}x`$). Assuming $`\mathrm{\Delta }^y^{}<\mathrm{\Delta }^x`$ (equivalently: $`\rho <0`$), we have $`N(y^{})N(x)N(z)`$. Hence, $`v=\widehat{u}=u^{p1}`$.
Form the $`u^0`$$`v`$ path $`R^{}`$ by deleting repeated consecutive elements in $`u^0\mathrm{}u^{p1}`$, and let $`\overline{R}`$ be the concatenation of $`R^{}`$, a shortest $`v`$$`y^{}`$ path $`R^{\prime \prime }`$, and the edge $`y^{}x`$. Clearly the $`\delta _1`$-length of each edge of $`R^{}`$ is zero, while the $`\delta _1`$-length of each edge of $`R^{\prime \prime }`$ is one. Also $`\delta _1(y^{}x)=1`$.
Comparing $`\overline{R}`$ with the path $`\overline{P}=x^0\mathrm{}x^{p1}`$ and using Claim 2, one can deduce that $`\overline{R}=p1`$ (i.e., $`\overline{R}`$ is a shortest path in $`K`$) and that $`\delta _1(\overline{R})=\delta _1(\overline{P})`$. Now let $`D`$ be the concatenation of $`R^{}`$, $`R^{\prime \prime }`$ and the edge $`y^{}z`$. Since $`\delta _1(y^{}x)=\delta _1(y^{}z)`$ and $`\delta _1(xy)=\delta _1(yz)=0`$, we have $`\delta _1(D)=\delta _1(\overline{R})=\delta _1(P_0)`$. Also $`|D|=|R|=p1`$ implies that $`D`$ has no peaks. Then, replacing in $`P`$ the part $`x^0\mathrm{}x^{p+1}`$ by $`D`$, we obtain the path $`P^{}`$ with $`\delta _1(P^{})=\delta (P)`$ and $`\omega (P^{})<\omega (P)`$ and can apply induction.
Thus, ((5.1)) is proven. In order to conclude that $`m_1`$ is an extension of $`d_1`$, it suffices to consider a path $`L`$ as in ((5.1)) and show the following:
1. (i) if $`z^\alpha T`$, then $`\delta _1(L)d_1(z_1^0z_1^\alpha )`$;
* $`\delta _1(L)+\delta _1(z^\alpha s_1)r_1(z_1^0)`$.
(In fact, (i) embraces the case of a path in $`G_1`$, with both ends in $`T_1`$, whose first and last edges belong to $`𝒜_1`$, while (ii) does the case when one of these edges is in $`𝒜_1`$ and the other in $`_1`$.) Case (i) is trivial because $`z^\alpha T`$ means that $`L`$ is a path in $`H`$, and therefore, the $`\delta _1`$-length of each of its edges in $`E_1`$ is equal to one. So let us prove (ii). One may assume that $`r(z^0)<\mathrm{}<r(z^\alpha )`$ (taking into account that $`d_1(z_1^0s_1)d_1(z_1^0z_1^\beta )+d_1(z_1^\beta z_1^\alpha )`$ and $`\delta _1(L)=\delta _1(L^{})+\delta _1(L^{\prime \prime })`$, where $`L^{}=z^0\mathrm{}z^\beta `$ and $`L^{\prime \prime }=z^\beta \mathrm{}z^\alpha `$, and assuming w.l.o.g. that $`L^{}`$ is $`\delta _1`$-shortest).
For $`i=0,\mathrm{},\alpha `$, let $`\mathrm{}_i`$ denote the $`\delta _1`$-length of the path $`z^0\mathrm{}z^i`$, and let $`\rho _i`$ and $`\mathrm{\Delta }_i`$ stand for $`r_1(z_1^i)`$ and $`\mathrm{\Delta }^{z^i}`$, respectively. By the definition of $`\delta _1`$ on $`_1`$, $`\delta _1(z^is_1)`$ is equal to $`\rho _i\mathrm{\Delta }_i`$. We show that
$$\mathrm{}_i+\rho \mathrm{\Delta }_i\rho _0,$$
(5.4)
using induction on $`i`$. This gives the desired inequality ((5.3))(ii) when $`i=\alpha `$. Since $`\mathrm{}_0=\mathrm{\Delta }_0=0`$, (5.4) holds for $`i=0`$. Assume it holds for $`i1`$ ($`0<i<\alpha `$), and let $`a:=\mathrm{}_i\mathrm{}_{i1}`$, $`b:=\rho _i\rho _{i1}`$ and $`c:=\mathrm{\Delta }_i\mathrm{\Delta }_{i1}`$. Then (5.4) for $`i`$ follows from $`a+bc0`$. To see the latter, consider four possible cases for $`e=z^{i1}z^i`$, taking into account that $`\mathrm{\Delta }_i\mathrm{\Delta }_{i1}`$ since $`r(z^i)>r(z^{i1})`$.
(a) Let $`eE_1`$ and $`\mathrm{\Delta }_i=\mathrm{\Delta }_{i1}`$. Then $`a+bc=1+1+0=2`$.
(b) Let $`eE_1`$ and $`\mathrm{\Delta }_i>\mathrm{\Delta }_{i1}`$. Then $`a+bc=0+11=0`$.
(c) Let $`eE_2`$ and $`\mathrm{\Delta }_i=\mathrm{\Delta }_{i1}`$. Then $`a+bc=0+00=0`$.
(d) Let $`eE_2`$ and $`\mathrm{\Delta }_i>\mathrm{\Delta }_{i1}`$. Then $`a+bc=1+01=0`$.
Thus, $`m_1`$ is an extension of $`d_1`$. It remains to show that $`m_i(e)=\delta _i(e)`$ for $`i=1,2`$ and $`e_i`$. This is obvious when $`eUU_i`$ or when $`\delta _i(e)=0`$. If $`e=xs_i_i`$, then $`m_i(e)=\delta _i(e)`$ follows from the fact that for $`tN(x)`$, the path in $`G_i`$ obtained by concatenationg the edge $`t_it`$, a shortest $`t`$$`x`$ path in $`K`$, and the edge $`xs_i`$ is $`\delta _i`$-shortest (this fact was shown at the beginning of this section). Finally, each edge $`eE(K)`$ belongs to a shortest $`t`$$`t^{}`$ path $`P`$ in $`K`$ with $`t,t^{}T`$. Since $`\delta _1(e^{})+\delta _2(e^{})=1`$ for all edges $`e^{}`$ of $`K`$, we have $`\delta _1(P)+\delta _2(P)=|P|=d(tt^{})=d_1(t_1t_1^{})+d_2(t_2t_2^{})`$, whence $`\delta _i(P)=d_i(t_it_i^{})`$, implying $`m_i(e)=\delta _i(e)`$.
This completes the proof of Lemma 4.2 and completes the proof of Theorems 1.4 and 1.3.
## 6 Intractable Cases
In this section we prove Theorem 1.6, considering a metric $`\mu `$ on a set $`T`$ such that either $`\mu `$ is non-modular or $`\mu `$ is modular but its underlying graph $`H=(T,U)`$ is non-orientable. W.l.o.g., one may assume $`\mu `$ is integer-valued. Our method borrows the idea from for the path metrics $`\mu =d^H`$ as in Theorem 1.5, which in turn generalizes the construction from for $`H=K_3`$.
Given a set $`VT`$, a function $`E_V\text{Z}_+`$, nodes $`s,tT`$, and points $`x,yVT`$, let $`\tau (s,x|t,y)`$ denote the minimum $`cm`$ among all $`m\mathrm{Ext}^0(\mu ,V)`$ such that $`m(xs)=m(yt)=0`$.
The core of the proof in that the 3-terminal cut problem is NP-hard is the construction of a “gadget” $`(V,c)`$ with specified $`s,t,x,y`$ satisfying the following property:
1. (i) $`\tau (s,x|t,y)=\tau (s,y|t,x)=\widehat{\tau }`$,
* $`\tau (s,x|s,y)=\tau (t,x|t,y)=\widehat{\tau }+\delta `$ for some $`\delta >0`$,
* $`\tau (s^{},x|t^{},y)\widehat{\tau }+\delta `$ for all other pairs $`\{s^{},t^{}\}`$ in $`T`$,
where $`\widehat{\tau }`$ stands for $`\tau (V,c,\mu )`$ (with $`\mu =d^{K_3}`$). Then the NP-hardness of the problem is easily shown by a reduction from MAX CUT.
Our aim is to construct corresponding “gadgets” satisfying ((6.1)) for $`\mu `$ as in Theorem 1.6; then the theorem will follow by a similar reduction.
First we consider the case when $`\mu `$ is modular but $`H`$ is non-orientable, which is technically simpler. In fact, the construction and arguments in this case are similar to those for the corresponding unweighted case ($`\mu =d^H`$) given in \[11, Sec. 6\]. More precisely, since $`H`$ is non-orientable, there exists a projective sequence $`(e_0,e_1,\mathrm{},e_{k1},e_k=e_0)`$ of edges of $`H`$ yielding the “twist” (or forming the orientation-reversing dual cycle). That is,
1. for $`i=0,\mathrm{},k1`$, $`e_i=s_it_i`$ and $`e_{i+1}=s_{i+1}t_{i+1}`$ are opposite edges in the 4-circuit $`C_i=s_it_it_{i+1}s_{i+1}s_i`$, and $`t_k=s_0`$ (and $`s_k=t_0`$).
(One can choose such a sequence with all edges (though not necessarily the nodes) distinct, but this is not important for us.) Since $`\mu `$ is modular, we have by ((2.2)) that
1. for $`i=0,\mathrm{},k1`$, $`\mu (e_i)`$ is a constant $`h`$, and $`\mu (s_is_{i+1})=\mu (t_it_{i+1})=:f_i`$.
We denote $`t_i`$ by $`s_{i+k}`$ and take indices modulo $`2k`$. The desired gadget is represented by the graph $`G=(V,E)`$ with the weights $`c(e)`$ of edges $`eE`$, where $`V=T\{z_0,\mathrm{},z_{2k1}\}`$ and for $`i=0,\mathrm{},2k1`$,
* $`z_i`$ is adjacent to both $`s_i`$ and $`s_{i+k}`$, and $`c(z_is_i)=c(z_is_{i+k})=N`$ for a positive integer $`N`$ (specified below);
* $`z_i`$ and $`z_{i+1}`$ are adjacent, and $`c(z_iz_{i+1})=1`$.
Figure 4 illustrates $`G`$ for $`k=4`$. We put $`s=s_0`$, $`t=t_0`$, $`x=z_0`$ and $`y=z_k`$, and formally extend $`c`$ by zero to $`E_VE`$. We assert that ((6.1)) holds.
Indeed, each $`m\mathrm{Ext}^0(\mu ,V)`$ is associated with the mapping $`\gamma :\{z_0,\mathrm{},z_{2k1}\}T`$, where $`\gamma (z_i)=s_j`$ if $`m(z_is_j)=0`$; we say that $`z_i`$ is attached by $`\gamma `$ to $`s_j`$ and denote $`m`$ by $`m^\gamma `$. If $`\gamma (z_i)=v`$, then, letting $`ϵ:=\mu (s_iv)+\mu (vs_{i+k})\mu (s_is_{i+k})`$, the contribution to the volume $`cm^\gamma `$ due to the edges $`e=z_is_i`$ and $`e^{}=z_is_{i+k}`$ is equal to
$$c(e)m^\gamma (e)+c(e^{})m^\gamma (e^{})=N(m^\gamma (e)+m^\gamma (e^{}))=Nh+Nϵ;$$
cf. ((6.3)). We have $`ϵ=0`$ if $`v\{s_i,s_{i+k}\}`$, and $`ϵ1`$ otherwise. Hence, every mapping $`\gamma `$ pretending to be optimal or nearly optimal must attach each $`z_i`$ to either $`s_i`$ or $`s_{i+k}`$ whenever $`N`$ is chosen sufficiently large (e.g., $`N=1+2k\mathrm{max}\{\mu (st):s,tT\}`$).
Next, if $`z_i`$ is attached to $`s_i`$ (resp. $`s_{i+k}`$) and $`z_{i+1}`$ to $`s_{i+1}`$ (resp. $`s_{i+1+k}`$), then the edge $`u=z_iz_{i+1}`$ contributes $`c(u)m^\gamma (u)=f_i`$ (cf. ((6.3)), letting $`f_j=f_{j+k}`$. On the other hand, if $`z_i`$ is attached to $`s_i`$ (resp. $`s_{i+k}`$) while $`z_{i+1}`$ to $`s_{i+1+k}`$ (resp. $`s_{i+1}`$), then the contribution becomes $`h+f_i`$ ($`=\mu (s_it_{i+1})`$).
So we can conclude that $`\widehat{\tau }=2khN+2(f_1+\mathrm{}+f_k)`$, and there are precisely two optimal 0-extensions, namely, $`m^{\gamma _1}`$ and $`m^{\gamma _2}`$, where $`\gamma _1(z_i)=s_i`$ and $`\gamma _2(z_i)=s_{i+k}`$ for $`i=0,\mathrm{},2k1`$. This gives (i) in ((6.1)). Furthermore, one can see that if $`m^\gamma `$ is the least-volume 0-extension induced by $`\gamma `$ that brings both $`x,y`$ either to $`s`$ or to $`t`$, then $`m^\gamma (z_jz_{j+1})=h+f_j`$ for precisely two numbers $`j\{0,\mathrm{},2k1\}`$ such that $`f_j=\mathrm{min}\{f_1,\mathrm{},f_k\}`$. So $`cm^\gamma =\widehat{\tau }+2h`$, yielding ((6.1))(ii). Finally, (iii) is ensured by the choice of $`N`$.
Thus, ((1.1)) with $`\mu `$ modular and $`H`$ non-orientable is NP-hard. Moreover, it is strongly NP-hard because the number $`N`$ is a constant depending only on $`\mu `$.
Next we consider the case when $`\mu `$ is not modular. Let $`\mathrm{\Delta }(x,y,z)`$ denote the value (perimeter) $`\mu (xy)+\mu (yz)+\mu (zx)`$ for $`x,y,zT`$. We fix a medianless triplet $`\{s_0,s_1,s_2\}`$ such that $`\mathrm{\Delta }(s_0,s_1,s_2):=\overline{\mathrm{\Delta }}`$ is minimum. By technical reasons, we put $`s_{i+3}=s_i`$, $`i=0,1,2`$, and take indices modulo 6. The gadget $`(G=(V,E),c)`$ that we construct has a somewhat more complicated structure compared with that for the corresponding unweighted case in \[11, Sec. 6\]. Here
$$V=TZ,Z=\{z_0,\mathrm{},z_5\}\text{and}E=E_1E_2E_3.$$
For $`i=1,2,3`$, the edges $`eE_i`$ are endowed with weights $`c_i(e)`$, and $`c(e)`$ is defined to be $`N_ic_i(e)`$. The factors $`N_1,N_2,N_3`$ are chosen so that $`N_1=1`$, $`N_2`$ is sufficiently large, and $`N_3`$ is sufficiently large with respect to $`N_2`$. Informally speaking, the “heavy” edges of $`E_3`$ provide that (at optimality or almost optimality) each point $`z_j`$ gets into the interval $`I_j:=\{vT:\mu (s_{j1}v)+\mu (vs_{j+1})=\mu (s_{j1}s_{j+1})\}`$, then the “medium” edges of $`E_2`$ make $`z_j`$ choose only between the endpoints $`s_{j1},s_{j+1}`$ of $`I_j`$, and finally the “light” edges of $`E_1`$ provide the desired property ((6.1)).
As before, $`m^\gamma `$ denotes the 0-extension of $`\mu `$ to $`V`$ induced by $`\gamma :ZT`$. Define $`d_i:=d_{i+3}=\mu (s_{i1}s_{i+1})`$. We say that a path $`P=(v_1,\mathrm{},v_k)`$ on $`T`$ is shortest if it is $`\mu `$-shortest.
The set $`E_3`$ consists of the edges $`e_j=z_js_{j1}`$ and $`e_j^{}=z_js_{j+1}`$ with $`c_3(e_j)=c_3(e_j^{})=1`$ for $`j=0,\mathrm{},5`$. Then the contribution to $`cm^\gamma `$ due to $`e_j`$ and $`e_j^{}`$ is $`N_3d_j`$ if $`\gamma (z_j)I_j`$, and at least $`N_3d_j+N_3`$ otherwise, yielding that $`z_j`$ should be mapped into $`I_j`$, by the choice of $`N_3`$. The minimality of $`\overline{\mathrm{\Delta }}`$ provides the following useful property.
###### Statement 6.1
For any $`vI_j`$, at least one of the paths $`P=(s_j,s_{j1},v)`$ and $`P^{}=(s_j,s_{j+1},v)`$ is shortest.
Proof. Let for definiteness $`j=1`$. Suppose $`P^{}`$ is not shortest. Then $`\mu (s_1v)<|P^{}|=\mu (s_1s_2)+\mu (s_2v)`$ and $`\mu (s_0v)=\mu (s_0s_2)\mu (s_2v)`$ imply $`\mathrm{\Delta }(s_1,v,s_0)<\overline{\mathrm{\Delta }}`$. So $`s_1,v,s_0`$ have a median $`w`$. If $`w=s_0`$, $`P`$ is shortest. Otherwise we have $`\mathrm{\Delta }(s_1,w,s_2)<\mathrm{\Delta }`$ (since $`\mu (s_1w)<\mu (s_1s_0)`$ and the path $`(s_2,v,w,s_0)`$ is, obviously, shortest). Then $`s_1,w,s_2`$ have a median $`q`$. It is easy to see that $`q`$ is a median for $`s_0,s_1,s_2`$; a contradiction.
We now explain the construction of $`E_2`$ and $`c_2`$. Each $`z=z_j`$ ($`j=0,\mathrm{},5`$) is connected to each $`s_i`$ ($`i=0,1,2`$) by edge $`u_i=zs_i`$ whose weight is defined by
$$c_2(u_i)=(d_{i1}+d_{i+1}d_i)/(d_{i1}d_{i+1})=:a_i$$
(6.4)
($`a_i`$ is positive and does not depend on $`j`$). Suppose $`z`$ is mapped by $`\gamma `$ to some $`s_i`$, say $`\gamma (z)=s_1`$. Then, up to a factor of $`N_2`$, the contribution to $`cm^\gamma `$ from the edges $`u_0,u_1,u_2`$ (concerning $`z`$) is
$`d_2a_0+d_0a_2=d_2(d_1+d_2d_0)/(d_1d_2)+d_0(d_1+d_0d_2)/(d_0d_1)`$ (6.5)
$`(d_1+d_2d_0)/d_1+(d_1+d_0d_2)/d_1=2.`$
On the other hand, the contribution grows when $`z_j`$ falls into the interior of any interval $`I_i`$.
###### Statement 6.2
Let $`vI_i\{s_{i1},s_{i+1}\}`$. Then $`\sigma :=(a_i\mu (s_iv):i=0,1,2)>2`$.
Proof. Let for definiteness $`i=0`$, $`\mu (s_1v)=ϵ`$ and $`\mu (s_0v)=d_2+ϵ`$ (cf. Statement 6.1). Then
$$\sigma =(d_2+ϵ)a_0+ϵa_1+(d_0ϵ)a_2=d_2a_0+d_0a_2+ϵ(a_0+a_1a_2)=2+ϵ(a_0+a_1a_2),$$
in view of (6.5). We observe that $`a_0+a_1a_2>0`$. Indeed,
$$\begin{array}{c}d_0d_1d_2(a_0+a_1a_2)=(d_0d_1+d_0d_2d_0^2)+(d_1d_0+d_1d_2d_1^2)(d_2d_0+d_2d_1d_2^2)\hfill \\ =2d_0d_1d_0^2d_1^2+d_2^2=d_2^2(d_0d_1)^2>0\end{array}$$
since $`d_2>d_0d_1`$. So $`\sigma >2`$.
Thus, by an appropriate choice of constants $`N_2`$ and $`N_3`$, each point $`z_j`$ must be mapped to either $`s_{j1}`$ or $`s_{j+1}`$. Such a mapping $`\gamma `$ is called feasible. We now construct the crucial set $`E_1`$ and function $`c_1`$. The set $`E_1`$ consists of six edges $`g_j=z_jz_{j+1}`$, $`j=0,\mathrm{},5`$, forming the 6-circuit $`C`$ (this is similar to the construction in motivated by ). The essense is how to assign $`c_1`$. For $`i=0,1,2`$, let $`h_i:=h_{i+3}:=(d_{i1}+d_{i+1}d_i)/2`$. These numbers would be just the distances from $`s_0,s_1,s_2`$ to their median if it existed, i.e.,
$$d_i=h_{i1}+h_{i+1}.$$
(6.6)
We define
$$c_1(z_jz_{j+1})=c_1(z_{j+3}z_{j+4})=h_{j1}\text{for}j=0,1,2.$$
(6.7)
For $`\gamma :ZT`$, let $`\zeta ^\gamma `$ denotes $`(c_1(g_j)m^\gamma (g_j):j=0,\mathrm{},5)`$, i.e., $`\zeta ^\gamma `$ is the contribution to $`cm^\gamma `$ from the edges of $`C`$. The analysis below will depend on the numbers
$$\rho =2(h_0h_1+h_1h_2+h_2h_0)\text{and}\alpha =2\mathrm{min}\{h_0^2,h_1^2,h_2^2\}.$$
(6.8)
W.l.o.g., assume $`h_0h_1,h_2`$, i.e., $`2h_0^2=\alpha `$. Our aim is to show that ((6.1)) holds if we take as $`s,t,x,y`$ the elements $`s_0,s_2,z_1,z_4`$, respectively.
To show this, consider the mapping $`\gamma _1`$ as drawn in Fig. 5a, i.e., $`\gamma _1(z_j)`$ is $`s_{j+1}`$ for $`j=0,2,4`$ and $`s_{j1}`$ for $`j=1,3,5`$. This $`\gamma _1`$ attaches $`x`$ to $`s`$ and $`y`$ to $`t`$. In view of (6.6)–(6.8), we have
$$\begin{array}{c}\zeta ^{\gamma _1}=c_1(g_0)\mu (\gamma _1(z_0)\gamma _1(z_1))+\mathrm{}+c_1(g_5)\mu (\gamma _1(z_5)\gamma _1(z_0))\hfill \\ =h_2d_2+h_00+h_1d_1+h_20+h_0d_0+h_10=h_2(h_0+h_1)+h_1(h_0+h_2)+h_0(h_1+h_2)=\rho .\end{array}$$
Similarly, $`\zeta ^{\gamma _2}=\rho `$ for the symmetric mapping $`\gamma _2`$ which is defined by $`\gamma _2(z_j)=\gamma _1(z_{j+3})`$, attaching $`x`$ to $`t`$ and $`y`$ to $`s`$. We shall see later that $`\gamma _1`$ and $`\gamma _2`$ are just optimal mappings for our gadget.
The mappings pretending to provide (ii) in ((6.1)) are $`\gamma _3`$ and $`\gamma _4`$ illustrated in Fig. 5b,c; here both $`x,y`$ are mapped by $`\gamma _3`$ to $`s`$, and by $`\gamma _4`$ to $`t`$. We have
$$\begin{array}{c}\zeta ^{\gamma _3}=h_2d_2+h_0d_2+h_10+h_2d_2+h_0d_2+h_10=(2h_2+2h_0)(h_0+h_1)\hfill \\ =2h_2h_0+2h_2h_1+2h_0^2+2h_0h_1=\rho +\alpha \end{array}$$
and
$$\begin{array}{c}\zeta ^{\gamma _4}=h_20+h_0d_1+h_1d_1+h_20+h_0d_1+h_1d_1=(2h_0+2h_1)(h_0+h_2)\hfill \\ =2h_0^2+2h_0h_2+2h_1h_0+2h_1h_2=\rho +\alpha .\end{array}$$
Now ((6.1)) is implied by the following.
###### Statement 6.3
Let $`\gamma `$ be a feasible mapping different from $`\gamma _1`$ and $`\gamma _2`$. Then $`\zeta ^\gamma \rho +\alpha `$.
Proof. By (6.6), $`\zeta ^\gamma `$ is representable as a nonnegative integer combination of products $`h_ih_j`$ for $`0i,j5`$ (including $`i=j`$). The contribution $`\zeta _j`$ to $`cm^\gamma `$ from a single edge $`g_j=z_jz_{j+1}`$ is as follows:
1. (i) if $`\gamma (z_j)=\gamma (z_{j+1})=s_{j1}`$, then $`\zeta _j=0`$;
* if $`\gamma (z_j)=s_{j+1}`$ and $`\gamma (z_{j+1})=s_j`$, then $`\zeta _j=h_{j1}d_{j1}=h_{j1}h_j+h_{j1}h_{j+1}`$;
* if $`\gamma (z_j)=s_{j+1}`$ and $`\gamma (z_{j+1})=s_{j1}`$, then $`\zeta _j=h_{j1}d_j=h_{j1}h_{j+1}+h_{j1}^2`$;
* if $`\gamma (z_j)=s_{j1}`$ and $`\gamma (z_{j+1})=s_j`$, then $`\zeta _j=h_{j1}d_{j+1}=h_{j1}h_j+h_{j1}^2`$;
We call $`g_j`$ slanting if it is as in case (iii) or (iv) of ((6.9)). If no edge of $`C`$ is slanting, then $`\gamma `$ is either $`\gamma _1`$ or $`\gamma _2`$. Otherwise $`C`$ contains at least two slanting edges. In this case we observe from ((6.9)) that the representation of $`\zeta ^\gamma `$ includes $`h_i^2+h_j^2`$ (or $`2h_i^2`$) for some $`i,j`$, which is at least $`\alpha `$. Now the result follows from the fact that the representation includes $`2h_ih_j`$ for each $`0i<j2`$.
To see the latter, w.l.o.g., assume $`i=0`$, $`j=2`$, and consider the edges $`g_0`$ and $`g_1`$. By (6.6), $`g_0`$ contributes $`h_0h_2`$ in cases (ii),(iv), i.e., when $`\gamma (z_1)=s_0`$. And if $`\gamma (z_1)=s_2`$, then $`g_1`$ contributes $`h_0h_2`$. Similarly, the pair $`g_3,g_4`$ contributes $`h_0h_2`$.
This completes the proof of Theorem 1.6.
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# Irreducible antifield formalism for reducible constrained Hamiltonian systems
## 1 Introduction
It is well-known that there are two BRST approaches to the quantization of arbitrary gauge theories. One of them is based on the Lagrangian formalism (known as the antifield formalism), while the other is dealing with Hamiltonian aspects . Both formulations can be applied to irreducible, as well as reducible gauge systems. For reducible theories, it is necessary to introduce ghosts of ghosts and their antifields in order to ensure the nilpotency of the BRST symmetry. The antifield treatment was extended to constrained Hamiltonian systems , allowing therefore a clearer connection between the Lagrangian and the Hamiltonian BRST symmetries.
In this paper we give a consistent procedure for quantizing reducible Hamiltonian systems with first-class constraints following an irreducible BRST mechanism. Although the idea of replacing a redundant set of first-class constraints by an irreducible one in a larger phase-space is known , , it has been neither consistently developed nor applied so far to the quantization of reducible gauge theories. Starting with a finite-stage reducible Hamiltonian first-class system, we perform the following steps: (i) we transform the original reducible theory into an irreducible one in a manner that allows the substitution of the BRST quantization of the reducible system by that of the irreducible theory, and (ii) we quantize the extended action of the irreducible system accordingly the antifield-BRST formalism. In consequence, the ghosts of ghosts, as well as their antifields do not appear within our formalism. By virtue of this, our method puts on equal footing the reducible and irreducible constrained Hamiltonian systems from the BRST formalism point of view. As far as we know, such an approach has not been previously published, hence our paper establishes a new result.
The paper is organized in six sections. Section 2 is dealing with enlarging the initial phase-space of an arbitrary first-stage reducible Hamiltonian system by adding some supplementary canonical pairs, and with further constructing an irreducible set of first-class constraints. The irreducible set is derived in a way that ensures the equivalence with the starting first-class set. In Section 3 we establish the physical equivalence between the reducible and irreducible systems. In this light, the physical observables and also the number of physical degrees of freedom associated with both theories are shown to coincide. The physical equivalence allows the replacement of the reducible BRST quantization by the irreducible one. The quantization of the resulting irreducible first-class Hamiltonian system is then performed on behalf of an appropriate gauge-fixing fermion. Section 4 exposes the generalization of our results from the previous sections to finite-stage reducible first-class constraints. In Section 5 we exemplify the general theory in the case of a reducible model describing the Stueckelberg coupling between abelian one- and two-form gauge fields. Section 6 ends the paper with some conclusions.
## 2 First-stage reducible Hamiltonian theories
In this section we show how one can construct a set of irreducible first-class constraints starting from a first-stage reducible one. We begin with a system described by $`N`$ canonical pairs $`(q^i,p_i)`$, subject to the first-class constraints
$$G_{a_0}(q,p)0,a_0=1,\mathrm{},M_0,$$
(1)
which are assumed to be first-stage reducible
$$Z_{a_1}^{a_0}G_{a_0}=0,a_1=1,\mathrm{},M_1,$$
(2)
and suppose that there are no second-class constraints in the theory. In (2) we used the strong equality because one can always define the first-stage reducibility functions such that to have off-shell reducibility. For the sake of simplicity, we assume that $`(q^i,p_i)`$ are bosonic, but the results can be extended to fermions by introducing some appropriate phases. We denote the first-class Hamiltonian by $`H`$, such that the gauge algebra is expressed by $`[G_{a_0},G_{b_0}]=C_{a_0b_0}^{c_0}G_{c_0}`$, $`[H,G_{a_0}]=V_{a_0}^{b_0}G_{b_0}`$. Relations (2) indicate that the functions $`G_{a_0}`$ are not all independent. Under these circumstances, we locally split these functions within the independent and dependent components, $`G_{\overline{a}_0}`$, respectively, $`G_{a_1}`$
$$G_{a_0}=\left(\begin{array}{c}G_{\overline{a}_0}\\ G_{a_1}\end{array}\right),\overline{a}_0=1,\mathrm{},M_0M_1,$$
(3)
with
$$G_{a_1}=M_{a_1}^{\overline{a}_0}G_{\overline{a}_0},$$
(4)
for some functions $`M_{a_1}^{\overline{a}_0}`$, such that $`G_{\overline{a}_0}0G_{a_1}0`$. All that is required is to choose the functions $`G_{a_0}`$ in such a way that the split can be achieved in principle. With the help of (4) we solve (2) with respect to $`Z_{a_1}^{a_0}`$.Accordingly, we find
$$Z_{b_1}^{a_0}=(M_{b_1}^{\overline{a}_0},\delta _{b_1}^{a_1}).$$
(5)
Next, we perform a transformation
$$G_{a_0}\stackrel{~}{G}_{a_0}=\left(\begin{array}{c}G_{\overline{a}_0}\\ \mathrm{𝟎}\end{array}\right),$$
(6)
with the help of an invertible matrix $`M_{a_0}^{b_0}`$,
$$\stackrel{~}{G}_{a_0}=M_{a_0}^{b_0}G_{b_0},$$
(7)
such that $`\stackrel{~}{G}_{a_0}0G_{a_0}0`$. This matrix allows the representation
$$M_{a_0}^{b_0}=\left(\begin{array}{cc}\delta _{\overline{a}_0}^{\overline{b}_0}& \mathrm{𝟎}\\ M_{a_1}^{\overline{b}_0}& \delta _{a_1}^{b_1}\end{array}\right),$$
(8)
while its inverse coincides with itself
$$\overline{M}_{a_0}^{b_0}=M_{a_0}^{b_0}.$$
(9)
If one inverses (7), one gets
$$G_{a_0}=M_{a_0}^{b_0}\stackrel{~}{G}_{b_0},$$
(10)
so, on account of (10) and (2), we consequently find
$$Z_{b_1}^{a_0}G_{a_0}=Z_{b_1}^{a_0}M_{a_0}^{b_0}\stackrel{~}{G}_{b_0}=0.$$
(11)
In this way, we can regard
$$\stackrel{~}{Z}_{c_1}^{b_0}=Z_{c_1}^{a_0}M_{a_0}^{b_0},$$
(12)
as the reducibility functions of $`\stackrel{~}{G}_{b_0}`$. Using (5) and (8) it follows that $`\stackrel{~}{Z}_{c_1}^{b_0}`$ is given by
$$\stackrel{~}{Z}_{c_1}^{b_0}=(\mathrm{𝟎},\delta _{c_1}^{b_1}).$$
(13)
If one splits the free index in (12) like $`b_0=(\overline{b}_0,b_1)`$ and uses (13), one derives (for $`b_0b_1`$)
$$Z_{c_1}^{a_0}M_{a_0}^{b_1}=\delta _{c_1}^{b_1},$$
(14)
hence
$$rank\left(Z_{c_1}^{a_0}M_{a_0}^{b_1}\right)=M_1,$$
(15)
where
$$M_{a_0}^{b_1}=\left(\begin{array}{c}0\\ \delta _{a_1}^{b_1}\end{array}\right).$$
(16)
Next, we transform the reducible constraints (1) into some irreducible ones. In this respect, we introduce a canonical pair $`(y^{a_1},\pi _{a_1})`$ associated with every (free index of) relation (2), which we impose to be constrained by
$$\pi _{a_1}0.$$
(17)
Obviously, the constraints (1) and (17) are first-class and reducible. The theory based on these constraints is physically equivalent with that based only on the constraints (1) as the two systems display the same number of physical degrees of freedom, and, moreover, it can be shown that they describe the same physical observables. Indeed, if $`f`$ denotes an observable of the theory with the constraints (1) and (17), then it is also an observable of the original one. The last statement arises in a simple manner by writing down the equations satisfied by $`f`$, namely,
$$[f,G_{a_0}]0,[f,\pi _{a_1}]0.$$
(18)
The equations (18) show that $`f`$ does not depend (at least weakly) on $`y^{a_1}`$, and, in addition, that the observables associated with this theory fulfill $`[f,G_{a_0}]0`$, which are nothing but the equations verified by the observables corresponding to the original system. The converse is also valid, i.e., any observable of the original theory satisfies (18) because it does not depend on $`(y^{a_1},\pi _{a_1})`$ and checks by definition $`[f,G_{a_0}]0`$.
Dropping out the trivial part of (6), we construct the irreducible first-class constraints
$$\stackrel{~}{\gamma }_{a_0}=\left(\begin{array}{c}G_{\overline{a}_0}\\ \pi _{a_1}\end{array}\right)0,$$
(19)
such that the momenta $`\pi _{a_1}`$ replace the dependent constraint functions. With the help of (6) and (16), the constraints (19) can be put under the form
$$\stackrel{~}{\gamma }_{a_0}=\stackrel{~}{G}_{a_0}M_{a_0}^{b_1}\pi _{b_1}0.$$
(20)
Now, we pass from (20) to the equivalent set of first-class constraints
$$\gamma _{a_0}=M_{a_0}^{b_0}\stackrel{~}{\gamma }_{b_0},$$
(21)
with $`M_{a_0}^{b_0}`$ the matrix (8). Making the notations $`A_{a_0}^{b_1}=M_{a_0}^{b_0}M_{b_0}^{b_1}`$ and using (10), we find from (21) the first-class constraints
$$\gamma _{a_0}G_{a_0}+A_{a_0}^{b_1}\pi _{b_1}0.$$
(22)
The matrix $`A_{a_0}^{b_1}`$ also verifies (14). Indeed, we have that
$$Z_{c_1}^{a_0}A_{a_0}^{b_1}=Z_{c_1}^{a_0}M_{a_0}^{b_0}M_{b_0}^{b_1}=\stackrel{~}{Z}_{c_1}^{b_0}M_{b_0}^{b_1}=\delta _{c_1}^{b_1}.$$
(23)
However, from practical reasons it is useful to weaken the condition (23) by taking $`A_{a_0}^{b_1}`$ such that $`Z_{c_1}^{a_0}A_{a_0}^{b_1}=D_{c_1}^{b_1}`$ is invertible, i.e.,
$$rank\left(D_{c_1}^{b_1}\right)=M_1.$$
(24)
We employ this choice throughout the paper. Moreover, the first-class constraints (22) are irreducible. Indeed, we have that $`Z_{a_1}^{a_0}\gamma _{a_0}=D_{a_1}^{b_1}\pi _{b_1}`$ is non-vanishing due to (24).
Within the above discussion we supposed that the split of the reducible constraints into independent and dependent ones can be done in principle, this assumption being useful for some technical purposes. Indeed, the split form of the original constraints represents an intermediate step in finally reaching the irreducible constraints (22) where the initial constraint functions appear in a covariant (not split) form. The derivation of the constraints (22) based on the above split is still useful in order to evidence the introduction of the intermediate reducible system possessing the constraints (1) and (17), which subsequently emphasizes in a suggestive manner how the dependent constraints can be replaced by some new degrees of freedom ensuring the irreducibility. However, the separation of the reducible constraints can spoil the covariance or destroy the locality of those relations where it is manifest. In fact, the split hypothesis is not crucial in arriving at (22) and can be replaced by homological arguments, as follows. It is well-known that the BRST symmetry $`s_R`$ associated with a Hamiltonian reducible theory contains two basic differentials
$$s_R=\delta _R+D_R+\mathrm{},$$
(25)
where $`\delta _R`$ denotes the Koszul-Tate differential and $`D_R`$ stands for a model of longitudinal derivative along the gauge orbits. In the case of first-stage reducible systems, the action of $`\delta _R`$ on the original phase-space variables and on the generators $`(𝒫_{a_0},𝒫_{a_1})`$ in the Koszul-Tate complex reads as
$$\delta _Rq^i=0,\delta _Rp_i=0,$$
(26)
$$\delta _R𝒫_{a_0}=G_{a_0},$$
(27)
$$\delta _R𝒫_{a_1}=Z_{a_1}^{a_0}𝒫_{a_0},$$
(28)
with $`𝒫_{a_0}`$ and $`𝒫_{a_1}`$ of antighost number one, respectively, two. The antighosts $`𝒫_{a_1}`$ are required in order to kill the antighost number one non-trivial co-cycles
$$\rho _{a_1}Z_{a_1}^{a_0}𝒫_{a_0},$$
(29)
in the homology of $`\delta _R`$. The idea with the help of which we can recover (22) is to redefine the antighosts $`𝒫_{a_0}`$ such that the non-trivial co-cycles of the type (29) vanish identically. If we succeed in doing this, the co-cycles (29) do not appear anymore, hence the antighosts $`𝒫_{a_1}`$ are no longer necessary such that the theory becomes indeed irreducible. In this light, we perform the transformation
$$𝒫_{a_0}\underset{a_0}{\overset{}{𝒫}}=D_{a_0}^{b_0}𝒫_{b_0},$$
(30)
where $`D_{a_0}^{b_0}`$ is chosen such that
$$Z_{a_1}^{a_0}D_{a_0}^{b_0}=0,D_{a_0}^{b_0}G_{b_0}=G_{a_0}.$$
(31)
From (27) and (3031) we obtain that
$$\delta \underset{a_0}{\overset{}{𝒫}}=G_{a_0},$$
(32)
which subsequently leads to
$$\delta \left(Z_{a_1}^{a_0}\underset{a_0}{\overset{}{𝒫}}\right)=0,$$
(33)
but with
$$Z_{a_1}^{a_0}\underset{a_0}{\overset{}{𝒫}}0.$$
(34)
In (3233) we redenoted $`\delta _R`$ by $`\delta `$ in order to outline that the new theory is irreducible. If we take
$$D_{a_0}^{b_0}=\delta _{a_0}^{b_0}Z_{a_1}^{b_0}\overline{D}_{b_1}^{a_1}A_{a_0}^{b_1},$$
(35)
where $`\overline{D}_{b_1}^{a_1}`$ is the inverse of $`D_{b_1}^{a_1}`$, the equations (31) are clearly satisfied. Substituting (35) in (32) we find that
$$\delta \left(𝒫_{a_0}Z_{a_1}^{b_0}\overline{D}_{b_1}^{a_1}A_{a_0}^{b_1}𝒫_{b_0}\right)=G_{a_0}.$$
(36)
As the co-cycles (34) vanish identically it results that (32) or (36) can be precisely associated with an irreducible system. In order to derive the form of the irreducible constraints we consider the new canonical pairs $`(y^{a_1},\pi _{a_1})`$, with $`\pi _{a_1}`$ the non-trivial solutions of the equations
$$D_{a_1}^{b_1}\pi _{b_1}=\delta \left(Z_{a_1}^{b_0}𝒫_{b_0}\right).$$
(37)
The equations (37) may have trivial, as well as non-trivial solutions. Initially, we notice that $`\pi _{b_1}=0`$ (trivial solutions) if and only if $`\delta \left(Z_{a_1}^{b_0}𝒫_{b_0}\right)=0`$. This case corresponds to the reducible theory with the constraints (1) and (17) (in this situation we obtain the co-cycles (29)). The non-trivial solutions $`\pi _{b_1}0`$ appear if and only if $`\delta \left(Z_{a_1}^{b_0}𝒫_{b_0}\right)0`$ (the quantities (29) are no longer co-cycles), hence if and only if the theory is irreducible. While within the split context the momenta $`\pi _{b_1}`$ replace the dependent constraint functions, in the homological approach they enforce the removal of the co-cycles (29). Expressing $`\pi _{b_1}`$ from (37) (in the irreducible case $`\pi _{b_1}0`$)
$$\pi _{b_1}=\delta \left(Z_{a_1}^{b_0}\overline{D}_{b_1}^{a_1}𝒫_{b_0}\right),$$
(38)
and replacing this result in (36) we get the relations
$$\delta 𝒫_{a_0}=G_{a_0}A_{a_0}^{b_1}\pi _{b_1}\gamma _{a_0}.$$
(39)
The last formulas are nothing but the definitions of $`\delta `$ on the antighost number one antighosts $`𝒫_{a_0}`$, that are attached to the irreducible system having the constraints (22). In conclusion, the first-class constraints (22) can be derived by requiring that the non-trivial co-cycles of the type (29) vanish identically under the redefinitions (30). For instance, in the case of free abelian two-form gauge fields the reducible first-class constraints read as
$$G_i^{(2)}2^j\pi _{ji}0.$$
(40)
The model is first-stage redundant, namely, $`Z^iG_i^{(2)}=0`$, with $`Z^i^i`$. The actions of the reducible Koszul-Tate differential on the antighost number one antighosts $`𝒫_i`$ are given by
$$\delta _R𝒫_i=2^j\pi _{ji}.$$
(41)
Redefining $`𝒫_i`$ such that
$$𝒫_i\underset{i}{\overset{}{𝒫}}=(\delta _i^j\frac{_i^j}{\mathrm{}})𝒫_jD_i^j𝒫_j,$$
(42)
from (41) we find
$$\delta \left(𝒫_i\frac{_i^j}{\mathrm{}}𝒫_j\right)=2^j\pi _{ji},$$
(43)
where $`\mathrm{}=_k^k`$. Introducing the canonical pair $`(\phi ,\pi )`$ playing the role of the variables $`(y^{a_1},\pi _{a_1})`$ and taking $`D_{a_1}^{b_1}`$ to be $`\mathrm{}`$, the equations (37) become
$$\mathrm{}\pi =\delta \left(^k𝒫_k\right),$$
(44)
so
$$\pi =\delta \left(\frac{^k}{\mathrm{}}𝒫_k\right).$$
(45)
Substituting (45) in (43), we find the relations
$$\delta 𝒫_i=2^j\pi _{ji}+_i\pi .$$
(46)
In this way the relations (46) emphasize the irreducible first-class constraints
$$\gamma _i2^j\pi _{ji}_i\pi 0,$$
(47)
which appear for instance in the example from Section 5 in the limit $`M=0`$ and in the absence of the fields $`H^\mu `$ (see the first relations in formula (133)).
Now we can show that the constraints (1) and (17) are equivalent with (22), i.e.,
$$\gamma _{a_0}0G_{a_0}0,\pi _{a_1}0.$$
(48)
It is simply to see that if (1) and (17) hold, then the constraints (22) also hold. The converse is valid, too. Indeed, we will see that
$$\gamma _{a_0}0G_{a_0}0,\pi _{a_1}0.$$
(49)
This can be shown as follows. First, we apply $`Z_{c_1}^{a_0}`$ on (22), which then yields
$$\overline{D}_{a_1}^{c_1}Z_{c_1}^{a_0}\gamma _{a_0}=\pi _{a_1}.$$
(50)
With the help of (22) and (50) we get
$$\left(\delta _{a_0}^{b_0}A_{a_0}^{b_1}\overline{D}_{b_1}^{c_1}Z_{c_1}^{b_0}\right)\gamma _{b_0}=G_{a_0}.$$
(51)
From (50-51) we reach (49). The Poisson brackets between the irreducible first-class constraints read
$$[\gamma _{a_0},\gamma _{b_0}]=\overline{C}_{a_0b_0}^{c_0}\gamma _{c_0},$$
(52)
where the new structure functions are expressed by
$`\overline{C}_{a_0b_0}^{c_0}=C_{a_0b_0}^{d_0}\left(\delta _{d_0}^{c_0}A_{d_0}^{b_1}\overline{D}_{b_1}^{c_1}Z_{c_1}^{c_0}\right)+`$
$`\left([G_{a_0},A_{b_0}^{b_1}]+[A_{a_0}^{b_1},G_{b_0}]+[A_{a_0}^{c_1},A_{b_0}^{b_1}]\pi _{b_1}\right)\overline{D}_{b_1}^{c_1}Z_{c_1}^{c_0}.`$ (53)
The first-class Hamiltonian of the new theory can be derived starting from the original one, $`H`$. Indeed, if we take
$$H^{}=H+h^{a_1}(q^i,p_i)\pi _{a_1},$$
(54)
with
$$[h^{a_1},G_{a_0}]=V_{a_0}^{b_0}A_{b_0}^{a_1},$$
(55)
we subsequently find
$$[H^{},\gamma _{a_0}]=\overline{V}_{a_0}^{b_0}\gamma _{b_0},$$
(56)
where
$$\overline{V}_{a_0}^{b_0}=V_{a_0}^{b_0}+\left([H,A_{a_0}^{a_1}]+[h^{b_1},A_{a_0}^{a_1}]\pi _{b_1}\right)\overline{D}_{a_1}^{c_1}Z_{c_1}^{b_0}.$$
(57)
It is clear that the first-class Hamiltonian (54) is not unique because we can always add to it any combinations of $`\gamma _{a_0}`$’s with coefficients that are arbitrary functions. The change induced by the modification of the Hamiltonian gives raise to a change in the structure functions (57).
In brief, in this section we constructed an irreducible first-class system associated with the original redundant one, described by the constraints (22) and the first-class Hamiltonian (54), displaying the gauge algebra (52) and (56). The irreducible theory built here will be important by virtue of the subsequent development.
## 3 Irreducible quantization of the reducible theory
Now, we show that the reducible, respectively, irreducible theories possess the same classical observables. We start from an observable $`F`$ of the irreducible theory. Accordingly, $`F`$ should verify the equations
$$[F,\gamma _{a_0}]0.$$
(58)
On account of (22) and (50), from (58) we deduce
$$[F,G_{a_0}]+[F,\pi _{a_1}]A_{a_0}^{a_1}0.$$
(59)
On the other hand, multiplying (59) with $`Z_{b_1}^{a_0}`$ and using (51), we arrive at
$$[F,\pi _{a_1}]D_{b_1}^{a_1}[F,Z_{b_1}^{a_0}]\left(\delta _{a_0}^{b_0}A_{a_0}^{c_1}\overline{D}_{c_1}^{d_1}Z_{d_1}^{b_0}\right)\gamma _{b_0}0.$$
(60)
Because $`D_{b_1}^{a_1}`$ has maximal rank (see (24)), from (60) we infer
$$[F,\pi _{a_1}]0,$$
(61)
such that
$$[F,\gamma _{a_0}]0[F,G_{a_0}]0.$$
(62)
In conclusion, if $`F`$ is an observable of the irreducible theory, then it is also an observable of the original reducible one. The converse is valid, too, because any observable of the reducible theory verifies the equations $`[F,G_{a_0}]0`$ and does not depend on the newly added canonical variables, such that (58) are indeed satisfied. Thus, both the irreducible and reducible models display the same physical observables. A simple count shows that the numbers of physical degrees of freedom of the reducible, respectively, irreducible theories are both equal to $`NM_0+M_1`$. The last conclusions prove that the original reducible theory is physically equivalent with the irreducible one. This makes permissible the replacement of the BRST quantization for the original redundant system by the BRST quantization of the irreducible theory.
The first attempt at quantizing the irreducible system is to apply the antifield-BRST formalism with respect to its extended action, namely,
$$S_0^E[q^i,p_i,y^{a_1},\pi _{a_1},u^{a_0}]=𝑑t\left(\dot{q}^ip_i+\dot{y}^{a_1}\pi _{a_1}H^{}u^{a_0}\gamma _{a_0}\right).$$
(63)
Action (63) is invariant under the gauge transformations
$$\delta _ϵF=[F,\gamma _{a_0}]ϵ^{a_0},\delta _ϵu^{a_0}=\dot{ϵ}^{a_0}\overline{V}_{b_0}^{a_0}ϵ^{b_0}\overline{C}_{b_0c_0}^{a_0}u^{b_0}ϵ^{c_0},$$
(64)
with $`ϵ^{a_0}`$ the gauge parameters associated with the irreducible constraints (22). In the absence of the newly introduced variables, the extended action (63) together with its gauge transformations, (64), should reduce to those from the reducible case. The gauge transformations of the Lagrange multipliers from (64) do not lead to the corresponding transformations from the reducible situation because the terms $`Z_{a_1}^{a_0}ϵ^{a_1}`$ are missing. The gauge parameters $`ϵ^{a_1}`$, which were attached to the first-stage reducibility functions, are absent within the irreducible approach. In order to restore these terms, it is necessary to further enlarge the phase-space by adding some supplementary canonical pairs $`(z_1^{a_1},p_{1a_1})`$, $`(z_2^{a_1},p_{2a_1})`$, subject to the constraints
$$p_{1a_1}0,p_{2a_1}0.$$
(65)
Obviously, (22) and (65) are still first-class and irreducible. Adding to the first set of constraints from (65) a combination of first-class constraints (see (50)), we obtain the equivalent first-class set
$$\gamma _{a_1}\pi _{a_1}p_{1a_1}0,\gamma _{a_1}^{}p_{2a_1}0.$$
(66)
The additional first-class constraints do not afflict the number of physical degrees of freedom of the former irreducible system. At the same time, the above established equivalence between the physical observables respectively associated with the reducible and irreducible theories remains valid. This is because an observable $`F`$ of the last irreducible model must check, beside (58), the equations
$$[F,\gamma _{a_1}]0,[F,\gamma _{a_1}^{}]0.$$
(67)
On account of (61) relations (67) indicate that in addition to the prior conditions $`F`$ does not depend (at least weakly) on the last added canonical pairs. As it will be seen below, the constraints (66) will imply the presence of the terms $`Z_{a_1}^{a_0}ϵ^{a_1}`$ within the gauge transformations of the Lagrange multipliers $`u^{a_0}`$. In this sense, the constraints (66) play in a certain way the role of the original reducibility relations. The first-class Hamiltonian with respect to the first-class set (22) and (66) can be taken as
$$H_0=H^{}+z_2^{a_1}Z_{a_1}^{a_0}\gamma _{a_0}+y^{a_1}\gamma _{a_1}^{},$$
(68)
such that the new irreducible gauge algebra reads
$$[\gamma _{a_0},\gamma _{b_0}]=\overline{C}_{a_0b_0}^{c_0}\gamma _{c_0},[\gamma _{a_0},\gamma _{a_1}]=0,[\gamma _{a_0},\gamma _{a_1}^{}]=0,$$
(69)
$$[\gamma _{a_1},\gamma _{b_1}]=0,[\gamma _{a_1},\gamma _{b_1}^{}]=0,[\gamma _{a_1}^{},\gamma _{b_1}^{}]=0,$$
(70)
$$[H_0,\gamma _{a_0}]=\stackrel{~}{V}_{a_0}^{b_0}\gamma _{b_0}+A_{a_0}^{b_1}\gamma _{b_1}^{},[H_0,\gamma _{a_1}]=\gamma _{a_1}^{},[H_0,\gamma _{a_1}^{}]=Z_{a_1}^{a_0}\gamma _{a_0}.$$
(71)
Simple calculations show that the functions $`\stackrel{~}{V}_{a_0}^{b_0}`$ from (71) are of the form
$$\stackrel{~}{V}_{a_0}^{b_0}=\overline{V}_{a_0}^{b_0}+z_2^{a_1}\left(\mu _{a_0a_1}^{b_0c_0}G_{c_0}+\lambda _{a_0a_1}^{b_1}Z_{b_1}^{b_0}+[Z_{a_1}^{b_0},A_{a_0}^{b_1}]\pi _{b_1}\right),$$
(72)
with $`\mu _{a_0a_1}^{b_0c_0}`$ and $`\lambda _{a_0a_1}^{b_1}`$ appearing in
$$[Z_{a_1}^{c_0},G_{d_0}]=Z_{a_1}^{a_0}C_{a_0d_0}^{c_0}+\mu _{d_0a_1}^{c_0b_0}G_{b_0}+\lambda _{d_0a_1}^{b_1}Z_{b_1}^{c_0},$$
(73)
and $`\mu _{a_0a_1}^{b_0c_0}`$ antisymmetric in the upper indices, i.e., $`\mu _{a_0a_1}^{b_0c_0}=\mu _{a_0a_1}^{c_0b_0}`$. The relations (73) can be inferred by taking the Poisson brackets between (2) and $`G_{d_0}`$, which leads to
$$\left(Z_{a_1}^{a_0}C_{a_0d_0}^{c_0}+[Z_{a_1}^{c_0},G_{d_0}]\right)G_{c_0}=0.$$
(74)
From (74) and (2) it follows directly (73). We outline that the functions $`\overline{C}_{a_0b_0}^{c_0}`$ and $`\stackrel{~}{V}_{a_0}^{b_0}`$ encode the reducible structure within the irreducible theory. The first-class Hamiltonian (68) is unique up to a combination in terms of the functions $`\gamma _{a_0}`$, $`\gamma _{a_1}`$ and $`\gamma _{a_1}^{}`$. The change of $`H_0`$ consequently implies the change of the structure functions from (71).
The extended action describing the new irreducible theory
$`S_0^E[q^i,p_i,y^{a_1},\pi _{a_1},z_1^{a_1},p_{1a_1},z_2^{a_1},p_{2a_1},u^{a_0},u^{a_1},v^{a_1}]={\displaystyle }dt(\dot{q}^ip_i+`$ (75)
$`\dot{y}^{a_1}\pi _{a_1}+\dot{z}_1^{a_1}p_{1a_1}+\dot{z}_2^{a_1}p_{2a_1}H_0u^{a_0}\gamma _{a_0}u^{a_1}\gamma _{a_1}v^{a_1}\gamma _{a_1}^{}),`$
is invariant under the gauge transformations
$$\delta _ϵF=[F,\gamma _{a_0}]ϵ^{a_0}+[F,\gamma _{a_1}]ϵ_1^{a_1}+[F,\gamma _{a_1}^{}]ϵ_2^{a_1},$$
(76)
$$\delta _ϵu^{a_0}=\dot{ϵ}^{a_0}\stackrel{~}{V}_{b_0}^{a_0}ϵ^{b_0}\overline{C}_{b_0c_0}^{a_0}u^{b_0}ϵ^{c_0}Z_{a_1}^{a_0}ϵ_2^{a_1},$$
(77)
$$\delta _ϵu^{a_1}=\dot{ϵ}_1^{a_1},\delta _ϵv^{a_1}=\dot{ϵ}_2^{a_1}A_{a_0}^{a_1}ϵ^{a_0}ϵ_1^{a_1}.$$
(78)
We emphasize that in this way the terms $`Z_{a_1}^{a_0}ϵ_2^{a_1}`$ are restored within the gauge transformations of the multipliers $`u^{a_0}`$. This is precisely the effect of introducing the supplementary pairs $`(z_1^{a_1},p_{1a_1})`$, $`(z_2^{a_1},p_{2a_1})`$ subject to the constraints (66). If in (75-78) we discard all the newly introduced canonical pairs, we get the extended action and the gauge transformations from the initial redundant case.
With these elements at hand, it appears clearly that we can substitute the quantization of the initial theory by the quantization of the last irreducible system. In the sequel we perform the antifield-BRST quantization with respect to the action (75). To this end, we introduce the ghosts
$$(\eta ^{a_0},\eta _1^{a_1},\eta _2^{a_1}),$$
(79)
and also the antifields
$$(q_i^{},p^i,y_{a_1}^{},\pi ^{a_1},z_{1a_1}^{},p_1^{a_1},z_{2a_1}^{},p_2^{a_1},u_{a_0}^{},u_{a_1}^{},v_{a_1}^{},\eta _{a_0}^{},\eta _{1a_1}^{},\eta _{2a_1}^{}).$$
(80)
The ghosts have ghost number one, the antifields associated with the variables involved with (75) possess ghost number minus one, while the antifields of the ghosts have ghost number minus two. The solution to the master equation is given by
$`S^E=S_0^E+{\displaystyle }dt(q_i^{}[q^i,\gamma _{a_0}]\eta ^{a_0}+p^i[p_i,\gamma _{a_0}]\eta ^{a_0}+`$ (81)
$`y_{a_1}^{}\left(A_{a_0}^{a_1}\eta ^{a_0}+\eta _1^{a_1}\right)z_{1a_1}^{}\eta _1^{a_1}+z_{2a_1}^{}\eta _2^{a_1}+`$
$`u_{a_0}^{}\left(\dot{\eta }^{a_0}\stackrel{~}{V}_{b_0}^{a_0}\eta ^{b_0}\overline{C}_{b_0c_0}^{a_0}u^{b_0}\eta ^{c_0}Z_{a_1}^{a_0}\eta _2^{a_1}\right)+`$
$`u_{a_1}^{}\dot{\eta }_1^{a_1}+v_{a_1}^{}(\dot{\eta }_2^{a_1}A_{a_0}^{a_1}\eta ^{a_0}\eta _1^{a_1}){\displaystyle \frac{1}{2}}\eta _{a_0}^{}\overline{C}_{b_0c_0}^{a_0}\eta ^{b_0}\eta ^{c_0}+\mathrm{}).`$
In order to derive a gauge-fixed action, it is necessary to fix the gauge. In this respect, it is useful to take a gauge-fixing fermion
$$\psi =\psi [q^i,p_i,y^{a_1},\pi _{a_1},z_1^{a_1},p_{1a_1},z_2^{a_1},p_{2a_1},\eta ^{a_0},\eta _1^{a_1},\eta _2^{a_1},u_{a_0}^{},u_{a_1}^{},v_{a_1}^{}],$$
(82)
implementing some irreducible gauge conditions, with the help of which we eliminate all the antifields excepting $`u_{a_0}^{}`$, $`u_{a_1}^{}`$, $`v_{a_1}^{}`$, that are maintained in favour of their fields. The possibility to build some irreducible gauge conditions is easier on behalf of the newly added canonical pairs, which play at this level the same role like the auxiliary variables from the reducible approach. We can put the gauge-fixed action under a form displaying a more direct link with the Hamiltonian BRST quantization of the irreducible system following the procedure exposed in . In this light, we declare the variables $`(\eta ^{a_0},u_{a_0}^{})`$, $`(\eta _1^{a_1},u_{a_1}^{})`$, $`(\eta _2^{a_1},v_{a_1}^{})`$ respectively conjugated in the Poisson bracket
$$[u_{a_0}^{},\eta ^{b_0}]=\delta _{a_0}^{b_0},[u_{a_1}^{},\eta _1^{b_1}]=\delta _{a_1}^{b_1},[v_{a_1}^{},\eta _2^{b_1}]=\delta _{a_1}^{b_1},$$
(83)
and regard the antifields like the momenta associated with the ghosts. Under these circumstances, the gauge-fixed action corresponding to (81) reads as
$`S_\psi ={\displaystyle }dt(\dot{q}^ip_i+\dot{y}^{a_1}\pi _{a_1}+\dot{z}_1^{a_1}p_{1a_1}+\dot{z}_2^{a_1}p_{2a_1}+`$
$`u_{a_0}^{}\dot{\eta }^{a_0}+u_{a_1}^{}\dot{\eta }_1^{a_1}+v_{a_1}^{}\dot{\eta }_2^{a_1}H_B+[\psi ,\mathrm{\Omega }]),`$ (84)
where the BRST charge, respectively, the BRST-extension of the first-class Hamiltonian start like
$$\mathrm{\Omega }=\gamma _{a_0}\eta ^{a_0}+\gamma _{a_1}\eta _1^{a_1}+\gamma _{a_1}^{}\eta _2^{a_1}+\frac{1}{2}u_{a_0}^{}\overline{C}_{b_0c_0}^{a_0}\eta ^{b_0}\eta ^{c_0}+\mathrm{},$$
(85)
$$H_B=H_0+u_{a_0}^{}\left(\stackrel{~}{V}_{b_0}^{a_0}\eta ^{b_0}+Z_{a_1}^{a_0}\eta _2^{a_1}\right)+v_{a_1}^{}\left(A_{a_0}^{a_1}\eta ^{a_0}+\eta _1^{a_1}\right)+\mathrm{}.$$
(86)
This completes our irreducible procedure in the case of first-stage reducible first-class Hamiltonian theories. Until now, we showed how a first-stage reducible first-class Hamiltonian system can be quantized in the framework of the irreducible antifield-BRST formalism, i.e., without introducing ghosts of ghosts.
## 4 $`L`$-stage reducible Hamiltonian theories
In this section we generalize the results from the first-stage case to higher-order-stage reducible systems. If the original Hamiltonian theory is $`L`$-stage reducible (with finite $`L`$), the construction of the corresponding irreducible system goes along the same line like that from the first-stage case. We assume the reducibility relations
$$Z_{a_1}^{a_0}G_{a_0}=0,Z_{a_1}^{a_0}Z_{a_2}^{a_1}=0,\mathrm{},Z_{a_{L1}}^{a_{L2}}Z_{a_L}^{a_{L1}}=0,$$
(87)
with $`a_k=1,\mathrm{},M_k`$. Next, we introduce the canonical pairs $`(y^{a_k},\pi _{a_k})_{k=1,\mathrm{},L}`$ corresponding to the free indices of the above reducibility relations, and constrain these new variables like
$$\pi _{a_k}0.$$
(88)
Constraints (1) and (88) are first-class and obviously reducible. In a manner similar with that from section 2, we derive the first-class constraints
$$\gamma _{a_0}G_{a_0}+A_{a_0}^{b_1}\pi _{b_1}0,$$
(89)
$$\gamma _{a_{2k}}Z_{a_{2k}}^{a_{2k1}}\pi _{a_{2k1}}+A_{a_{2k}}^{a_{2k+1}}\pi _{a_{2k+1}}0,k=1,\mathrm{},a,$$
(90)
$$\overline{\gamma }_{a_{2k}}\pi _{a_{2k}}0,k=1,\mathrm{},a,$$
(91)
which are equivalent with (1) and (88). Acting like in the first-stage situation, we find that
$$\pi _{a_{2k+1}}=m_{a_{2k+1}}^{b_{2k}}\gamma _{b_{2k}},G_{a_0}=m_{a_0}^{b_{2k}}\gamma _{b_{2k}},k=0,\mathrm{},b,$$
(92)
for some appropriate functions $`m_{a_{2k+1}}^{b_{2k}}`$ and $`m_{a_0}^{b_{2k}}`$, such that the equivalence between (1), (88) and (8991) is direct. We employed the notations
$$a=\{\begin{array}{c}\frac{L}{2},\mathrm{for}L\mathrm{even},\\ \frac{L1}{2},\mathrm{for}L\mathrm{odd},\end{array}$$
(93)
$$b=\{\begin{array}{c}\frac{L}{2}1,\mathrm{for}L\mathrm{even},\\ \frac{L1}{2},\mathrm{for}L\mathrm{odd}.\end{array}$$
(94)
In (90) the functions $`A_{a_k}^{a_{k+1}}`$ depend only on $`(q^i,p_i)`$ and possess the property
$$rank\left(Z_{a_k}^{a_{k1}}A_{a_{k1}}^{b_k}\right)\underset{i=k}{\overset{L}{}}()^{k+i}M_i.$$
(95)
Moreover, the $`A_{a_{k1}}^{b_k}`$’s can be taken to satisfy the relations
$$A_{a_{k1}}^{b_k}A_{b_k}^{a_{k+1}}=0.$$
(96)
The last relations are based on the fact that we can always choose the $`A_{a_{k1}}^{b_k}`$’s proportional with the transposed of $`Z_{b_k}^{a_{k1}}`$’s. On account of (96), one finds that the first-class set (8991) is irreducible. We remark that (91) are irreducible. Thus, it remains to be proved that (8990) are so. This can be seen by multiplying (89) with $`Z_{b_1}^{a_0}`$ and (90) with $`Z_{b_{2k+1}}^{a_{2k}}`$, which induce
$$Z_{b_1}^{a_0}\gamma _{a_0}=Z_{b_1}^{a_0}A_{a_0}^{c_1}\pi _{c_1},Z_{b_{2k+1}}^{a_{2k}}\gamma _{a_{2k}}=Z_{b_{2k+1}}^{a_{2k}}A_{a_{2k}}^{a_{2k+1}}\pi _{a_{2k+1}}.$$
(97)
With the help of (97) and (96) we infer that $`Z_{b_1}^{a_0}\gamma _{a_0}=0`$, $`Z_{b_{2k+1}}^{a_{2k}}\gamma _{a_{2k}}=0`$ if and only if
$$\pi _{a_{2k+1}}=A_{a_{2k+1}}^{a_{2k+2}}\nu _{a_{2k+2}},k=0,\mathrm{},b,$$
(98)
with $`\nu _{a_{2k+2}}`$ some functions. Replacing (98) in (8990) we obtain
$$G_{a_0}0,Z_{b_{2k}}^{a_{2k1}}A_{a_{2k1}}^{a_{2k}}\nu _{a_{2k}}0,$$
(99)
which leads, by virtue of (9596), to
$$\nu _{a_{2k}}A_{a_{2k}}^{a_{2k+1}}\lambda _{a_{2k+1}},$$
(100)
for some $`\lambda _{a_{2k+1}}`$. Substituting (100) in (98) we derive that (8990) are reducible with the reducibility functions $`Z_{b_{2k+1}}^{a_{2k}}`$ if and only if $`\pi _{a_{2k+1}}0`$. In this situation the constraints (8990) and (91) are nothing but (1) and (88). Thus, (8991) are reducible with the reducibility functions $`Z_{b_{2k+1}}^{a_{2k}}`$ if and only if they have the form (1) and (88). On the other hand, if one multiplies (90) by $`A_{b_{2k1}}^{a_{2k}}`$, one gets
$$A_{b_{2k1}}^{a_{2k}}\gamma _{a_{2k}}=A_{b_{2k1}}^{a_{2k}}Z_{a_{2k}}^{a_{2k1}}\pi _{a_{2k1}},$$
(101)
due to (96). From (101), it results that $`A_{b_{2k1}}^{a_{2k}}\gamma _{a_{2k}}=0`$ if and only if
$$\pi _{a_{2k1}}=Z_{a_{2k1}}^{a_{2k2}}\mu _{a_{2k2}},$$
(102)
for some $`\mu _{a_{2k2}}`$. Inserting (102) in (90) we find
$$A_{a_{2k}}^{a_{2k+1}}Z_{a_{2k+1}}^{b_{2k}}\mu _{b_{2k}}0,$$
(103)
which leads to
$$\mu _{b_{2k}}Z_{b_{2k}}^{a_{2k1}}\sigma _{a_{2k1}},$$
(104)
for some $`\sigma _{a_{2k1}}`$. Introducing (104) in (102) we conclude that (8991) are reducible with the reducibility functions $`A_{b_{2k1}}^{a_{2k}}`$ if and only if they reduce to (1) and (88). In this way, the irreducibility of (8991) is completely proved.
In the meantime, it is still necessary to add the pairs $`(z_1^{a_{2k+1}},p_{1a_{2k+1}})`$, $`(z_2^{a_{2k+1}},p_{2a_{2k+1}})`$, with $`k=0,\mathrm{},b`$. With the help of the last pairs we build the supplementary first-class constraints
$$\gamma _{a_{2k+1}}\pi _{a_{2k+1}}p_{1a_{2k+1}}0,\gamma _{a_{2k+1}}^{}p_{2a_{2k+1}}0.$$
(105)
The equivalence between the observables of the original redundant theory and those of the irreducible system is gained like in the first-stage situation. We illustrate the proof of the equivalence in the case $`L`$ odd, the other situation being treated in a similar fashion. If $`F`$ stands for an observable of the irreducible system, the conditions $`[F,\overline{\gamma }_{a_{2k}}]0`$ indicate that it does not depend, at least weakly, on $`y^{a_{2k}}`$. In the meantime, $`F`$ should verify
$$[F,\gamma _{a_{2k}}]0,k=0,\mathrm{},a.$$
(106)
We start from the last relation (106) (for $`k=a`$). On account of (92), we obtain
$$[F,\pi _{a_{L2}}]Z_{a_{L1}}^{a_{L2}}+[F,\pi _{a_L}]A_{a_{L1}}^{a_L}0.$$
(107)
Multiplying the above equation by $`Z_{b_L}^{a_{L1}}`$, on behalf of (95), and as $`M_{L+1}=0`$, we derive
$$[F,\pi _{a_L}]0,$$
(108)
such that (107) becomes
$$[F,\pi _{a_{L2}}]Z_{a_{L1}}^{a_{L2}}0.$$
(109)
Multiplying the next equation from (106) (for $`k=a1`$) with $`Z_{b_{L2}}^{a_{L3}}`$, we further infer
$$[F,\pi _{a_{L2}}]A_{a_{L3}}^{a_{L2}}Z_{b_{L2}}^{a_{L3}}0.$$
(110)
With the help of (96), from (110) we deduce
$$[F,\pi _{a_{L2}}]=n_{b_{L1}}A_{a_{L2}}^{b_{L1}},$$
(111)
for some functions $`n_{b_{L1}}`$. Replacing (111) in (109) it follows that the weak relations $`n_{b_{L1}}A_{a_{L2}}^{b_{L1}}Z_{a_{L1}}^{a_{L2}}0`$ imply
$$n_{b_{L1}}\rho _{b_L}A_{b_{L1}}^{b_L},$$
(112)
for some $`\rho _{b_L}`$. Inserting (112) in (111) we get
$$[F,\pi _{a_{L2}}]0,$$
(113)
due to (96). Reprising the same steps on the remaining equations (106) we consequently arrive to
$$[F,\pi _{a_{L2k}}]0,$$
(114)
which lead to
$$[F,G_{a_0}]0.$$
(115)
Moreover, the equations $`[F,\gamma _{a_{2k+1}}]0`$ and $`[F,\gamma _{a_{2k+1}}^{}]0`$ express the fact that $`F`$ does not depend on the $`z`$’s. Thus, any observable of the irreducible theory does not involve, at least weakly, the newly introduced variables, and, in addition, it satisfies (115), which are nothing but the equations that should be checked by any observable of the original redundant system, which show that $`F`$ is also an observable of the reducible theory. Conversely, if $`F^{}`$ denotes an observable associated with the reducible system, then it is obviously an observable of the irreducible theory.
The first-class Hamiltonian with respect to the irreducible first-class constraints (8991) and (105) can be taken under the form
$`H_0^{}=H^{}+{\displaystyle \underset{k=1}{\overset{a}{}}}y^{a_{2k}}\gamma _{a_{2k}}+{\displaystyle \underset{k=0}{\overset{b}{}}}y^{a_{2k+1}}p_{2a_{2k+1}}+`$
$`{\displaystyle \underset{k=0}{\overset{b}{}}}z_2^{a_{2k+1}}\left(Z_{a_{2k+1}}^{a_{2k}}\gamma _{a_{2k}}+A_{a_{2k+1}}^{a_{2k+2}}\gamma _{a_{2k+2}}\right),`$ (116)
with $`H^{}`$ given by (54), where we understood the convention $`f^{a_k}=0`$ if $`k<0`$ or $`k>L`$. The first-class Hamiltonian (4) is again unique up to adding a combination in the first-class constraint functions.
With all these elements at hand, the quantization of the irreducible theory goes from now on along the standard antifield-BRST rules. The ghost spectrum contains only the ghost number one variables associated with the corresponding constraint functions
$$\eta ^{a_0}\gamma _{a_0},\eta _1^{a_{2k}}\overline{\gamma }_{a_{2k}},\eta _2^{a_{2k}}\gamma _{a_{2k}},k=1,\mathrm{},a,$$
(117)
$$\eta _1^{a_{2k+1}}\gamma _{a_{2k+1}},\eta _2^{a_{2k+1}}\gamma _{a_{2k+1}}^{},k=0,\mathrm{},b,$$
(118)
while the antifield sector is given by
$$(q_i^{},p^i),(y_{a_k}^{},\pi ^{a_k})_{k=1,\mathrm{},L},(z_{1a_{2k+1}}^{},p_1^{a_{2k+1}})_{k=0,\mathrm{},b},(z_{2a_{2k+1}}^{},p_2^{a_{2k+1}})_{k=0,\mathrm{},b},$$
(119)
$$u_{a_0}^{},(u_{a_{2k}}^{},v_{a_{2k}}^{})_{k=1,\mathrm{},a},(u_{a_{2k+1}}^{},v_{a_{2k+1}}^{})_{k=0,\mathrm{},b},$$
(120)
$$\eta _{a_0}^{},(\eta _{1a_{2k}}^{},\eta _{2a_{2k}}^{})_{k=1,\mathrm{},a},(\eta _{1a_{2k+1}}^{},\eta _{2a_{2k+1}}^{})_{k=0,\mathrm{},b}.$$
(121)
The antifields $`(u_{a_{2k}}^{},v_{a_{2k}}^{})`$ correspond to the Lagrange multipliers of the constraint functions $`\overline{\gamma }_{a_{2k}}`$, respectively, $`\gamma _{a_{2k}}`$, and $`(u_{a_{2k+1}}^{},v_{a_{2k+1}}^{})`$ are associated with $`\gamma _{a_{2k+1}}`$, respectively, $`\gamma _{a_{2k+1}}^{}`$. The variables (119-120) have ghost number minus one, while in (121) there appear only ghost number minus two antifields. The gauge-fixing fermion should be taken to depend on the $`\mathrm{\Phi }^A`$’s, on the ghosts, and also on the antifields of the Lagrange multipliers, where
$$\mathrm{\Phi }^A=(q^i,p_i,y^{a_k},\pi _{a_k},z_1^{a_{2k+1}},p_{1a_{2k+1}},z_2^{a_{2k+1}},p_{2a_{2k+1}}).$$
(122)
With the help of the gauge-fixing fermion we eliminate all the antifields except the antifields of the multipliers, and also the Lagrange multipliers. The gauge-fixed action will be expressed by
$`S_\psi ={\displaystyle }dt(\dot{q}^ip_i+{\displaystyle \underset{k=1}{\overset{L}{}}}\dot{y}^{a_k}\pi _{a_k}+{\displaystyle \underset{k=0}{\overset{b}{}}}(\dot{z}_1^{a_{2k+1}}p_{1a_{2k+1}}+\dot{z}_2^{a_{2k+1}}p_{2a_{2k+1}})+`$
$`u_{a_0}^{}\dot{\eta }^{a_0}+{\displaystyle \underset{k=1}{\overset{L}{}}}(u_{a_k}^{}\dot{\eta }_1^{a_k}+v_{a_k}^{}\dot{\eta }_2^{a_k})H_B+[\psi ,\mathrm{\Omega }]),`$ (123)
where the BRST charge and the BRST extension of the first-class Hamiltonian (4) respectively start like
$`\mathrm{\Omega }=\gamma _{a_0}\eta ^{a_0}+{\displaystyle \underset{k=1}{\overset{a}{}}}\left(\overline{\gamma }_{a_{2k}}\eta _1^{a_{2k}}+\gamma _{a_{2k}}\eta _2^{a_{2k}}\right)+`$
$`{\displaystyle \underset{k=0}{\overset{b}{}}}\left(\gamma _{a_{2k+1}}\eta _1^{a_{2k+1}}+\gamma _{a_{2k+1}}^{}\eta _2^{a_{2k+1}}\right)+{\displaystyle \frac{1}{2}}u_{a_0}^{}\overline{C}_{b_0c_0}^{a_0}\eta ^{b_0}\eta ^{c_0}+`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{a}{}}}{\displaystyle \underset{k=1}{\overset{a}{}}}{\displaystyle \underset{i=1}{\overset{a}{}}}v_{a_{2j}}^{}\overline{C}_{b_{2k}c_{2i}}^{a_{2j}}\eta _2^{b_{2k}}\eta _2^{c_{2i}}+\mathrm{},`$ (124)
$`H_B=H_0^{}+u_{a_0}^{}\left(\stackrel{~}{V}_{b_0}^{a_0}\eta ^{b_0}+Z_{a_1}^{a_0}\eta _2^{a_1}\right)+`$
$`{\displaystyle \underset{k=1}{\overset{a}{}}}v_{a_{2k}}^{}\left(\eta _1^{a_{2k}}+{\displaystyle \underset{j=1}{\overset{a}{}}}\stackrel{~}{V}_{b_{2j}}^{a_{2k}}\eta _2^{b_{2j}}+Z_{a_{2k+1}}^{a_{2k}}\eta _2^{a_{2k+1}}+A_{a_{2k1}}^{a_{2k}}\eta _2^{a_{2k1}}\right)+`$
$`v_{a_1}^{}\left(\eta _1^{a_1}+A_{a_0}^{a_1}\eta ^{a_0}+Z_{a_2}^{a_1}\eta _2^{a_2}\right)+`$
$`{\displaystyle \underset{k=1}{\overset{b}{}}}v_{a_{2k+1}}^{}\left(\eta _1^{a_{2k+1}}+A_{a_{2k}}^{a_{2k+1}}\eta ^{a_{2k}}+Z_{a_{2k+2}}^{a_{2k+1}}\eta _2^{a_{2k+2}}\right)+\mathrm{}.`$ (125)
The functions $`\overline{C}_{b_{2k}c_{2i}}^{a_{2j}}`$ and $`\stackrel{~}{V}_{b_{2k}}^{a_{2j}}`$ are those involved with the irreducible gauge algebra arising in the $`L`$-stage reducible case. In this way, we realized the BRST quantization of arbitrary $`L`$-stage reducible first-class Hamiltonian systems in an irreducible manner, i.e., without introducing ghosts of ghosts. This completes our analysis.
## 5 Example
Here we exemplify the general theory exposed above in the case of abelian one- and two-form gauge fields with Stueckelberg coupling. We start with the Lagrangian action
$$S_0^L[H^\mu ,A^{\mu \nu }]=d^4x\left(\frac{1}{12}F_{\mu \nu \rho }^2+\frac{1}{4}\left(MA_{\mu \nu }F_{\mu \nu }\right)^2\right),$$
(126)
where $`F_{\mu \nu }`$ and $`F_{\mu \nu \rho }`$ denote the field strengths associated with $`H_\mu `$, respectively, $`A_{\mu \nu }`$, and the notation $`F_{\mu \nu \rho }^2`$ signifies $`F_{\mu \nu \rho }F^{\mu \nu \rho }`$. (We used a similar notation for the other square.) The system described by action (126) possesses the first-class constraints
$$G_i^{\left(1\right)}\pi _{0i}0,G^{\left(1\right)}\mathrm{\Pi }_00,$$
(127)
$$G_i^{\left(2\right)}2^l\pi _{li}+M\mathrm{\Pi }_i0,G^{\left(2\right)}^i\mathrm{\Pi }_i0,$$
(128)
and the first-class Hamiltonian
$`H={\displaystyle }d^3x(\pi _{ij}^2{\displaystyle \frac{1}{2}}\mathrm{\Pi }_i^2+A^{0i}G_i^{\left(2\right)}+H^0G^{\left(2\right)}+`$
$`{\displaystyle \frac{1}{12}}F_{ijk}^2+{\displaystyle \frac{1}{4}}(MA_{ij}F_{ij})^2).`$ (129)
In (1275), the $`\pi `$’s and $`\mathrm{\Pi }`$’s are the canonical momenta associated with the corresponding $`A`$’s and $`H`$’s. The first-class constraints (128) are first-stage reducible
$$^iG_i^{\left(2\right)}+MG^{\left(2\right)}=0,$$
(130)
with the reducibility functions
$$Z_{a_1}^{a_0}=(^i,M).$$
(131)
The functions $`A_{a_0}^{a_1}`$ read
$$A_{a_0}^{a_1}=\left(\begin{array}{c}_i\\ M\end{array}\right),$$
(132)
such that $`Z_{a_1}^{a_0}A_{a_0}^{b_1}=\left(^i_i+M^2\right)`$ is invertible. The variables $`(y^{a_1},\pi _{a_1})`$ will be denoted in this case by $`(\phi ,\pi )`$. The irreducible first-class constraints are given by (127) and
$$\gamma _i^{\left(2\right)}2^l\pi _{li}+M\mathrm{\Pi }_i_i\pi 0,\gamma ^{\left(2\right)}^i\mathrm{\Pi }_iM\pi 0,$$
(133)
while the first-class Hamiltonian $`H^{}`$ (see (54)) reads
$$H^{}=H+d^3x\left(A^{0i}_i\pi MH^0\pi \right).$$
(134)
We introduce the pairs $`(\phi _1,\pi _1)`$, $`(\phi _2,\pi _2)`$, and set the constraints
$$\gamma \pi \pi _10,\gamma ^{}\pi _20.$$
(135)
The momentum $`\pi `$ is indeed a combination of the first-class constraints (133)
$$\pi =\frac{1}{\left(^i_i+M^2\right)}\left(^i\gamma _i^{\left(2\right)}+M\gamma ^{\left(2\right)}\right).$$
(136)
The first-class Hamiltonian with respect to (127), (133) and (135) has the form
$$H_0=H^{}+d^3x\left(\phi _2\left(^i\gamma _i^{\left(2\right)}+M\gamma ^{\left(2\right)}\right)\phi \pi _2\right)d^3xh_0.$$
(137)
The extended action
$`S_0^E={\displaystyle }d^4x(\dot{A}^{0i}\pi _{0i}+\dot{A}^{ij}\pi _{ij}+\dot{H}^0\mathrm{\Pi }_0+\dot{H}^i\mathrm{\Pi }_i+\dot{\phi }\pi +\dot{\phi }_1\pi _1+`$
$`\dot{\phi }_2\pi _2h_0u^iG_i^{\left(1\right)}uG^{\left(1\right)}u^{}\gamma v^i\gamma _i^{\left(2\right)}v\gamma ^{\left(2\right)}v^{}\gamma ^{}),`$ (138)
is invariant under the gauge transformations
$$\delta _ϵA^{0i}=ϵ_1^i,\delta _ϵH^0=ϵ_1,\delta _ϵA^{ij}=^{[i}ϵ_2^{j]},\delta _ϵH^i=^iϵ_2+Mϵ_2^i,$$
(139)
$$\delta _ϵ\phi =_iϵ_2^iMϵ_2+\stackrel{~}{ϵ}_1,\delta _ϵ\phi _1=\stackrel{~}{ϵ}_1,\delta _ϵ\phi _2=\stackrel{~}{ϵ}_2,\delta _ϵu=\dot{ϵ}_1,\delta _ϵu^{}=\underset{1,}{\overset{.}{\stackrel{~}{ϵ}}}$$
(140)
$$\delta _ϵv^i=\dot{ϵ}_2^i^i\stackrel{~}{ϵ}_2ϵ_1^i,\delta _ϵv=\dot{ϵ}_2+M\stackrel{~}{ϵ}_2ϵ_1,\delta _ϵv^{}=\underset{2}{\overset{.}{\stackrel{~}{ϵ}}}+_iϵ_2^iMϵ_2+\stackrel{~}{ϵ}_1,$$
(141)
the gauge variations of all the momenta being identically vanishing. In (139141) the gauge parameters $`ϵ_1^i`$, $`ϵ_1`$, $`ϵ_2^i`$, $`ϵ_2`$, $`\stackrel{~}{ϵ}_1`$ and $`\stackrel{~}{ϵ}_2`$ are respectively associated with the constraint functions $`G_i^{\left(1\right)}`$, $`G^{\left(1\right)}`$, $`\gamma _i^{\left(2\right)}`$, $`\gamma ^{\left(2\right)}`$, $`\gamma `$ and $`\gamma ^{}`$. From (139141) we can derive the Lagrangian gauge transformations associated with the irreducible theory (including, of course, the gauge transformations of the original fields). In view of this we should consider a model of irreducible Hamiltonian theory. In this light we assume that (127) and the former constraint in (135) are primary, while (133) and the latter constraint from (135) are secondary. Passing from the extended action (5) to the corresponding total one (obtained by taking $`v^i=0`$, $`v=0`$ and $`v^{}=0`$ in (5)) we derive its gauge invariances in the standard manner. Indeed, the equations $`v^i=0`$, $`v=0`$ and $`v^{}=0`$ imply $`\delta _ϵv^i=0`$, $`\delta _ϵv=0`$ and $`\delta _ϵv^{}=0`$. The last three equations lead via (141) to
$$ϵ_1^i=\dot{ϵ}_2^i^i\stackrel{~}{ϵ}_2,ϵ_1=\dot{ϵ}_2+M\stackrel{~}{ϵ}_2,\stackrel{~}{ϵ}_1=\underset{2}{\overset{.}{\stackrel{~}{ϵ}}}_iϵ_2^i+Mϵ_2.$$
(142)
Replacing $`ϵ_1^i`$, $`ϵ_1`$ and $`\stackrel{~}{ϵ}_1`$ from (142) in (139140) we get
$$\delta _ϵA^{0i}=\dot{ϵ}_2^i^i\stackrel{~}{ϵ}_2,\delta _ϵH^0=\dot{ϵ}_2+M\stackrel{~}{ϵ}_2,\delta _ϵA^{ij}=^{[i}ϵ_2^{j]},\delta _ϵH^i=^iϵ_2+Mϵ_2^i,$$
(143)
$$\delta _ϵ\phi =\underset{2}{\overset{.}{\stackrel{~}{ϵ}}},\delta _ϵ\phi _1=\underset{2}{\overset{.}{\stackrel{~}{ϵ}}}+_iϵ_2^iMϵ_2,\delta _ϵ\phi _2=\stackrel{~}{ϵ}_2,$$
(144)
$$\delta _ϵu=\ddot{ϵ}_2+M\underset{2}{\overset{.}{\stackrel{~}{ϵ}}},\delta _ϵu^{}=\underset{2}{\overset{..}{\stackrel{~}{ϵ}}}_i\dot{ϵ}_2^i+M\dot{ϵ}_2.$$
(145)
The Lagrangian action corresponding to the above total action coincides with the original one, and its gauge transformations, which derive from (143144), read as
$$\delta _ϵA^{\mu \nu }=^\mu ϵ^\nu ^\nu ϵ^\mu ,\delta _ϵH^\mu =^\mu ϵ+Mϵ^\mu ,\delta _ϵ\phi _1=_\mu ϵ^\mu Mϵ.$$
(146)
The gauge transformations for $`\phi `$ and $`\phi _2`$ were omitted as these fields play the role of Lagrange multipliers (see (137)) and are not relevant in the Lagrangian context. In order to write down (146) we employed the notations
$$ϵ^\mu =(\stackrel{~}{ϵ}_2,ϵ_2^i),ϵ=ϵ_2.$$
(147)
In consequence, our formalism reproduces via (139141) the original gauge transformations and outputs some new gauge transformations (for $`\phi _1`$) that make the gauge transformation set (146) irreducible also at the Lagrangian level. The Lorentz covariance of the gauge transformations (146) is due to the introduction in the theory of the pairs $`(\phi _1,\pi _1)`$ and $`(\phi _2,\pi _2)`$.
In the sequel we approach the antifield BRST treatment of (5). Straightforward calculation then yield the solution to the master equation
$`S^E=S_0^E+{\displaystyle }d^4x(A_{0i}^{}\eta _1^i+H_0^{}\eta _1+A_{ij}^{}^{[i}\eta _2^{j]}+H_i^{}(^i\eta _2+M\eta _2^i)+`$
$`\phi ^{}\left(_i\eta _2^iM\eta _2+\stackrel{~}{\eta }_1\right)\phi _1^{}\stackrel{~}{\eta }_1+\phi _2^{}\stackrel{~}{\eta }_2+u_i^{}\dot{\eta }_1^i+u^{}\dot{\eta }_1+u^{}\underset{1}{\overset{.}{\stackrel{~}{\eta }}}+`$
$`v_i^{}\left(\dot{\eta }_2^i^i\stackrel{~}{\eta }_2\eta _1^i\right)+v^{}\left(\dot{\eta }_2+M\stackrel{~}{\eta }_2\eta _1\right)+`$
$`v^{}(\underset{2}{\overset{.}{\stackrel{~}{\eta }}}+_i\eta _2^iM\eta _2+\stackrel{~}{\eta }_1)).`$ (148)
All the ghosts from (5) have ghost number one, and all the antifields ghost number minus one. We take the gauge fixing fermion
$`\psi ={\displaystyle }d^4x(u_i^{}(_jA^{ji}+MH^i+^i\phi _1)+u^{}(_iH^iM\phi _1)`$ (149)
$`u^{}(_jA^{j0}+MH^0)),`$
which implements the irreducible gauge conditions $`_jA^{ji}+MH^i+^i\phi _1=0`$, $`_iH^iM\phi _1=0`$, and $`_jA^{j0}+MH^0=0`$. After some computation we are led to the gauge-fixed action
$`S_\psi ^E=S_0^L+{\displaystyle }d^4x(B_\mu (_\nu A^{\nu \mu }+MH^\mu +^\mu \phi _1)+b(_\nu H^\nu M\phi _1)+`$
$`u_\mu ^{}(\mathrm{}+M^2)\eta _2^\mu +u^{}(\mathrm{}+M^2)\eta _2),`$ (150)
such that the resulting path integral is given by
$$Z_\psi =𝒟A^{\nu \mu }𝒟H^\mu 𝒟B_\mu 𝒟\phi _1𝒟u_\mu ^{}𝒟\eta _2^\mu 𝒟u^{}𝒟\eta _2\mathrm{exp}iS_\psi ^E.$$
(151)
In (5-151) we employed the identifications
$$B_\mu =(\pi _1,\pi _{0i}),b=\mathrm{\Pi }_{0,}u_\mu ^{}=(u^{},u_j^{})\eta _2^\mu =(\stackrel{~}{\eta }_2,\eta _2^j).$$
(152)
One can check that there are no residual gauge invariances in action (5). Moreover, the gauge-fixed action (5) is Lorentz covariant. This is due precisely to the introduction in the theory of the pairs $`(\phi _1,\pi _1)`$ and $`(\phi _2,\pi _2)`$ subject to the constraints (135).
While the gauge-fixing fermion (149) is useful in getting the covariant path integral (151), the fermion
$`\psi ^{}={\displaystyle }d^4x(u_i^{}(_jA^{ji}\dot{A}^{0i}+^i\phi _1)+`$
$`u^{}(_iH^i\dot{H}^0)u^{}(_jA^{j0}\dot{\phi }_1)),`$ (153)
is appropriate in order to make the reduction to the physical degrees of freedom in the path integral. Starting with the solution (5) and on behalf of (5) we find after some computation the path integral over physical degrees of freedom for the model under consideration of the form
$`Z_\psi ^{}={\displaystyle }𝒟A^{ij}𝒟\pi _{ij}𝒟H^i𝒟\mathrm{\Pi }_i𝒟\phi _1𝒟\pi _1\delta (^j\pi _{ji}+_i\pi _1)\times `$
$`\delta \left(_jA^{ji}+^i\phi _1\right)\delta \left(^i\mathrm{\Pi }_i\right)\delta \left(_iH^i\right)\mathrm{exp}i\overline{S}_\psi ^{},`$ (154)
where $`\overline{S}_\psi ^{}`$ is given by
$$\overline{S}_\psi ^{}=d^4x\left(\dot{A}^{ij}\pi _{ij}+\dot{H}^i\mathrm{\Pi }_i+\pi _{ij}^2+\frac{1}{2}\mathrm{\Pi }_i^2\frac{1}{12}F_{ijk}^2\frac{1}{4}\left(MA_{ij}F_{ij}\right)^2\right).$$
(155)
The delta functions from the constraint functions and their gauge conditions in the path integral (5) show that the independent fields and momenta are precisely the transverse components of $`H^i`$ and $`\mathrm{\Pi }_i`$ and also the longitudinal components of $`A^{ij}`$ and $`\pi _{ij}`$. It is clear that the conditions $`^i\mathrm{\Pi }_i=0`$ and $`_iH^i=0`$ restrict the integration only over the two transverse degrees of freedom for the vector fields and their momenta (typically for electromagnetism). Related to the remaining conditions from the measure of (5), it can be shown that they enforce the longitudinal parts as independent components of the tensor fields and their momenta. Indeed, $`A^{ij}`$ and $`\pi _{ij}`$ can be decomposed into longitudinal and transverse components
$$A_{ij}=_iA_j^T_jA_i^T+\epsilon _{ijk}^kA^L,\pi _{ij}=_i\pi _j^T_j\pi _i^T+\epsilon _{ijk}^k\pi ^L,$$
(156)
where the transverse components satisfy $`^iA_i^T=0`$ and $`^i\pi _i^T=0`$. Then, the conditions $`_jA^{ji}+^i\phi _1=0`$ and $`^j\pi _{ji}+_i\pi _1=0`$ imply via (156) that
$$^i_iA_j^T+_j\phi _1=0,^i_i\pi _j^T+_j\pi _1=0,$$
(157)
hence
$$A_j^T=\frac{_j}{\mathrm{}}\phi _1,\pi _j^T=\frac{_j}{\mathrm{}}\pi _1.$$
(158)
On the other hand, from (157) it follows that $`^i_i\phi _1=0`$ and $`^i_i\pi _1=0`$, which then yield $`\phi _1=0`$, $`\pi _1=0`$ by virtue of the boundary conditions for the unphysical degrees of freedom $`(\phi _1,\pi _1)`$ (vacuum to vacuum). Inserting the last relations back in (158) we find that the conditions checked by the tensor fields and their momenta lead to $`A_j^T=0`$ and $`\pi _j^T=0`$, so the only physical degrees of freedom are described by the longitudinal pair $`(A^L,\pi ^L)`$. In this way the conditions implemented in the measure of (5) lead to transverse degrees of freedom for the vector fields, respectively to a longitudinal one for the tensor fields, like in the reducible approach. This completes the analysis of the investigated model.
## 6 Conclusion
In conclusion, we succeeded in giving a systematic irreducible procedure of quantizing reducible first-class Hamiltonian systems accordingly the antifield BRST method. This new result was inferred by means of constructing an irreducible first-class Hamiltonian theory in a larger phase-space that remains physically equivalent to the original redundant one. The above equivalence makes legitimate the replacement of the quantization of the reducible theory by that of the irreducible system. As a consequence of our irreducible approach, the ghosts of ghosts, their antifields, as well as the pyramidal structure of auxiliary fields are no longer necessary. We further illustrate in detail the theoretical part of the paper in the case of the Stueckelberg coupled abelian one- and two-form gauge fields.
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# The Effect of Hydrostatic Weighting on the Vertical Temperature Structure of the Solar Corona
## 1 INTRODUCTION
Attempts to solve the elusive coronal heating problem have been undertaken by determining the heating function $`E_H(h)`$ as function of height $`h`$, inferred from the vertical temperature structure $`T_e(h)`$ of the solar corona. In this context, a systematic temperature increase $`T(h)`$ with height $`h`$ has been reported from numerous observations of the quiet diffuse corona, coronal arcades, or coronal loops (Mariska & Withbroe 1978; Kohl et al. 1980; Falconer 1994; Foley et al. 1996; Sturrock, Wheatland, and Acton 1996a; 1996b; Wheatland, Sturrock, & Acton 1997; Fludra et al. 1999; Priest et al. 1999; 2000). A common method that is chosen to infer the vertical temperature structure $`T_e(h)`$ is the extraction of soft X-ray fluxes in different wavelengths as function of height, say $`F_1(h)`$ and $`F_2(h)`$ from two different wavelengths 1 and 2, and then to use the filter-ratio method $`Q(h)=F_2(h)/F_1(h)`$ to determine the temperature as function of height, $`T(h)`$, by inverting the filter-ratio function $`Q(T)`$. The filter-ratio method has some obvious limitations, such as the limited range where the function $`Q(T)`$ is unique and thus permits only an inversion within this range, but the method has also some more subtle drawbacks in the case of a multi-temperature plasma, as it exists in the solar corona. In principle, the filter-ratio method is only exact for an isothermal plasma, within the uniqueness range of $`Q(T)`$. The solar corona consists of myriads of open and closed field lines filled with plasmas of almost every temperature in the range of $`10^4\stackrel{<}{}T\stackrel{<}{}10^7`$ K, which is usually quantified with a differential emission measure distribution $`dEM(T)/dT`$. This multi-thermal nature can cause systematic errors in the determination of an average vertical temperature profile $`T_e(h)`$, due to a systematic weighting bias of the temperature-dependent pressure and density scale heights (Fig.1). The purpose of this Letter is to demonstrate this systematic error in the determination of the vertical temperature profile $`T_e(h)`$, for some typical observations of active regions and the quiet corona, using a broadband-filter instrument, such as the Yohkoh Soft X-Ray Telescope (SXT).
## 2 MODEL
The soft X-ray flux measured along a given line-of-sight represents an integral over emission measure contributions from plasmas with different temperatures, which can be expressed by the differential emission measure distribution $`dEM(T)/dT`$, where the emission measure contribution at a given temperature $`[T,T+dT]`$ itself represents an integration along the line-of-sight $`z`$,
$$\left(\frac{dEM(T)}{dT}\right)dT=n_e^2(T,z)𝑑z.$$
(1)
The flux measured by a detector $`i`$ is then given by the product of the differential emission measure function $`dEM(T)/dT`$ with the instrumental temperature response function $`R_i(T)`$,
$$F_i=\frac{dEM(T)}{dT}R_i(T)𝑑T.$$
(2)
We characterize now the solar corona by a superposition of many different flux tubes (along open or closed magnetic field lines), each one having its own temperature and density function. For the purpose of this demonstration we make the simplest assumption that is compatible with observations, namely (1) that each flux tube is near-isothermal (as it has been established for many observed EUV loops in the temperature range of $`T_e1.02.0`$ MK, e.g. Neupert et al. 1998; Lenz et al. 1999; Aschwanden et al. 1999; 2000a; 2000b), and that each flux tube is in near-hydrostatic equilibrium (a condition that has been verified for EUV loops within factors of $``$1-3, Schrijver et al. 1999; Aschwanden et al. 1999; 2000a; 2000b). Thus, the density structure of a (near-isothermal) fluxtube can be approximated by
$$n_e(h,T_e)=n_{e0}\mathrm{exp}[\frac{h}{\lambda (T_e)}].$$
(3)
where the density (or pressure) scale height $`\lambda (T_e)`$ in hydrostatic equilibrium is proportional to the temperature $`T_e`$,
$$\lambda (T_e)=\frac{k_BT_e}{\mu m_pg_{\mathrm{}}}=\lambda _0(\frac{T_e}{1\mathrm{MK}})$$
(4)
with $`\lambda _0=47`$ Mm for coronal conditions, with $`\mu m_p`$ the average ion mass (i.e. $`\mu 1.4`$ for H:He=10:1), and $`g_{\mathrm{}}`$ the solar gravitation. The differential emission measure $`dEM(T,h)/dT`$ is proportional to $`n_e(h)^2`$, and thus has approximatly an exponential height dependence with a scale height half of the density scale height. In first order, the height dependence of the line-of-sight integrated emission measure of plasma with a particular temperature $`T`$ can then be characterized by the half density scale height, neglecting the curvature of the corona. The fluxes $`F_1(h)`$ and $`F_2(h)`$ recorded in two detectors (i=1,2) with different wavelengths, characterized by the temperature response functions $`R_1(T)`$ and $`R_2(T)`$, is then
$$F_i(h)=\frac{dEM(T,h=0)}{dT}\mathrm{exp}[\frac{2h}{\lambda _0T}]R_i(T)𝑑T,$$
(5)
When the filter-ratio method is applied, one takes the flux ratio of the two fluxes at every pixel (along a chosen altitude path $`h`$)
$$Q(h)=\frac{F_1(h)}{F_2(h)},$$
(6)
which is now height-dependent, so that the resulting filter-ratio temperature $`T(Q)=T(Q[h])`$ yields the height dependence of the temperature, $`T(h)`$.
## 3 OBSERVATIONS AND SIMULATION
Some typical differential emission measure distributions $`dEM(T)/dT`$ have been determined with the NASA/GSFC Solar EUV Rocket Telescope and Spectrograph (SERTS), using densitiy-sensitive line ratios from 8 different ionization states of iron between Fe<sup>+9</sup> (Fe X) and Fe<sup>+16</sup> (Fe XVII), during two flights in 1991 and 1993 (Brosius et al. 1996). These line ratios provide density diagnostics between temperatures of $`log(T_e)=5.0`$ and $`log(T_e)=6.7`$ (i.e. $`T_e0.15.0`$ MK). Brosius et al. (1996, see their Fig.8 and 9) derived a differential emission measure curve $`dEM(T)/dT`$ in the temperature range of $`log(T_e)=4.87.0`$, which is reproduced in Fig.2 (top panel), for two observations of active regions (AR93, AR91) and two observations of Quiet Sun regions (QR93, QR91).
We consider now the instrumental response functions of Yohkoh SXT. For active regions, the two filters sensitive to the lowest temperatures are the thin alumninium (Al 1265 Å) and the Al/Mg/Mn composite filter (Tsuneta et al. 1991). The corresponding response functions $`R_1(T)`$ and $`R_2(T)`$ are shown in Fig.2 (second panel), and their filter ratio $`Q(T)=R_2(T)/R_1(T)`$ is given in Fig.2 (bottom panel). In order to understand the temperature contributions to the observed flux we show the differential soft X-ray flux $`dF(T)/dT=[dEM(T)/dT]R(T)`$ (Fig.1, third panel), for both filters and for all 4 regions. The differential soft X-ray flux exhibits a peak at a temperature of $`T_e10^{6.65}=4.5`$ MK for the active regions, and at $`T_e10^{6.3}=2.0`$ MK for the quiet Sun regions.
We calculate now the fluxes $`F_1(h)`$ and $`F_2(h)`$ in the two filters as function of the height $`h`$ above the limb, using the hydrostatic distribution defined in Eqs.5-6, where each fluxtube with (different) temperature $`T`$ has a (different) density scale height of $`\lambda =\lambda _0T`$, while the total ensemble of fluxtubes is summed up by an integration over the entire temperature range (i.e. temperature integral in Eqs.5-6). The resulting SXR fluxes as function of height are shown in Fig.3 top, illustrating that the SXR flux drops exponentially with height. We derive now the filter ratio $`Q(h)=F_2(h)/F_1(h)`$, shown in Fig.3 (middle panel) for all 4 regions. The filter ratio $`Q(h)`$ clearly varies as function of height $`h`$, although each fluxtube is assumed to be isothermal.
We demonstrate now what effect this filter ratio variation $`Q(h)`$ has on the inference of a single-temperature model $`T(h)`$, as it is assumed in the classical filter-ratio method by definition. To invert the filter ratio $`Q(T)`$ as function of the temperature $`T`$, we find the following analyical approximation (accurate within $`\stackrel{<}{}0.7\%`$) in the temperature range of $`T=1.56.0`$ MK (see fit in Fig.2 bottom),
$$Q(T):=\frac{R_2(T)}{R_1(T)}0.39+0.27[log(T)6.18]^{1/2}.$$
(7)
This analytical approximation allows us conveniently to invert the filter-ratio temperature in the range of $`Q=0.40.6`$, i.e.
$$log(T[Q])=6.18+\left(\frac{Q0.39}{0.27}\right)^2.$$
(8)
The inverted temperatures $`T[Q(h)]`$ are shown in Fig.3 bottom for all 4 regions. The filter ratio temperature $`T(h)`$ shows a height dependence from $`T(h=0)2.1`$ MK to $`T(h=0.5r_{\mathrm{}})3.1`$ MK for the quiet regions, and from $`T(h=0)4.14.4`$ MK to $`T(h=0.5r_{\mathrm{}})5.46.3`$ MK for the active regions. Thus, the weighting effect of temperature scale heights over the broadband response function introduces an apparent temperature gradient of $`dT/dh0.003`$ K m<sup>-1</sup> for the quiet corona regions, and about $`dT/dh0.005`$ K m<sup>-1</sup> for active regions. This corresponds about to a doubling of the apparent temperature over a distance of a solar radius $`r_{\mathrm{}}`$,
$$\mathrm{\Delta }T^{SXT}T_0(\frac{h}{r_{\mathrm{}}}).$$
(9)
## 4 DISCUSSION AND CONCLUSIONS
We have investigated the effect of hydrostatic density scale heights in coronal loops on the inference of a filter-ratio temperature from a broadband instrument, in particular for the two thinnest filters of Yohkoh/SXT, which are generally used to derive electron temperatures in active regions and in the quiet corona. The principal effect is that, with increasing altitude $`h`$ (above the solar surface), the emission measure-weighted temperature $`T_e`$ becomes systematically more weighted by the larger scale heights $`\lambda `$ (Fig.1), which are associated with loops of higher temperature, and thus mimic an average temperature increase with height. We used differential emission measure distributions $`dEM(T)/dT`$ that have been observed in active regions and in quiet Sun regions and simulated the temperature bias on $`T(h)`$ for the instrumental response functions of Yohkoh/SXT. The resulting temperature bias can be quantified approximately as $`\mathrm{\Delta }T^{SXT}T_0(h/r_{\mathrm{}})`$. We discuss now the consequences of this result.
The radial variation of temperature in the inner corona (out to 0.7 and 0.95 solar radii) has been examined for the diffuse corona from long-exposure Yohkoh/SXT images by Wheatland, Sturrock, & Acton (1997). These authors find a systematic temperature increase from $`T_e1.6`$ MK near the solar surface to $`T_e2.4`$ at a height of 0.5 solar radii for the 7-9 May 1992 active region, and from $`T_e1.8`$ MK to $`T_e23`$ MK at 1 solar radius for the 26 August 1992 region. This systematic temperature increase of the solar corona was interpreted in terms of a downward heat flux, leading to the conclusion of a heat deposition above the observed height. According to our model (Eq.9), we estimate fully consistent temperature increases \[$`T_e(h=0)=1.6`$ MK $`T_e(h=0.5r_{\mathrm{}})=2.4`$ MK for the first case, and ($`T_e(h=0)=1.8`$ MK $`T_e(h=r_{\mathrm{}})=2.7`$ MK for the second case\] from the emission measure-weighted hydrostatic scale heights alone, even if all fluxtubes are isothermal. Therefore, if the hydrostatic weighting effect on the Yohkoh/SXT filter ratio method would be corrected, no net temperature increase would result, and thus no support for a heating function in the upper corona is warranted.
With the same measurement technique, Priest et al. (1999; 2000) analyzed large-scale arcades and loops and found a temperature increase from $`T_e(h=0)=1.6`$ MK to $`T_e(h=0.5r_{\mathrm{}})=2.22.3`$ MK for a first loop observed on 1992 Oct 3, an increase from $`T_e(h=0)=1.6`$ MK to $`T_e(h=500`$ Mm$`)=2.42.6`$ MK in a second loop, and an increase from $`T_e(h=0)=1.6`$ MK to $`T_e(h=350`$ Mm$`)2.1`$ MK in a third loop. The authors fitted three heating models to these temperature increases $`T_e(h)`$ and found that a uniform heating function provides the best fit for all 3 cases, while heating functions localized at the loop top was found to be less likely, and a heating function localized near the loop footpoints was rejected. From our model (Eq.9) we can reproduce the same temperature increases for these 3 cases, so that virtually no net temperature increase remains, if the Yohkoh/SXT filter ratios would be corrected for the hydrostatic emission measure weighting. Comparing the corrected temperature profiles $`T_e(h)const`$ with the heating models shown of Figs.8 and 9 in Priest et al. (2000), one would conclude that the data are most consistent with the theoretical model of footpoint heating, a conclusion that would also be more in line with other recent observations from TRACE (Schrijver et al. 1999; Aschwanden et al. 2000b).
In summary, we like to point out that filter ratio temperatures from broadband instruments may lead to systematic errors in the determination of vertical temperature profiles $`T_e(h)`$, that can only be corrected properly by forward-fitting of models which contain both temperature $`T_e(h)`$ and density profiles $`n_e(h)`$. The systematic effects are larger for broadband filter ratios (e.g. Yohkoh/SXT) than for narrowband filters (e.g. SoHO/EIT or TRACE). Any detected temperature increase derived from an emission measure-weighted temperature definition is subject to the hydrostatic weighting of a multi-temperature plasma, and does not directly describe a variation ($`dT/dh`$) of the electron temperature along a magnetic field line.
Acknowledgements: We appreciate helpful comments from Drs. David Alexander, Jim Lemen, Karel Schrijver, Greg Slater, and the referee. This work was supported by Yohkoh/SXT NASA contract NAS8-40801.
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## Figure Captions
Fig.1: This cartoon illustrates scale height-weighted contributions of hydrostatic loops or open fluxtubes to the emission measure observed along two line-of-sights above the solar limb. The left line-of-sight at a height of $`h=100`$ Mm above the limb samples significant emission from the 3 loops with temperatures of 1.5-2.5 MK. The right line-of-sight at a height of $`h=200`$ Mm above the limb samples significant emission only from the hottest loop with $`T=2.5`$ MK.
Fig.2: The differential emission measure distribution $`dEM(T)/dT`$ of two active regions (AR93, AR91) and two quiet Sun regions (QR93, QR91) measured by Brosisus et al. (1996) with SERTS data (top panel). The Yohkoh/SXT response function for the two thinnest filters (second panel). The differential SXR fluxes $`dF(T)/dT=[dEM(T)/dT]R(T)`$ for the two SXT filters (thin and thick linestyles) for all 4 regions (third panel). The filter ratio $`Q(T)=R_2(T)/R_1(T)`$ for the two Yohkoh/SXT filters and an analytical approximation in the range of $`T=1.56.0`$ MK (bottom panel).
Fig.3:The height dependence of the observed SXT fluxes $`F(h)`$ for the two filters (thin and thick linestyles) and all 4 regions (different linestyles) (top panel). The resulting filter ratio $`Q(h)`$ for all 4 regions (second panel), and the inferred filter-ratio temperatures (bottom panel). Note that the filter-ratio temperature $`T(h)`$ shows a systematic increase with height, although a model with isothermal loops was assumed.
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# RELAXING NEAR THE CRITICAL POINT
## I Introduction and Motivation
The program of relativistic heavy ion collisions both at Brookhaven and at Cern seeks to understand the phase diagram of QCD in conditions of temperatures that were achieved during the first 10$`\mu `$s after the Big Bang or densities several times that of nuclear matter which could exist at the center of neutron stars. Current theoretical understanding leads to the conclusion that QCD could undergo two phase transitions: a confinement-deconfinement (or hadronization) and the chiral phase transition. Current lattice data seem to suggest that both occur at about the same temperature $`T_c160\text{MeV}`$. The consensus emerging in the field is that several types of observables will have to be studied simultaneously and event-by-event analysis of data will have to be carried out to extract unambiguous signals both hadronic and electromagnetic to reveal the presence of a Quark-Gluon Plasma phase. Recent results reported from CERN-SPS seem to indicate a strong evidence for the existence of the QGP in Pb-Pb collisions, and RHIC at Brookhaven will begin operation soon with Au-Au collisions with four dedicated detectors capable of event-by-event analysis of hadronic and electromagnetic observables.
For QCD with only two flavors of massless quarks (u,d) it has been argued that the chiral phase transition at finite temperature but vanishing baryon number density is of second order and described by the universality class of O(4) Heisenberg ferromagnets. It has also been suggested recently that at finite baryon density there is a second order critical point described by the Ising universality class. Second order critical points are characterized by strong critical long-wavelength fluctuations and a diverging correlation length that could lead to important experimental signatures. These signatures would be akin to critical opalescence near the critical point in binary fluids and could be observed in an event-by-event analysis of the fluctuations of the charged particle transverse momentum distribution (mainly pions). These fluctuations are characterized by the typical correlation length of the order parameter and it has been suggested that the phenomenon of critical slowing down, ubiquitous near the critical point of second order phase transitions, can lead to strong departures from equilibrium that will determine the value of the correlation length when long-wavelength fluctuations freeze out. Critical slowing down of long-wavelength fluctuations near a second order critical point is the statement that the long-wavelength Fourier components of the order parameter relax very slowly towards equilibrium. In mean-field theory in classical critical phenomena, the relaxation time diverges proportional to the susceptibility near the critical temperature but thermal fluctuations renormalize the relaxation time to be of the form $`\tau (\stackrel{}{k}=0)\xi ^z`$ with $`\xi `$ the correlation length, or at critical point $`\tau (k)k^z`$ with $`z`$ a dynamical critical exponent . Another similar manifestation of an anomalously slow relaxation of long-wavelength fluctuations arises in weakly first order phase transitions when the system enters into the mixed phase where the (isothermal) speed of sound, which determines the velocity of propagation of long-wavelength pressure waves, becomes anomalously small resulting in the softest point of the equation of state. In this case there have also been suggestions that there are experimental consequences of this softening in relativistic heavy ion collisions in observables related to collective flow and the transverse momentum distributions of particles at freezeout.
In classical normal fluids near the critical point the vanishing of the (isothermal) speed of sound, critical opalescence (strong scattering of light by long-wavelength fluctuations) and critical slowing down are all related, and in ferromagnets the spin diffusion constant vanishes near the critical point again signaling critical slowing down.
The softening of the equation of state near the critical point of QCD could also have important cosmological implications. When the the QGP enters the mixed phase with hadrons, the speed of sound becomes anomalously small and the time scale for propagation of pressure waves over a given critical wavelength becomes longer than the free fall time for gravitational collapse which is then unhindered by the pressure of the hadronic gas. This could lead to the formation of primordial black holes with a possible imprint in the acoustic peaks in the cosmic microwave background. Other possible cosmological relics from the QCD phase transition with a mixed phase had been predicted, from strange quark nuggets to MACHO’s.
A familiar argument is typically invoked to state that while the QCD phase transition in the Early Universe occurred in local thermodynamic equilibrium (LTE) this may not be the case in Relativistic Heavy Ion Collisions. The argument compares the typical collisional relaxation time scale obtained from a strong interaction process $`\tau _{coll}10^{22}\text{secs}`$ to the time scale for cooling near the critical temperature $`160\text{MeV}`$, i.e. $`T/\dot{T}H^110^5\text{secs}`$. The argument is that since $`\tau _{coll}H^1`$ the phase transition occurs in LTE in cosmology, whereas in relativistic heavy ion collisions at RHIC and LHC energies these time scales will be comparable. However this argument completely neglects the possibility that long-wavelength fluctuations could undergo critical slowing down and freeze out, i.e. fall out of local thermal equilibrium, even before the phase transition. The freeze-out of long-wavelength fluctuations during the phase transition could result in important non-equilibrium effects on the size and distribution of primordial black holes or any other cosmological relic just as they could lead to important observables in the momentum distributions of charged pions in relativistic heavy ion collisions.
Indeed there are simpler experimental situations where this is the case, in typical normal fluids a collisional relaxation time (away from the critical point) is of the order of $`10^9\text{secs}`$ while near the critical point (even at $`10\%`$ of the critical temperature) critical slowing down becomes very dramatic and thermalization time scales become of the order of minutes if not hours.
Thus the phenomenological importance of critical slowing down for the QCD phase transition both in Relativistic Heavy Ion Collisions as well as in Early Universe Cosmology motivates us to study this phenomenon in a model Quantum Field Theory that bears on the low energy (chiral) phenomenology of QCD, the $`O(N)`$ linear sigma model. Furthermore the study of critical slowing down is the precursor to a more complete program to understand transport phenomena and the relaxation of hydrodynamic modes at or near a critical point.
Goal: Our goal is to provide a consistent microscopic description of critical slowing down at or near criticality directly from an underlying quantum field theory that is at least phenomenologically motivated to study the QCD phase transitions. This will be a first step in a program that seeks to offer a consistent description of transport near critical points that eventually may be merged with a hydrodynamic description to obtain a more reliable picture of critical phenomena near the deconfinement and chiral phase transitions and an assessment of the potential phenomenological observables both in early universe cosmology as well as in relativistic heavy ion collisions. We begin this program in this article by focusing on the relaxation rate of long-wavelength fluctuations of the order parameter at and near the critical point in a consistent non-perturbative framework.
Strategy: We begin our study of critical slowing down by analyzing the relaxation rate of long-wavelength fluctuations of the order parameter at and near the critical point in an $`O(N)`$ scalar field theory, which is a phenomenological arena to study the relaxation of sigma mesons and pions. Our first step is to obtain the relaxation rate to lowest order in perturbation theory (two loops). This calculation reveals clearly the breakdown of the perturbative expansion for long-wavelength fluctuations at or near the critical point as a result of the strong infrared behavior for soft loop momentum and the necessity for a non-perturbative treatment. We then implement a non-perturbative resummation of bubble-type diagrams via the large $`N`$ approximation to obtain the damping rate in the next-to-leading order in the large $`N`$ limit. The resummation implied by the large $`N`$ limit to order $`1/N`$ is akin to that obtained via the renormalization group with the one loop beta function and reveals the softening of the scattering amplitude and the crossover to an effective three dimensional theory for momenta $`q\lambda T`$ with $`\lambda `$ the quartic coupling.
Summary of main results:
We have obtained the relaxation rate for long-wavelength fluctuations of the order parameter at the critical point and for homogeneous fluctuations near criticality both to lowest order in perturbation theory (two loops) and near the critical point to next to leading order in the large $`N`$ limit.
The two-loop results for the relaxation rate for a fluctuation of wavevector $`\stackrel{}{k}`$ of the order parameter at the critical point is found to be $`\mathrm{\Gamma }(k,T)\lambda ^2T^2/k`$ whereas near the critical point, homogeneous fluctuations (with $`\stackrel{}{k}=0`$) relax with a rate $`\mathrm{\Gamma }_0(m_T,T)\lambda ^2T^2/m_T`$. Here $`m_T|TT_c|^{1/2}T_c`$ is the effective thermal mass. These results clearly reveal the breakdown of the perturbative expansion in the long wavelength limit $`k0`$ at $`T=T_c`$ and for $`TT_c`$ and $`k=0`$.
A detailed analysis of the different contributions to these results for the relaxation rate shows that the rate is dominated by very soft loop momentum $`q\lambda T`$ which in the weak coupling limit $`\lambda 1`$ are classical. The implementation of a non-perturbative resummation via the large $`N`$ limit explicitly leads to an effective scattering amplitude that vanishes in the static long-wavelength limit as a consequence of the crossover to a three dimensional theory for loop momenta $`q\lambda T`$. This effective scattering amplitude allows us to recognize that the effective three dimensional coupling for soft momenta approaches the three dimensional non-trivial (Wilson-Fisher) fixed point in the long-wavelength limit near the critical point. The large $`N`$ resummation for the relaxation rate incorporates this effective three dimensional coupling in the spectral density that determines the imaginary part of the retarded self-energy for the order parameter. Since the effective three dimensional coupling is driven to its fixed point at long wavelength, the contribution from very soft loop momenta $`q\lambda T`$ which give the strongest infrared behavior in lowest order in perturbation theory is effectively screened by this renormalization of the coupling. Consequently the most important contribution to the relaxation rate arises both from the semisoft classical region of loop momentum $`Tq\lambda T`$ and also from the hard region $`qT`$. A detailed analysis of the contribution from the loop momenta reveals a non-perturbative ultrasoft scale
$$k_{us}\frac{\lambda T}{4\pi }e^{\frac{4\pi }{\lambda }}.$$
We find that for soft momenta $`kk_{us}`$ the damping rate is dominated by classical semisoft loop momenta and given by
$$\mathrm{\Gamma }(k,T)\stackrel{kk_{us}}{=}\frac{\lambda T}{2\pi N}\left[1+𝒪\left(\frac{1}{\mathrm{ln}\frac{\lambda T}{k}}\right)\right].$$
(1)
For $`k,m_Tk_{us}`$ the classical approximation breaks down and the damping rate at the critical point $`m_T=0`$ for $`kk_{us}`$ is given by
$$\mathrm{\Gamma }(k,T)\stackrel{kk_{us}}{=}\frac{4\pi T}{3N\mathrm{ln}\frac{T}{k}}\left[1+𝒪\left(\frac{1}{\mathrm{ln}\frac{T}{k}}\right)\right].$$
(2)
For homogeneous fluctuations near the critical point ($`k=0,m_T|TT_c|^{1/2}0`$) the damping rate is given by
$$\mathrm{\Gamma }_0(m_T,T)\stackrel{m_Tk_{us}}{=}\frac{4\pi T}{3N\mathrm{ln}\frac{T}{m_T}}\left[1+𝒪\left(\frac{1}{\mathrm{ln}\frac{T}{m_T}}\right)\right].$$
(3)
Thus critical slowing down, i.e, the vanishing of the quasiparticle width $`\mathrm{\Gamma }`$ for long-wavelengths emerges in the ultrasoft limit $`kk_{us}`$ or very near the critical point $`m_Tk_{us}`$ where it vanishes logarithmically slow in the $`k,m_T0`$ limit to this order in $`1/N`$. Notice that in such regimes the rate is independent of the coupling $`\lambda `$.
The large $`N`$ approximation is not limited to weak coupling and our results apply just as well to a strong coupling case $`\lambda 1`$ wherein we find that the relaxation rate is given by eq.(2)-(3). However this analysis clearly reveals that for weak coupling there emerges a hierarchy of widely separated scales for loop momenta: from hard $`qT`$ to semisoft $`Tq\lambda T`$, and soft $`\lambda Tq`$ that lead to different contributions to the relaxation rate. Which is the relevant scale for the damping rate is determined by the wavevector of the fluctuation of the order parameter and the proximity to the critical temperature. For $`k,m_Tk_{us}`$ the classical approximation does apply and the damping rate is dominated by the soft and semisoft classical loop momenta \[with the result (1)\], whereas for $`k,m_Tk_{us}`$ the classical approximation breaks down and the damping rate is dominated by hard loop momenta $`qT`$ \[with the results (2)-(3)\].
A similar hierarchy exists in non-abelian plasmas- and we compare and contrast our results in the scalar theory with those in the hard thermal loop approximation in abelian and non-abelian plasmas-.
This article is organized as follows: in section II we introduce the model, obtain the real-time equation of motion for the order parameter and describe the strategy followed to obtain the relaxation rate. In section III we carry out a perturbative analysis of the relaxation rate to two loops order, recognize the breakdown of perturbation theory and compare to the case of the hard thermal loop resummation program in gauge theories. In sections IV and V we introduce the large $`N`$ limit, obtain the effective static scattering amplitude in leading order in the large $`N`$ and discuss the dimensional crossover for soft momenta and the effective three dimensional coupling being driven to the three dimensional fixed point. We then use these results to obtain the relaxation rate and near criticality to order $`𝒪(1/N)`$ in the large $`N`$ limit and explicitly discuss the screening of the soft loop momenta. The contribution from classical soft and semisoft momenta and that of hard loop momenta are analyzed separately to highlight the important differences. In this section we discuss further the validity of a quasiparticle interpretation of the collective long-wavelength fluctuations of the order parameter. In section VI we summarize our conclusions and results and discuss the next step of the program. In appendix A the equations of motion in the large $`N`$ limit are derived and in appendix B the polarization integral is computed.
## II Preliminaries: the model and the strategy
We study the model of scalar fields $`\stackrel{}{\mathrm{\Phi }}(x)`$ in the vector representation of $`O(N)`$, which is conjectured to describe the equilibrium universality class for the chiral phase transition with two light quarks for $`N=4`$. The Lagrangian density is given by
$$=\frac{1}{2}(_\mu \stackrel{}{\mathrm{\Phi }})^2\frac{1}{2}\left[m_T^2+\delta m^2(T)\right]\stackrel{}{\mathrm{\Phi }}^2(x)\frac{\lambda }{2N}[\stackrel{}{\mathrm{\Phi }}^2(x)]^2+\stackrel{}{J}\stackrel{}{\mathrm{\Phi }}$$
(4)
where the external current $`\stackrel{}{J}`$ has been introduced to generate an expectation value for the scalar field (i.e. the order parameter) by choosing it to be nonzero along a particular (sigma) direction.
The counterterm $`\delta m^2(T)`$ is introduced to cancel the tadpole contributions (local terms) so that perturbation theory (or the large $`N`$ expansion) is carried out in terms of the effective thermal mass $`m_T`$. In particular to leading order in the large $`N`$ expansion there is the hard thermal loop contribution given by the usual tadpole term $`\lambda \stackrel{}{\mathrm{\Phi }}^2/N\lambda T^2`$ which combined with the zero temperature (negative) mass squared leads to an effective finite temperature mass $`m_T^2(T^2T_c^2)`$. The critical theory corresponds to $`T=T_c`$, i.e. $`m_T=0`$. In this case the counterterm $`\delta m^2(T)`$ is adjusted consistently order by order to set the effective finite temperature mass equal to zero.
As stated in the introduction, our goal is to obtain the relaxation rate (damping rate) of the order parameter at and near the critical point. This will be achieved by obtaining the equation of motion for the expectation value of the scalar field, i.e, the order parameter and treating its evolution in real time as an initial value problem. This is achieved by coupling an external source that serves the purpose of preparing the initial state. From the equation of motion we recognize the self-energy and compute the relaxation rate from its imaginary part on shell. We write
$$\mathrm{\Phi }^a(\stackrel{}{x},t)=\phi (\stackrel{}{x},t)\delta ^{a,1}+\eta ^a(\stackrel{}{x},t);\stackrel{}{\eta }(\stackrel{}{x},t)=0$$
(5)
where we chose the particular direction “1” by choosing external source term to be different from zero along this direction to give the field an expectation value (see below). The equation of motion for $`\phi (\stackrel{}{x},t)`$ is obtained by imposing that $`\eta ^i(\stackrel{}{x},t)=0`$ consistently in the perturbative expansion. In terms of the spatial Fourier transform of the order parameter $`\phi `$ and following the steps detailed in appendix A.2 (see also) we find
$$\ddot{\phi }_k(t)+[k^2+m_T^2+\delta m^2(T)+m_{tad}^2(T)]\phi _k(t)+_{\mathrm{}}^{\mathrm{}}\mathrm{\Sigma }_{ret,k}(tt^{})\phi _k(t^{})𝑑t^{}=J_k(t)$$
where $`J_k(t)`$ is the external source that generates the initial value problem and $`\mathrm{\Sigma }_{ret,k}(tt^{})`$ is the two-loops retarded self-energy without the tadpole contributions. The one and two-loops tadpole contributions (local) are accounted for in $`m_{tad}^2(T)`$. As described above, the counterterm $`\delta m^2(T)`$ is fixed consistently in perturbation theory by requesting that it cancels all constant (in space and time) contributions to the self-energy (such as the tadpoles) i.e,
$$\delta m^2(T_c)+m_{tad}^2(T_c)=0$$
The retarded self-energy has a dispersive representation in terms of the spectral density $`\stackrel{~}{\rho }(\omega ,k)`$ given by
$$\mathrm{\Sigma }_{ret,k}(tt^{})=\frac{d\omega ^{}}{2\pi }e^{i\omega ^{}(tt^{})}𝑑\omega \frac{\stackrel{~}{\rho }(\omega ,k)}{\omega \omega ^{}iϵ}$$
(6)
in terms of which the relaxation (damping) rate is given by
$$\mathrm{\Gamma }(k,T)=\frac{\pi }{2\omega _p(k)}\stackrel{~}{\rho }(\omega _p(k),k,T)$$
(7)
where $`\omega _p(k)`$ is the position of the pole in the propagator, i.e, the true dispersion relation. For the perturbative two loops or to leading order in the large $`N`$ limit as studied here $`\omega _p^2(k)=k^2+m_T^2`$, which at $`T=T_c`$ takes the form $`\omega _p(k)=|\stackrel{}{k}|`$.
With the purpose of clearly revealing the breakdown of the perturbative expansion for soft momenta $`k\lambda T`$ we will begin our analysis by focusing first on the perturbative evaluation of the damping rate. At one loop order the only contribution to the self energy is given by the tadpole term $`\lambda \stackrel{}{\mathrm{\Phi }}^2(\stackrel{}{x},t)/N`$ which is local, determines to lowest order the temperature dependent mass $`m_T|TT_c|^{1/2}`$ and determines the counterterm. Furthermore this is the leading contribution in the hard thermal loop limit. The lowest order contribution to the absorptive (imaginary) part of the self-energy arises at two loops and is studied in detail in the next section.
## III Perturbation Theory: two loops
We begin our study by carrying out a perturbative evaluation of the damping rate to lowest order, i.e. to two loops to reveal several important features of the soft momentum limit, and to pave the way to implement a non-perturbative evaluation of the self-energy in the large $`N`$ limit. Furthermore, as it will become clear during the course of the calculation, the lowest order contribution contains some of the important ingredients of the large $`N`$ limit and will highlight the contribution to the relaxation rate from different regions of loop momentum.
After substracting the one and two loops tadpole contributions which are cancelled by the counterterm, the spatial Fourier transform of the retarded self-energy reads:
$`\mathrm{\Sigma }_{ret,k}(tt^{})`$ $`=`$ $`8\lambda ^2{\displaystyle \frac{N+2}{N^2}}{\displaystyle }{\displaystyle \frac{d^3p}{(2\pi )^3}}{\displaystyle \frac{d^3q}{(2\pi )^3}}\{G_{\stackrel{}{k}+\stackrel{}{q}}^>(tt^{})G_{\stackrel{}{p}+\stackrel{}{q}}^>(tt^{})G_\stackrel{}{p}^>(tt^{})`$ (8)
$``$ $`G_{\stackrel{}{k}+\stackrel{}{q}}^<(tt^{})G_{\stackrel{}{p}+\stackrel{}{q}}^<(tt^{})G_\stackrel{}{p}^<(tt^{})\}\mathrm{\Theta }(tt^{})`$ (9)
where the Wightmann functions $`G^>,G^<`$ are given in Appendix A.1.
With the purpose of comparing with the results of later sections, it proves convenient to introduce the intermediate quantities
$`𝒢_q^>(tt^{})`$ $`=`$ $`2\lambda {\displaystyle \frac{N+2}{N}}{\displaystyle \frac{d^3p}{(2\pi )^3}G_{\stackrel{}{p}+\stackrel{}{q}}^>(tt^{})G_\stackrel{}{p}^>(tt^{})}`$ (10)
$`=`$ $`{\displaystyle 𝑑q_0e^{iq_0(tt^{})}S^>(q_0,q)}`$ (11)
$`𝒢_q^<(tt^{})`$ $`=`$ $`{\displaystyle \frac{2\lambda }{N}}(N+2){\displaystyle \frac{d^3p}{(2\pi )^3}G_{\stackrel{}{p}+\stackrel{}{q}}^<(tt^{})G_\stackrel{}{p}^<(tt^{})}`$ (12)
$`=`$ $`{\displaystyle 𝑑q_0e^{iq_0(tt^{})}S^<(q_0,q)}`$ (13)
and using the expression for the Wightmann functions $`G_\stackrel{}{k}^>(tt^{}),G_\stackrel{}{k}^<(tt^{})`$ given in appendix A.1, it is a straightforward exercise to show that the spectral functions $`S^<(q_0,q);S^>(q_0,q)`$ obey the KMS condition
$$S^<(q_0,q)=e^{\beta q_0}S^>(q_0,q)$$
Introducing the spectral density
$$\sigma (q_0,q)=S^>(q_0,q)S^<(q_0,q)$$
which at the critical point $`T=T_c`$, i.e. $`m_T=0`$ is found to be given by
$`\sigma (q_0,q)=2\lambda {\displaystyle \frac{N+2}{N}}{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{4p|\stackrel{}{p}+\stackrel{}{q}|}}`$ $`\{[1+n_{\stackrel{}{q}+\stackrel{}{p}}+n_\stackrel{}{p}][\delta (q_0|\stackrel{}{p}+\stackrel{}{q}|p)\delta (q_0+|\stackrel{}{p}+\stackrel{}{q}|+p)]`$ (15)
$`+[n_\stackrel{}{p}n_{\stackrel{}{q}+\stackrel{}{p}}][\delta (q_0|\stackrel{}{p}+\stackrel{}{q}|+p)\delta (q_0+|\stackrel{}{p}+\stackrel{}{q}|p)]\}`$
we find
$$S^>(q_0,q)=[1+n(q_0)]\sigma (q_0,q),S^<(q_0,q)=n(q_0)\sigma (q_0,q)$$
(16)
where
$$n(q_0)=\frac{1}{e^{\beta q_0}1},n_\stackrel{}{k}=\frac{1}{e^{\beta |\stackrel{}{k}|}1}.$$
(17)
The near critical case is considered in sec. III.C.
A lengthy but straighforward calculation with the Bose-Einstein distribution functions for massless particles leads to the following expression for $`\sigma (q_0,q)`$
$$\sigma (q_0,q)=\frac{\lambda }{8\pi ^2}\frac{N+2}{N}\left\{\mathrm{\Theta }(|q_0|q)\text{sign}(q_0)+\frac{2T}{q}\mathrm{ln}\left[\frac{1e^{\frac{|q_0+q|}{2T}}}{1e^{\frac{|q_0q|}{2T}}}\right]\right\}.$$
(18)
It is important to emphasize that the second, finite temperature term is the combination of two different contributions: a two massless particle cut with support in the region $`q_0>q`$ and a Landau damping cut with support in the region $`qq_0q`$.
Introducing the Fourier representation of the theta function
$$\mathrm{\Theta }(tt^{})=\frac{d\omega }{2\pi i}\frac{e^{i\omega (tt^{})}}{\omega +iϵ}$$
(19)
in the expression for the self-energy (9), we find the spectral density that enters in the dispersive representation of the retarded self-energy (6) to be given by
$`\stackrel{~}{\rho }(\omega ,k)`$ $`=`$ $`{\displaystyle \frac{4\lambda }{N}}{\displaystyle \frac{d^3q}{(2\pi )^3}\frac{dq_0}{2|\stackrel{}{k}+\stackrel{}{q}|}\sigma (q_0,q)\left[1+n_{\stackrel{}{q}+\stackrel{}{k}}+n(q_0)\right]\left[\delta (\omega q_0|\stackrel{}{k}+\stackrel{}{q}|)\delta (\omega +q_0+|\stackrel{}{k}+\stackrel{}{q}|)\right]}`$ (20)
$`=`$ $`{\displaystyle \frac{4\lambda }{N}}{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}{\displaystyle \frac{1}{2|\stackrel{}{k}+\stackrel{}{q}|}}\{\sigma (q,\omega |\stackrel{}{k}+\stackrel{}{q}|)[n_{\stackrel{}{q}+\stackrel{}{k}}n(|\stackrel{}{k}+\stackrel{}{q}|\omega )]`$ (21)
$`+`$ $`\sigma (q,\omega +|\stackrel{}{k}+\stackrel{}{q}|)[n_{\stackrel{}{q}+\stackrel{}{k}}n(\omega +|\stackrel{}{k}+\stackrel{}{q}|)]\}`$ (22)
with $`\sigma (q_0,q)`$ given by (18) and we used that,
$$\sigma (q_0,q)=\sigma (q_0,q),1+n(q_0)=n(q_0).$$
Although this seems to be a cumbersome manner to write down the two loop contribution, it does prove convenient to establish contact with the large $`N`$ description in the next section.
We are interested in the relaxation of long-wavelength fluctuations of the order parameter, hence we will consider the soft limit $`kT`$.
We study the contributions coming from soft $`qT`$ and hard $`qT`$ loop-momenta separately.
### A The soft momenta contribution (classical region)
This is the classical region where the Bose-Einstein distribution functions can be approximated as $`n(\omega )T/\omega ,n_\stackrel{}{k}T/k`$.
In this regime the contribution of the soft momenta $`qT`$ yields
$$\stackrel{~}{\rho }_{cl}(\omega ,k)=\frac{\lambda ^2T^2\omega }{2\pi ^2}\frac{N+2}{N^2}\frac{d^3q}{(2\pi )^3}\frac{1}{q|\stackrel{}{k}+\stackrel{}{q}|^2}\left[\frac{\mathrm{ln}\left|\frac{\omega |\stackrel{}{k}+\stackrel{}{q}|+q}{\omega |\stackrel{}{k}+\stackrel{}{q}|q}\right|}{\omega |\stackrel{}{k}+\stackrel{}{q}|}+\frac{\mathrm{ln}\left|\frac{\omega +|\stackrel{}{k}+\stackrel{}{q}|+q}{\omega +|\stackrel{}{k}+\stackrel{}{q}|q}\right|}{\omega +|\stackrel{}{k}+\stackrel{}{q}|}\right]$$
where we only kept the contribution of order $`T`$ to $`\sigma (q_0,q)`$. We evaluate the spectral density at $`\omega =k`$ leading to the following form for the soft momenta contribution to the damping rate (7)
$$\mathrm{\Gamma }_{cl}(k,T)=\left(\frac{\lambda T}{4}\right)^2\frac{N+2}{2\pi ^4N^2}\left[J_1(k,T)+J_2(k,T)\right]$$
(23)
with $`J_1(k,T)`$ and $`J_2(k,T)`$ given by the following expressions
$`J_1(k,T)`$ $`=`$ $`{\displaystyle d^3q\frac{\mathrm{ln}\left|\frac{|\stackrel{}{k}+\stackrel{}{q}|+qk}{|\stackrel{}{k}+\stackrel{}{q}|qk}\right|}{q|\stackrel{}{k}+\stackrel{}{q}|^2(|\stackrel{}{k}+\stackrel{}{q}|k)}}`$ (24)
$`J_2(k,T)`$ $`=`$ $`{\displaystyle d^3q\frac{\mathrm{ln}\left|\frac{|\stackrel{}{k}+\stackrel{}{q}|+q+k}{|\stackrel{}{k}+\stackrel{}{q}|q+k}\right|}{q|\stackrel{}{k}+\stackrel{}{q}|^2(|\stackrel{}{k}+\stackrel{}{q}|+k)}}.`$ (26)
The angular integrals can be performed analytically and $`k`$ can be scaled out of the integral by introducing the variable $`x=q/k`$ leading to the final expression
$$\mathrm{\Gamma }_{cl}(k,T)=\frac{\lambda ^2T^2}{16\pi ^3k}\frac{N+2}{N^2}_0^{\mathrm{}}𝑑xF[x]$$
(27)
with the function $`F[x]`$ given by
$$F[x]=\frac{2}{x}\left[2x\frac{\mathrm{ln}x}{x^21}+\mathrm{ln}\left|\frac{x+1}{x1}\right|\right].$$
and is depicted in Fig. 1. In principle, the upper limit in the integral in (27) should be $`\alpha T/k`$ with $`\alpha 1`$ to restrict the integral to the soft momenta where the classical approximation is valid, but the integrand falls off as $`1/x^2`$ for $`x1`$ and the integral is dominated by the small $`x`$ region $`0<x1`$ as shown in fig. 1. The integral from $`x=0`$ to $`x=\mathrm{}`$ in eq.(27) gives $`2\pi ^2`$ and the contribution from the classical loop momenta $`qT`$ to the damping rate to two loops order is thus given by
$$\mathrm{\Gamma }_{cl}(k,T)=\frac{\lambda ^2T^2}{8\pi k}\frac{N+2}{N^2}.$$
(28)
As it will be seen in detail in the next section, the large $`N`$ resummation leads to a screening of the soft loop momentum which cuts off the contribution of momentum $`q<\lambda T`$. Hence, with the purpose of comparing the perturbative two-loop result with that in the large $`N`$ limit, it proves convenient to obtain the contribution to the damping rate from the region of semisoft loop momenta with $`\lambda Tq\alpha T`$ with $`\lambda \alpha 1`$ for the case of soft external momentum $`\lambda T/k1`$.
The contribution to the integral of the function $`F[x]`$ from this region is given by
$$_{\lambda T/k}^{\alpha T/k}𝑑xF[x]\frac{4k}{\lambda T}\mathrm{ln}\frac{\lambda T}{k}$$
therefore this region of loop momentum gives a contribution to the damping rate given by
$$\mathrm{\Gamma }_{Tq>\lambda T}\frac{\lambda T}{2\pi }\frac{N+2}{N^2}\mathrm{ln}\frac{\lambda T}{k}.$$
(29)
Thus we see that the contribution from the soft region of internal loop momentum $`q<\lambda T`$ contributes a factor $`\lambda T/k1`$ larger than the region of momentum $`q>\lambda T`$ for soft external momentum $`k\lambda T`$. This observation will become important when we compare with the large $`N`$ result, because as we will show explicitly below, the resummation of the effective scattering amplitude will lead to a softening of the effective vertex and hence screening of very soft momenta in the loop.
### B The hard momenta contribution
We now focus on obtaining the contribution to the damping rate from hard loop momenta $`qT`$. From the expression of the finite temperature contribution to the spectral density $`\sigma (q_0,q)`$ (18) it is clear that hard momenta $`q_0,qT`$ will be exponentially suppressed unless either $`|qq_0|T`$ or $`|q+q_0|T`$. Consider the expression for the spectral density (22) for $`\omega =k`$ and consider the contribution from the delta function with support for $`q_0=k|\stackrel{}{k}+\stackrel{}{q}|`$; it is straightforward to see that the other delta function will give a similar contribution. For $`qTk`$ we find that $`|q_0+q|=k|1\mathrm{cos}\theta |`$ where $`\theta `$ is the angle between $`\stackrel{}{k}`$ and $`\stackrel{}{q}`$ and $`|q_0q|2T`$, hence the region of loop momentum that dominates corresponds to the emission (or absorption) of a pair of scalars (the particles in the loop) with total center of mass momentum collinear with the external momentum.
Keeping the leading term $`\mathrm{ln}(|q_0+q|/2T)`$ in $`\sigma (q_0,q)`$ and the full occupation factors in the expression for $`\stackrel{~}{\rho }`$ the spectral density (22) becomes
$$\stackrel{~}{\rho }(\omega =k,k)\lambda ^2Tk\mathrm{ln}\frac{T}{k}$$
and the contribution to the damping rate from the hard loop momentum region is
$$\mathrm{\Gamma }_{hard}(k,T)\lambda ^2T\mathrm{ln}\frac{T}{k}$$
which is a factor $`(k/T)\mathrm{ln}\frac{T}{k}1`$ smaller than the contribution from the classical loop momenta (28) for $`k/T1`$. In summary, when the external momentum is soft $`k\lambda T`$ the damping rate is completely determined by the classical region $`qT`$ of loop momenta and given by eq.(28). That in a scalar theory the leading temperature effects are determined by the classical region of loop momenta was already anticipated in references but the computation above identifies the contribution to the relaxation rate from several different regions of loop momentum. This identification will become important to understand the result obtained from the large $`N`$ limit.
### C Near criticality
A calculation very similar to that in the critical inhomogeneous case can be carried out for homogeneous fluctuations ($`k=0`$) near the critical point with the effective thermal mass $`m_T^2\lambda (T^2T_c^2)`$ straightforwardly. In this case the angular integrals are trivial and most of the steps are similar to the critical case leading to (see also)
$$\mathrm{\Gamma }_0(m_T,T)=\frac{\lambda ^2T^2}{8\pi m_T}\frac{N+2}{N^2}$$
A similar analysis for contributions from different regions of loop momentum is obtained by replacing $`km_T`$ in the arguments above.
The resonance parameter $`\mathrm{\Gamma }(k,T)/\omega _p`$ with $`\omega _p`$ the position of the single (quasi) particle pole determines how broad is the resonance. If $`\mathrm{\Gamma }(k,T)\omega _p`$ the quasiparticle can be described by a narrow resonance and its decay occurs on time scales much longer than those of the microscopic oscillations $`\omega _p^1`$. On the other hand for $`\mathrm{\Gamma }(k,T)\omega _p`$ the notion of quasiparticle is not appropriate and the excitation is described by a very short lived broad resonance.
The two loops calculation reveals that at the critical point $`m_T=0`$
$$\frac{\mathrm{\Gamma }(k,T)}{k}\stackrel{Tk}{=}\frac{N+2}{8\pi N^2}\left(\frac{\lambda T}{k}\right)^2;$$
analogously, near criticality for homogeneous fluctuations $`k=0`$
$$\frac{\mathrm{\Gamma }(m_T,T)}{m_T}\stackrel{Tm_T}{=}\frac{N+2}{8\pi N^2}\left(\frac{\lambda T}{m_T}\right)^2.$$
Hence up to this order a quasiparticle interpretation is not reliable for $`k,m_T\lambda T`$.
Moreover, for very soft external momenta or very near the critical temperature, $`k;m_T\lambda T`$ the perturbative expansion clearly breaks down and a non-perturbative scheme must be invoked to study the damping rate.
This situation is similar to that in the hard thermal loops (HTL) program in the sense that for external momenta $`k\lambda T`$ (in gauge theories $`\lambda `$ must be replaced by the gauge coupling squared) a non-perturbative resummation is needed. However, here the similarity ends and the major difference with the HTL program is revealed: whereas in the HTL case the non-perturbative region is dominated by hard internal loop momenta $`qT`$, the relaxational dynamics at the critical point is dominated by soft classical (and as it will become clear below, also semi-soft) internal loop momenta $`qT`$. The difference can also be clearly seen formally by restoring the $`\mathrm{}`$ in the contributions: the temperature always appears in the combination $`T/\mathrm{}`$ (from the distribution functions), in the HTL program the gauge coupling constant squared $`e^2e^2\mathrm{}`$ (since this is the loop counting parameter) hence the HTL scale $`e^2T^2e^2T^2/\mathrm{}`$. However in the scalar case the loop counting parameter is $`\lambda \lambda \mathrm{}`$, hence the contribution $`\lambda ^2T^2`$ is classical i.e. independent of $`\mathrm{}`$. Therefore whereas in the HTL program perturbation theory breaks down at a semiclassical scale $`keT/\sqrt{\mathrm{}}`$, at the critical point of a scalar field theory the perturbative expansion breaks down at a classical scale $`k\lambda T`$. In the HTL program the damping rate of collective excitations is typically of order $`e^2T`$ and the quasiparticle poles (plasmons and plasminos) are of order $`\omega _peT/\sqrt{\mathrm{}}`$ hence for weak coupling the long-wavelength quasiparticles are always relatively narrow resonances. This is in striking contrast with the case of a critical scalar theory where the long-wavelength excitation of the order parameter is gapless.
## IV Large N
Having recognized the non-perturbative nature of the relaxation for the long-wavelength components of the order parameter, we seek to use a consistent non-perturbative description and study the relaxation of the order parameter in the large $`N`$ limit. This limit is best studied by introducing an auxiliary field that replaces the quartic interaction via a Gaussian integration (Hubbard-Stratonovich transformation), hence the Lagrangian density becomes
$$=\frac{1}{2}(_\mu \stackrel{}{\mathrm{\Phi }})^2\left[\sqrt{\frac{\lambda }{N}}\chi (x)+\frac{1}{2}\left(m_T^2+\delta m^2(T)\right)\right]\stackrel{}{\mathrm{\Phi }}^2(x)+\frac{1}{2}\chi ^2(x)+\stackrel{}{J}(x)\stackrel{}{\mathrm{\Phi }}(x).$$
(30)
Before we engage in a study of the damping rate, it is important to highlight that the large $`N`$ expansion effectively provides a reorganization of the perturbative series which for example at leading order and at zero temperature is akin to the resummation of the leading logarithms via the renormalization group for the scattering amplitude. We now study in detail this resummation at finite temperature which will reveal the screening of the scattering amplitude for soft momenta which in turn will be responsible for screening the infrared behavior of the spectral functions and the damping rate.
In this section we will focus on the critical theory $`m_T=0`$. The analysis of the off-critical case is given in section V.
### A Effective Scattering Amplitude
To leading order in the large $`N`$ the two particle to two particle scattering amplitude is dominated by $`s`$-channel exchange and is completely determined by the propagator of the auxiliary field $`\chi `$. Fig. 2 shows the Dyson sum for the propagator of the auxiliary field in leading order in the large $`N`$ limit and fig. 3 shows the $`s`$-channel scattering amplitude in leading order, the $`t`$ and $`u`$-channel contributions are subleading. The bubble diagram which is the building block of the propagator of the auxiliary field (and therefore the $`s`$-channel scattering amplitude) is simpler to be calculated in the Matsubara formulation of finite temperature field theory with an external frequency $`\nu _n=2\pi nT`$ and given by
$$I_{bub}(\nu _n,q)=2\lambda T\underset{\nu _m}{}\frac{d^3p}{(2\pi )^3}\frac{1}{\nu _m^2+\stackrel{}{p}^2}\frac{1}{(\nu _m+\nu _n)^2+(\stackrel{}{p}+\stackrel{}{q})^2}.$$
(31)
To illustrate the resummation in a more clear manner, we focus on the static limit which is obtained by setting the external Matsubara frequency to zero. The strongest infrared behavior and leading contribution in the high temperature limit arises from the term $`m=0`$ in the Matsubara sum, the remaining spatial momentum integral is carried out leading to
$$I_{bub}(0,\stackrel{}{q})=\frac{\lambda T}{4q}$$
(32)
and the $`s`$-channel scattering amplitude in the static limit is given by (see Fig. 3)
$$M(0,\stackrel{}{k}+\frac{\stackrel{}{q}}{2};0,\stackrel{}{k}+\frac{\stackrel{}{q}}{2})=\frac{1}{N}\frac{\lambda }{1+I_{bub}(0,\stackrel{}{q})}.$$
We see that the effective temperature and momentum dependent coupling constant, defined as the coefficient of $`1/N`$ in the s-channel scattering amplitude in the static limit, is given by
$$\lambda _{eff}(q)=\frac{\lambda }{1+\frac{\lambda T}{4q}}.$$
(33)
This expression reveals several noteworthy features. Firstly, we see that at high temperature in the critical region the actual expansion parameter in the sum of bubbles is $`\lambda T/q`$ with $`q`$ the spatial momentum tranferred into the loop. The factor $`T`$ is a consequence of the dimensional reduction and the factor $`\lambda T`$ can be interpreted as the dimensionful three dimensional coupling. Since the expansion is in terms of dimensionless quantities the factor $`q`$ in the denominator is required for dimensional reasons. In fact this can be understood via a parallel with the calculation at zero temperature in $`4ϵ`$ space-time euclidean dimensions with a coupling $`\lambda T^ϵ`$ with T now some dimensionful scale, the loop integral for the massless theory produces a factor $`q^ϵ`$ and for $`ϵ=1`$ i.e. the three dimensional theory one finds the result for the finite temperature loop in the static limit. Secondly, the expression for the effective coupling (33) is a result of the large $`N`$ resummation to leading order and is the same as that obtained from the solution to the renormalization group equation for the running coupling using the one-loop beta function obtained in the $`ϵ`$ expansion and setting $`ϵ=1`$, i.e. the large $`N`$ resummation is akin to the resummation obtained from the renormalization group in euclidean field theory, in the sense that the leading order in the large $`N`$ leads to a running coupling which is the same as that obtained from the one-loop beta function. Thirdly, since the effective expansion parameter in the sum of bubbles is $`\lambda T/4q`$ it is convenient to introduce the three dimensional coupling $`\lambda _3(q)=\lambda T/4q`$ and its effective counterpart
$$\lambda _{3,eff}(q)=\frac{\lambda _3(q)}{1+\lambda _3(q)}=\frac{\lambda T}{4q+\lambda T}.$$
(34)
The main point is that this effective three-dimensional coupling is driven to the three-dimensional Wilson-Fisher fixed point $`\lambda ^{}=1`$ in the soft momentum limit $`q0`$, while the effective four dimensional coupling (33) is driven to the trivial fixed point in this limit. Hence, whereas the three dimensional coupling $`\lambda _3(q)=\lambda T/4q`$ diverges in the $`q0`$ limit, the large $`N`$ (equivalent to the renormalization group) effective coupling $`\lambda _{3,eff}(q)`$ is driven to a finite fixed point in the soft momentum limit. Therefore the large $`N`$ resummation is effectively screening the infrared divergences associated with the soft momentum limit much in the same manner as the resummation implied by the renormalization group within the $`ϵ`$ expansion. Obviously, in exactly three euclidean dimensions one could hardly justify the validity of an $`ϵ`$ expansion, but the large $`N`$ limit provides a non-perturbative framework that includes a similar resummation. The main point of this discussion is the realization that the resummation implied by the large $`N`$ limit provides an effective coupling constant that is well behaved in the infrared limit, thus leading to the conclusion that the simple point-like scattering vertex must be resummed before attempting to compute the damping rate or any other transport coefficient near the critical region.
The analysis in this section reveals the role played by the scale $`\lambda T`$: internal loop momenta $`q\lambda T`$ lead to non-perturbative contributions, in the weak coupling limit $`\lambda 1`$ these non-perturbative scales are classical, on the other hand for $`q\lambda T`$ the effective couplings (either four or three dimensional) are small for weak coupling $`\lambda `$ and the effective vertices coincide with the bare vertices. The implications of this discussion will be important to understand the different contributions to the relaxation rate.
### B The relaxation rate
As discussed above at leading order in the large $`N`$ limit the only contribution to the scalar self-energy is a tadpole $`\lambda \stackrel{}{\mathrm{\Phi }}^2/N𝒪(1)`$, which results in the effective thermal mass $`m_T|TT_c|^{1/2}`$ and is cancelled by the mass counterterm. In this section we consider the theory at the critical temperature where the renormalized temperature dependent mass exactly vanishes.
At next-to-leading order $`𝒪(1/N)`$ the self-energy obtains an absorptive part and is given by the diagram shown in fig. 4.
In appendix A.2 we provide the details necessary to obtain the retarded self-energy in terms of a dispersion relation as in eq.(6), with the spectral density
$`\stackrel{~}{\rho }(\omega ,k)`$ $`=`$ $`{\displaystyle \frac{4\lambda }{N}}{\displaystyle \frac{d^3q}{(2\pi )^3}\frac{dq_0}{2|\stackrel{}{k}+\stackrel{}{q}|}\rho (q_0,q)\left[1+n_{\stackrel{}{q}+\stackrel{}{p}}+n(q_0)\right]\left[\delta (\omega q_0|\stackrel{}{k}+\stackrel{}{q}|)\delta (\omega +q_0+|\stackrel{}{k}+\stackrel{}{q}|)\right]}`$ (35)
$`=`$ $`{\displaystyle \frac{2\lambda }{N}}{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}{\displaystyle \frac{1}{|\stackrel{}{k}+\stackrel{}{q}|}}\{\rho (\omega |\stackrel{}{k}+\stackrel{}{q}|,q)[n_{\stackrel{}{q}+\stackrel{}{k}}n(|\stackrel{}{k}+\stackrel{}{q}|\omega )]`$ (36)
$`+`$ $`\rho (\omega +|\stackrel{}{k}+\stackrel{}{q}|,q)[n_{\stackrel{}{q}+\stackrel{}{k}}n(\omega +|\stackrel{}{k}+\stackrel{}{q}|)]\}.`$ (37)
We performed here the integral over $`q_0`$ by using the delta functions thereby setting the combination $`q_0+|\stackrel{}{q}+\stackrel{}{k}|=\pm \omega `$ for the respective delta functions.
In appendix A.1 we show in detail that
$$\rho (q_0,q)=\frac{1}{\pi }\frac{\mathrm{\Pi }_I(q_0,q)}{\left[1+\mathrm{\Pi }_R(q_0,q)\right]^2+\mathrm{\Pi }_I^2(q_0,q)},$$
(38)
where $`\mathrm{\Pi }_I(q_0,q)`$ is given by the leading order in the large $`N`$ limit of the two loop spectral density (15), as
$`\mathrm{\Pi }_I(q_0,q)=2\lambda \pi {\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{4p|\stackrel{}{p}+\stackrel{}{q}|}}`$ $`\{[1+n_{\stackrel{}{q}+\stackrel{}{p}}+n_\stackrel{}{p}][\delta (q_0|\stackrel{}{p}+\stackrel{}{q}|p)\delta (q_0+|\stackrel{}{p}+\stackrel{}{q}|+p)]`$ (39)
$`+`$ $`[n_\stackrel{}{p}n_{\stackrel{}{q}+\stackrel{}{p}}][\delta (q_0|\stackrel{}{p}+\stackrel{}{q}|+p)\delta (q_0+|\stackrel{}{p}+\stackrel{}{q}|p)]\}.`$ (40)
The first term, proportional to the sum of the occupation factors, corresponds to the two particle cut while the second term proportional to the difference is obviously only present in the medium and corresponds to Landau damping. The real and imaginary parts of the polarization of the auxiliary field $`\mathrm{\Pi }(q_0,q)`$ are related by a dispersion relation, i.e,
$$\mathrm{\Pi }_R(q_0,q)=\frac{1}{\pi }𝑑\omega \mathrm{\Pi }_I(q,\omega )𝒫\frac{1}{\omega q_0}.$$
(41)
Keeping the leading temperature dependence we obtain
$$\mathrm{\Pi }_I(q_0,q)=\frac{\lambda T}{4\pi q}\mathrm{ln}\left[\frac{1e^{\frac{|q_0+q|}{2T}}}{1e^{\frac{|q_0q|}{2T}}}\right]+𝒪\left(\lambda T^0\right).$$
(42)
In appendix B we show explicitly that the leading temperature dependence for the real part of the polarization operator of the auxiliary field is given by
$$\mathrm{\Pi }_R(q_0,q)=\frac{\lambda T}{4q}\left[\mathrm{\Theta }(qq_0)\mathrm{\Theta }(qq_0)\right]+𝒪\left(\lambda \mathrm{ln}T\right)$$
(43)
which in the static limit reduces to (32).
It is clear from eq.(42) that just like in the case of perturbation theory up to two loops, there are two important regions to consider: i) the classical region with $`q_0,qT`$ and ii) the hard region with $`q_0,qT`$ but with either $`|qq_0|T`$ or $`|q+q_0|T`$ the other regions of hard momentum being exponentially suppressed. We will be primarily interested in the case of soft external momentum $`k\lambda T`$ i.e. long-wavelength fluctuations of the order parameter.
In the high temperature limit the spectral density $`\rho (q_0,q)`$ takes the explicit form
$`\rho (k+|\stackrel{}{k}+\stackrel{}{q}|,q)`$ $`=`$ $`{\displaystyle \frac{\frac{4q}{\lambda T}L_+}{\pi ^2\left[\frac{4q}{\lambda T}+1\right]^2+L_+^2}}`$ (44)
$`\rho (k|\stackrel{}{k}+\stackrel{}{q}|,q)`$ $`=`$ $`{\displaystyle \frac{\frac{4q}{\lambda T}L_{}}{\left[\frac{4q\pi }{\lambda T}\right]^2+L_{}^2}}`$ (46)
where we used eqs.(42)-(43) analyzing carefully the support of the theta functions in $`\mathrm{\Pi }_R(q_0,q)`$ and defined
$$L_\pm \mathrm{log}\left|\frac{k\pm |\stackrel{}{k}+\stackrel{}{q}|+q}{k\pm |\stackrel{}{k}+\stackrel{}{q}|q}\right|.$$
Introducing in the integral (37) the dimensionless variable $`y=q/k`$ and $`x=\mathrm{cos}\theta `$ (where $`\theta `$ is the angle between the vectors $`\stackrel{}{k}`$ and $`\stackrel{}{q}`$) and setting $`\omega =k`$ yields for the damping rate (7)
$`\mathrm{\Gamma }(k,T)`$ $`=`$ $`{\displaystyle \frac{k^2}{\pi NT}}{\displaystyle _0^{\mathrm{}}}y^3dy{\displaystyle _1^{+1}}{\displaystyle \frac{dx}{w(x,y)}}\{[{\displaystyle \frac{1}{e^{\beta kw(x,y)}1}}{\displaystyle \frac{1}{e^{\beta k[w(x,y)1]}1}}]{\displaystyle \frac{L_+(x,y)}{\pi ^2\left[\frac{4ky}{\lambda T}+1\right]^2+L_+^2(x,y)}}`$ (47)
$`+`$ $`[{\displaystyle \frac{1}{e^{\beta kw(x,y)}1}}{\displaystyle \frac{1}{e^{\beta k[w(x,y)+1]}1}}]{\displaystyle \frac{L_{}(x,y)}{\left[\frac{4ky\pi }{\lambda T}\right]^2+L_{}^2(x,y)}}\}`$ (48)
where
$$L_\pm (x,y)=\mathrm{log}\left|\frac{1w(x,y)+y}{1w(x,y)y}\right|\text{and}w(x,y)\sqrt{1+y^2+2xy}.$$
For $`\lambda T/4k1`$ the region of small $`y`$ (small momentum) is screened by the resummation of the scattering amplitude and leads to a small contribution to the damping rate of order $`k`$. For $`y>\lambda T/k`$ the screening is not effective and the integrals in eq.(48) are dominated by the neighborhood of the point $`y=y^{}`$ at which $`L_\pm ^2(x,y)`$ is of the same order as the other square in the denominators, i.e, $`[4ky/\lambda T]^2`$. For large $`y`$ we can expand $`w(x,y)`$ as follows
$$w(x,y)=y+x+𝒪\left(\frac{1}{y}\right)=y\left[1+𝒪\left(\frac{1}{y}\right)\right]$$
(49)
We thus find
$$y^{}\frac{\lambda T}{4\pi k}\mathrm{ln}\frac{\lambda T}{4\pi k}$$
\[For $`kT`$ the $`1`$ in the denominator of the first term of eq.(48) is irrelevant in the determination of $`y^{}`$\].
The classical approximation for the occupation numbers, i.e, loop momenta $`T`$ can be used at $`y=y^{}`$ provided $`ky^{}T`$ and $`\lambda 1`$, i.e, for $`kk_{us}`$ where
$$k_{us}\frac{\lambda T}{4\pi }e^{\frac{4\pi }{\lambda }}$$
is the ultrasoft scale.
As it will be discussed in detail below, for $`kk_{us}`$ hard momenta dominate the contributions to the width $`\mathrm{\Gamma }(k,T)`$.
We analyze in the subsequent section the width $`\mathrm{\Gamma }(k,T)`$ in the two regimes : soft for which $`Tkk_{us}`$ and ultrasoft $`kk_{us}`$.
### C Classical contribution
To obtain the contribution from classical momenta $`qT`$ we perform the following approximations: i) approximate $`\mathrm{\Pi }_I(q_0,q)`$ by its limit for $`q_0,qT`$
$$\mathrm{\Pi }_I(q_0,q)=\frac{\lambda T}{4\pi q}\mathrm{ln}\left|\frac{q_0+q}{q_0q}\right|$$
with the real part given by the leading temperature contribution (43) and ii) approximate the Bose-Einstein occupation factors by their classical limit, i.e. in (37) we replace
$$1+n_{\stackrel{}{q}+\stackrel{}{p}}+n(q_0)=T\frac{q_0+|\stackrel{}{q}+\stackrel{}{p}|}{q_0|\stackrel{}{q}+\stackrel{}{p}|}.$$
We thus find from eq.(48) for the classical limit contribution to the damping rate.
$`\mathrm{\Gamma }_{cl}(k,T)`$ $`=`$ $`{\displaystyle \frac{k}{\pi N}}{\displaystyle _0^{\alpha T/k}}y^3dy{\displaystyle _1^{+1}}{\displaystyle \frac{dx}{w^2(x,y)}}\{{\displaystyle \frac{1}{1w(x,y)}}{\displaystyle \frac{L_+(x,y)}{\frac{1}{\lambda _{3,eff}^2(yk)}+L_+^2(x,y)}}`$ (50)
$`+`$ $`{\displaystyle \frac{1}{1+w(x,y)}}{\displaystyle \frac{L_{}(x,y)}{\left[\frac{1}{\lambda _{3,eff}(ky)}\frac{1}{\lambda ^{}}\right]^2+L_{}^2(x,y)}}\}`$ (52)
with $`\lambda _{3,eff}(q)`$ being the effective three dimensional coupling given by (34) and $`\lambda ^{}=1`$ the three dimensional fixed point. We have introduced an explicit upper momentum cutoff $`q_{max}=\alpha T`$ with $`\alpha 1`$ that restricts the integration domain to the region where the classical approximation is valid.
The expression (52) clearly reveals the role of the three dimensional effective coupling $`\lambda _{3,eff}(q)`$ given by (34) and its non-trivial (three dimensional Wilson-Fisher) fixed point $`\lambda ^{}=1`$ reached in the soft limit $`q0`$. The phenomenon of screening of the infrared behavior by the renormalization of the coupling is now explicit, the region of soft loop momentum $`q\lambda T`$ is independent of the coupling and temperature because the effective three dimensional coupling is near its non-trivial fixed point. The only scale in the integral in the soft-momentum region is $`k`$ and a dimensional analysis reveals that the contribution from this region is proportional to $`k`$.
On the other hand, when the loop momentum is $`q\lambda T`$ the renormalization of the coupling is ineffective and the effective coupling coincides with the three dimensional coupling $`\lambda T/4q`$. For weak coupling $`\lambda 1`$ and loop momenta $`q\lambda T`$ the effective three dimensional coupling is $`\lambda _{3,eff}(q)\lambda _3(q)=\lambda T/4q1`$ and the denominators in eq.(52) are dominated by the terms $`1/\lambda _{3,eff}(q)`$. If the logarithms can be neglected we clearly see that this contribution is the same as that given by the integrals $`J_1(k,T)`$ and $`J_2(k,T)`$ given by eqs. (24) in the two loop computation of the damping rate (23), which is proportional to $`\lambda ^2T^2/k`$. This region begins to dominate for $`q>\lambda T\mathrm{ln}(\lambda T/k)`$ when $`1/\lambda _{3,eff}(q)`$ becomes larger than the logarithm in the denominators in both integrals in eq.(52). If the loop momenta are such that $`Tq`$ in this region, i.e. the classical approximation is valid, the contribution of this region to the integral can be estimated by cutting off the integrals in eq.(52) at a lower momentum of order $`q_{min}\lambda T\mathrm{ln}(\lambda T/k)`$. Hence following the same arguments as for the two loops case that led to the estimate (29) we conclude that the region of classical semi-soft loop momentum $`Tqq_{min}`$ leads to a contribution to the damping rate $`\lambda T`$. However as $`k`$ becomes smaller, $`q_{min}`$ approaches the cutoff $`\alpha T`$, i.e, the limit of validity of the classical approximation and the logarithmic terms cannot be neglected. In particular for $`kk_{us}`$ with $`k_{us}`$ the ultrasoft scale introduced above the integral becomes sensitive to momenta of order of $`T`$ and the classical approximation breaks down. For these ultrasoft momenta of the fluctuations the damping rate is determined by the region of hard loop momentum $`qT`$.
This analysis yields to a preliminary assessment of how different regions of loop momentum will contribute to the damping rate:
i) the soft region of loop momentum $`q\lambda T`$ is dominated by the three dimensional fixed point and contributes to the damping rate
$$\mathrm{\Gamma }_{q\lambda T}k$$
ii) the semisoft region of loop momentum $`Tq\lambda T`$ is still dominated by classical modes but the renormalization of the scattering amplitude is irrelevant. If $`q_{min}\lambda T\mathrm{ln}(\lambda T/k)T`$ the logarithms can be neglected and the integrals in eq.(52) behave similarly to the perturbative two loops computation. For momenta $`Tqq_{min}\lambda T\mathrm{ln}(\lambda T/k)`$ the integrals are dominated by the terms in the denominators proportional to $`1/\lambda _3^2(q)`$ leading to a contribution
$$\mathrm{\Gamma }_{Tqq_{min}}\lambda T$$
The validity of the classical approximation and the dominance of semisoft loop momenta is warranted for weak coupling when there is a clear separation between the hard scales with $`qT`$ the semisoft scales with $`Tq\lambda T`$ and the soft scales for which $`\lambda Tq`$.
iii) However when $`q_{min}T`$ the logarithmic terms are the dominant terms in the denominators of the integrands and this region is sensitive to the hard loop momenta $`qT`$ and the classical approximation is not warranted. This region requires the full Bose-Einstein distributions and will be studied in detail below.
After this preliminary assessment, we now provide a quantitative analysis of the different regions.
As argued above in the region $`y\lambda _3(k)=\lambda T/(4k),`$ i.e. $`q\lambda T`$, the integrals in eq.(52) are independent of $`\lambda _3(k)`$ (i.e, of both the coupling constant and temperature) and infrared finite, contributing to the damping rate a term that is proportional to $`k`$. This is the region of loop momenta for which the effective coupling is near the three dimensional Wilson-Fisher fixed point. The screening of the scattering amplitude is ineffective for loop momenta $`q\lambda T`$ i.e, the semisoft scales, in this region $`y\lambda _3(k)\mathrm{ln}\lambda _3(k)`$ and the term $`[y\lambda _3(k)]^2`$ dominates over the logarithms leading to a contribution to the damping rate proportional to $`k\lambda _3(k)=\lambda T/4`$. Since in this region $`y>\lambda _3(k)1`$ we can approximate $`w(x,y)`$ according to eq.(49) and the integrals simplify considerably.
Up to $`1/\lambda _3(k)`$ corrections we find
$`\mathrm{\Gamma }_{cl}(k,T)`$ $`=`$ $`{\displaystyle \frac{k}{N\pi }}{\displaystyle _{\lambda _3(k)}^{\alpha T/k}}dy{\displaystyle _1^{+1}}dx\{{\displaystyle \frac{\mathrm{ln}\frac{2y}{1x}}{\pi ^2\left[\frac{y}{\lambda _3(k)}+1\right]^2+\mathrm{ln}^2\frac{2y}{1+x}}}`$ (53)
$`+`$ $`{\displaystyle \frac{\mathrm{ln}\frac{2y}{1+x}}{\frac{\pi ^2y^2}{\lambda _3^2(k)}+\mathrm{ln}^2\frac{2y}{1+x}}}\left\}\right[1+𝒪\left({\displaystyle \frac{1}{\lambda _3(k)}}\right)].`$ (55)
Now the angular integrals (over the variable $`x`$) can be performed changing the integration variables to $`y=\lambda _3(k)u`$ and expanding in inverse powers of $`\mathrm{ln}\lambda _3(k)`$, since for $`\lambda \alpha 1`$
$$\mathrm{ln}\left[\frac{2\lambda _3(k)u}{1\pm x}\right]=\mathrm{ln}\left[2\lambda _3(k)\right]\left[1+𝒪\left(\frac{1}{\mathrm{ln}\lambda _3(k)}\right)\right].$$
This is certainly a slowly converging approximation for soft momenta but numerical calculations show it to be reliable (see figs. 5-7 below). Then we find from eq.(53)
$`\mathrm{\Gamma }_{cl}(k,T)`$ $`=`$ $`{\displaystyle \frac{2k\lambda _3(k)}{N\pi }}{\displaystyle _1^{\mathrm{}}}𝑑u\left[{\displaystyle \frac{\mathrm{ln}[2\lambda _3(k)]}{\pi ^2(u+1)^2+\mathrm{ln}^2[2\lambda _3(k)]}}+{\displaystyle \frac{\mathrm{ln}[2\lambda _3(k)]}{(\pi u)^2+\mathrm{ln}^2[2\lambda _3(k)]}}\right]\left[1+𝒪\left({\displaystyle \frac{1}{\mathrm{ln}\lambda _3(k)}}\right)\right]`$ (56)
$`=`$ $`{\displaystyle \frac{\lambda T}{2N\pi }}\left[1+𝒪\left({\displaystyle \frac{1}{\mathrm{ln}\lambda _3(k)}}\right)\right].`$ (58)
A detailed analysis of both integrals above reveal the presence of two important scales: a) the cutoff scale $`y=\alpha T/k`$, that is $`q=\alpha T`$ with $`\alpha 1`$ that determines the regime of validity of the classical approximation when the Bose-Einstein occupation factors may be replaced by their classical counterparts and b) a scale $`y^{}\lambda _3(k)\mathrm{ln}\lambda _3(k)+\mathrm{}`$ at which there is a crossover of behavior in the denominators of the integrals. For $`y<y^{},y/\lambda _3(k)\mathrm{ln}y`$ and the denominator is dominated by the logarithm, whereas for $`yy^{}`$, $`y/\lambda _3(k)\mathrm{ln}y`$ and the integrands behave as $`\mathrm{ln}y/y^2`$ which is the same behavior as that in the integrals $`J_1(k,T)`$ and $`J_2(k,T)`$ in eqs.(24) for the two loop computation. In the case when $`y^{}\alpha T/k`$ the integrand falls off very fast and the integral is independent of the upper cutoff.
The result (58) is confirmed by a careful numerical study of the integrals in this range and displayed in figs. 5-7.
The condition that $`y^{}\alpha T/k`$ translates into the following condition for $`k`$
$$k\lambda Te^{\frac{\alpha }{\lambda }}$$
It is clear that $`y^{}`$ becomes of the order of $`T/k`$ and therefore the crossover scale becomes of order $`T`$ for $`kk_{us}\frac{\lambda T}{4\pi }e^{\frac{4\pi }{\lambda }}`$. Hence for wavevectors $`kk_{us}`$ the classical approximation is valid and the damping rate is dominated by semisoft classical loop momenta $`Tq\lambda T`$ and given by
$$\mathrm{\Gamma }_{cl}(k,T)=\frac{\lambda T}{2N\pi }[1+𝒪\left(\frac{1}{\mathrm{ln}\frac{\lambda T}{k}}\right)];\text{for}k\frac{\lambda T}{4\pi }e^{\frac{4\pi }{\lambda }}.$$
(59)
In the opposite limit, i.e. for $`kk_{us}`$ the crossover scale $`y^{}\alpha T/k`$ and the terms $`[y/\lambda _3(k)]^2`$ in the denominators are negligible as compared to $`(\mathrm{ln}y)^2`$ in this range. In this case the integrals can be evaluated by neglecting the $`[y/\lambda _3(k)]^2`$ in the denominators with the result
$$\mathrm{\Gamma }_{cl}(k,T)\frac{\alpha T}{\mathrm{ln}(\alpha T/k)}.$$
This cutoff dependence signals the breakdown of the classical approximation since the integrand is sensitive to hard momenta of order $`q\alpha T`$. For $`kk_{us}`$ the crossover scale $`y^{}`$ becomes of order $`T/k`$, i.e, $`qT`$ and we must keep the full occupation numbers, this is the regime dominated by the hard loop momentum, which is studied below.
Thus we conclude that the non-perturbative region of wavevectors for which the classical approximation is valid is $`\lambda Tk\frac{\lambda T}{4\pi }e^{\frac{4\pi }{\lambda }}`$ and in this region the relaxation rate is given by (59). However for long-wavelength fluctuations with wavevectors $`kk_{us}\frac{\lambda T}{4\pi }e^{\frac{4\pi }{\lambda }}`$ the classical approximation breaks down and we must consider the contribution from the hard loop momenta.
This analysis of the classical contribution reveals that i) the screening of loop momenta $`q\lambda T`$ by the infrared renormalization of the scattering amplitude makes the damping rate a factor $`k/\lambda T1`$ smaller than the lowest order (two loops) computation, ii) the damping rate is independent of momentum for $`k\frac{\lambda T}{4\pi }e^{4\pi /\lambda }`$ and given by en. (59) i.e. there is no critical slowing down in the regime of validity of the classical approximation to this order in the large $`N`$ expansion.
### D Ultrasoft Scale: $`k\frac{\lambda T}{4\pi }e^{4\pi /\lambda }`$
We now focus on the computation of the damping rate in the regime of ultra-soft fluctuations of the order parameter, i.e, $`k\frac{\lambda T}{4\pi }e^{4\pi /\lambda }`$.
In this limit we expand the difference of the occupation numbers to order $`\frac{k}{T}`$ inside the integrand in eq.(48)
$$\frac{1}{e^{\beta kw(x,y)}1}\frac{1}{e^{\beta k[w(x,y)1]}1}=\frac{k}{T}\frac{e^{\beta kw(x,y)}}{\left(e^{\beta kw(x,y)}1\right)^2}+𝒪\left(\frac{k}{T}\right)^2$$
Since in this region the logarithm dominates, we neglect the terms $`4\pi q/\lambda T`$ in the denominators. We thus find
$`\mathrm{\Gamma }(k,T)`$ $`\stackrel{Tk}{=}`$ $`{\displaystyle \frac{k^3}{N\pi T^2}}{\displaystyle _0^{\mathrm{}}}y^3dy{\displaystyle _1^{+1}}{\displaystyle \frac{dx}{w(x,y)}}{\displaystyle \frac{e^{\beta kw(x,y)}}{\left(e^{\beta kw(x,y)}1\right)^2}}\times `$ (62)
$`\left[{\displaystyle \frac{L_+(x,y)}{\pi ^2+L_+^2(x,y)}}+{\displaystyle \frac{1}{L_{}(x,y)}}\right].`$
In order to perform the integration it is convenient to change variables to
$$v\frac{k}{T}[w(x,y)1],\sigma \frac{2k}{T}\frac{w(x,y)1}{y+1w(x,y)}.$$
The width then takes the form
$`\mathrm{\Gamma }(k,T)`$ $`\stackrel{Tk}{=}`$ $`{\displaystyle \frac{2T}{N\pi }}[{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{e^vv^3dv}{\left(e^v1\right)^2}}{\displaystyle _v^{\mathrm{}}}{\displaystyle \frac{d\sigma }{\sigma ^2}}{\displaystyle \frac{\mathrm{ln}\left(1+\frac{T\sigma }{k}\right)}{\pi ^2+\mathrm{ln}^2\left(1+\frac{T\sigma }{k}\right)}}`$ (63)
$`+`$ $`{\displaystyle _{2k/T}^{\mathrm{}}}{\displaystyle \frac{e^vv^3dv}{\left(e^v1\right)^2}}{\displaystyle _v^{\mathrm{}}}{\displaystyle \frac{d\sigma }{\sigma ^2}}{\displaystyle \frac{1}{\mathrm{ln}\left(\frac{T\sigma }{k}1\right)}}\left]\right[1+𝒪\left({\displaystyle \frac{k}{T}}\right)].`$ (65)
We can further approximate these expressions by expanding in inverse powers of $`\mathrm{ln}\frac{T}{k}`$. We set,
$`{\displaystyle _v^{\mathrm{}}}{\displaystyle \frac{d\sigma }{\sigma ^2}}{\displaystyle \frac{\mathrm{ln}\left(1+\frac{T\sigma }{k}\right)}{\pi ^2+\mathrm{ln}^2\left(1+\frac{T\sigma }{k}\right)}}={\displaystyle _v^{\mathrm{}}}{\displaystyle \frac{d\sigma }{\sigma ^2}}{\displaystyle \frac{1}{\mathrm{ln}\frac{T\sigma }{k}}}\left[1+𝒪\left({\displaystyle \frac{1}{\mathrm{ln}\frac{T}{k}}}\right)\right]={\displaystyle \frac{1}{v\mathrm{ln}\frac{T}{k}}}\left[1+𝒪\left({\displaystyle \frac{1}{\mathrm{ln}\frac{T}{k}}}\right)\right]`$ (66)
(67)
$`{\displaystyle _v^{\mathrm{}}}{\displaystyle \frac{d\sigma }{\sigma ^2}}{\displaystyle \frac{1}{\mathrm{ln}\left(\frac{T\sigma }{k}1\right)}}={\displaystyle \frac{1}{v\mathrm{ln}\frac{T}{k}}}\left[1+𝒪\left({\displaystyle \frac{1}{\mathrm{ln}\frac{T}{k}}}\right)\right].`$ (68)
The asymptotic form of the width thus becomes
$$\mathrm{\Gamma }(k,T)\stackrel{kk_{us}}{=}\frac{4\pi T}{3N\mathrm{ln}\frac{T}{k}}\left[1+𝒪\left(\frac{1}{\mathrm{ln}\frac{T}{k}}\right)\right],$$
(69)
where we used the integral
$$_0^{\mathrm{}}\frac{e^vv^2dv}{\left(e^v1\right)^2}=\frac{\pi ^2}{3}.$$
There are two important noteworthy features of this result: i) the damping rate for ultrasoft fluctuations is independent of the coupling and ii) critical slowing down of long-wavelength fluctuations emerges in the ultrasoft momentum limit with the damping rate vanishing only logarithmically as $`k0`$.
The intermediate regime between the soft and ultrasoft scales is difficult to study analytically, we therefore studied the damping rate in a wide range of momentum $`k`$ numerically.
Figs. 5-7 display the dimensionless ratio $`N\mathrm{\Gamma }(k,T)/T`$ as a function of $`\mathrm{ln}\frac{T}{k}`$ for three fixed values of the coupling: $`\lambda =12.0,\mathrm{\hspace{0.33em}1.0}`$ and $`0.2`$ respectively as obtained via a numerical integration of eq. (48). We see that the damping rate is a monotonically decreasing function of $`\frac{T}{k}`$ for large enough values of $`\frac{T}{k}`$. The smaller is the coupling, the slower $`N\mathrm{\Gamma }(k,T)/T`$ decreases as a function of $`\frac{T}{k}`$. Furthermore we have established numerically the reliability of the results in the soft and ultrasoft regimes, thus confirming our detailed analysis in these cases.
## V Near critical regime: $`\stackrel{}{k}=0,|TT_c|T_c`$
Having understood in detail the critical case we are now in position to complete the study of relaxation by considering the near-critical case, i.e. $`|TT_c|T_c`$. The general case of $`\stackrel{}{k}0,TT_c`$ is rather complicated but we can learn much by focusing on the homogeneous case $`\stackrel{}{k}=0`$. There are two important modifications of the previous results that are required to study in detail the near-critical case:
i:) To leading order in the large $`N`$ limit, the finite temperature effective mass squared is given by
$$m_T^2\lambda (T^2T_c^2)$$
therefore the effective mass (inverse of the correlation length) vanishes near the critical temperature as $`m_T|TT_c|^{1/2}`$ which is the mean-field behavior consistent. Since the absorptive part of the self-energy is next to leading order $`𝒪(1/N)`$ we consistently use this effective mass near the critical point. Therefore to this order the frequencies are given by $`\omega _\stackrel{}{p}=\sqrt{\stackrel{}{p}^2+m_T^2}`$.
ii:) The effective static scattering amplitude can be obtained by replacing the massless Matsubara propagators in (31) by the corresponding massive ones, and is now given by
$$\lambda _{eff}(\frac{q}{m_T},\frac{T}{m_T})=\frac{\lambda }{1+\frac{\lambda T}{2\pi q}\text{arctg}\left[\frac{q}{2m_T}\right]}$$
which now reveals the vanishing of the effective coupling
$$\lambda _{eff}(0,\frac{T}{m_T})=\frac{\lambda }{1+\frac{\lambda T}{4\pi m_T}}$$
(70)
for $`m_T0`$ \[compare with eq.(33\].
The critical region of relevance corresponds to $`m_T\lambda T`$ for $`TT_c^+`$. In the case under consideration, for homogeneous fluctuations of the order parameter ($`\stackrel{}{k}=0`$) the only dimensionful quantity is the effective mass $`m_T`$ which regulates the infrared behavior of the integrals. Thus, just as in the critical case studied above two different regimes emerge which we refer to as: i) the semicritical regime $`m_T\lambda T/(4\pi )`$; ii) the ultracritical regime $`m_T\lambda T/(4\pi )e^{4\pi /\lambda }`$. It will become clear below that the semicritical and the ultracritical regimes correspond respectively to the soft and the ultrasoft regimes discussed at $`T=T_c`$, $`k0`$. Since the relevant loop momenta are semisoft, $`Tq\lambda T`$ and hard $`qT`$ it is clear that the effective coupling (70) behaves just as in the critical case for $`\stackrel{}{k}0`$ studied above since for this range of loop momenta $`q/m_T1`$.
In order to compute the damping rate we need the general expression for the resummed spectral density $`\stackrel{~}{\rho }(q_0,q)`$ in presence of a non-zero thermal mass, which is now given by \[see eq.(37)\]
$`\stackrel{~}{\rho }(\omega ,k)`$ $`=`$ $`{\displaystyle \frac{4\lambda }{N}}{\displaystyle \frac{d^3q}{(2\pi )^3}\frac{dq_0}{2\omega _{|\stackrel{}{k}+\stackrel{}{q}|}}\rho (q_0,q)\left[1+n(\omega _{|\stackrel{}{q}+\stackrel{}{k}|})+n(q_0)\right]\left[\delta (\omega q_0\omega _{|\stackrel{}{k}+\stackrel{}{q}|})\delta (\omega +q_0+\omega _{|\stackrel{}{k}+\stackrel{}{q}|})\right]}`$ (71)
$`=`$ $`{\displaystyle \frac{4\lambda }{N}}{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}{\displaystyle \frac{1}{2\omega _{|\stackrel{}{k}+\stackrel{}{q}|}}}\{\rho (\omega \omega _{|\stackrel{}{k}+\stackrel{}{q}|},q)[n(\omega _{|\stackrel{}{q}+\stackrel{}{k}|})n(\omega _{|\stackrel{}{k}+\stackrel{}{q}|}\omega )]`$ (72)
$`+`$ $`\rho (\omega +\omega _{|\stackrel{}{k}+\stackrel{}{q}|},q)[n(\omega _{|\stackrel{}{q}+\stackrel{}{k}|})n(\omega +\omega _{|\stackrel{}{k}+\stackrel{}{q}|})]\}`$ (73)
where $`\omega _k^2m_T^2+k^2`$, $`n(q_0)`$ is the Bose-Einstein distribution function (17) and $`\rho (\omega \pm \omega _{|\stackrel{}{k}+\stackrel{}{q}|},q)`$ is the massive spectral density at leading order in the $`1/N`$ expansion, given by equation (38) with the following expression for $`\mathrm{\Pi }_I(\omega \pm \omega _{|\stackrel{}{k}+\stackrel{}{q}|},q)`$ :
$`\mathrm{\Pi }_I(\omega ,p)`$ $`=`$ $`\left\{{\displaystyle \frac{\lambda }{8\pi }}\sqrt{1{\displaystyle \frac{4m_T^2}{\omega ^2p^2}}}\text{sgn}(\omega )+{\displaystyle \frac{\lambda T}{4\pi p}}\mathrm{ln}\left[{\displaystyle \frac{1e^{\beta \omega _p^+}}{1e^{\beta \omega _p^{}}}}\right]\right\}\mathrm{\Theta }(\omega ^2p^24m_T^2)`$ (75)
$`+{\displaystyle \frac{\lambda T}{4\pi p}}\mathrm{ln}\left[{\displaystyle \frac{1e^{\beta \omega _p^+}}{1e^{\beta \omega _p^{}}}}\right]\mathrm{\Theta }(p^2\omega ^2)`$
$`\omega _p^\pm `$ $`=`$ $`\left|{\displaystyle \frac{\omega }{2}}\pm {\displaystyle \frac{p}{2}}\sqrt{1{\displaystyle \frac{4m_T^2}{\omega ^2p^2}}}\right|.`$ (76)
We notice that (75) reduces to (42) in the critical limit $`m_T/T0`$.
Keeping the leading correction in $`m_T^2`$ yields in the high temperature regime,
$`\mathrm{\Pi }_I(\omega ,p)=`$ $`{\displaystyle \frac{\lambda T}{4\pi p}}\left\{\mathrm{log}\right|{\displaystyle \frac{\omega +p}{\omega p}}|{\displaystyle \frac{4m_T^2p^2\text{sgn}(\omega )}{(\omega ^2p^2)^2}}+[\theta (\omega ^2p^24m_T^2)\theta (\omega ^2p^2)]\mathrm{log}|{\displaystyle \frac{\omega +p}{\omega p}}|`$ (78)
$`+𝒪(m_T^4)\}`$
The real part $`\mathrm{\Pi }_R(\omega ,p)`$ can now be obtained via the dispersion relation (41),
$`\mathrm{\Pi }_R(q_0,q)`$ $`=`$ $`{\displaystyle \frac{\lambda T}{4q}}\left\{\mathrm{\Theta }(qq_0)\mathrm{\Theta }(qq_0)+{\displaystyle \frac{4m_T^2}{\pi ^2(q^2q_0^2)}}\left[\mathrm{ln}{\displaystyle \frac{q^2}{m_T^2}}+1{\displaystyle \frac{q^2}{q^2q_0^2}}\left({\displaystyle \frac{q_0^2}{q^2}}1+\mathrm{ln}{\displaystyle \frac{q^2}{q_0^2}}\right)\right]\right\}`$ (80)
$`+𝒪(m_T^4)`$
The damping rate for homogeneous configurations is given by
$$\mathrm{\Gamma }_0(m_T,T)=\underset{k0}{lim}\frac{\pi \stackrel{~}{\rho }(\omega _k,k)}{2\omega _k}=\frac{\pi \stackrel{~}{\rho }(m_T,0)}{2m_T},$$
which is now given explicitly by
$`\mathrm{\Gamma }_0(m_T,T)={\displaystyle \frac{\pi \lambda }{m_TN}}{\displaystyle \frac{d^3q}{(2\pi )^3}}`$ $`{\displaystyle \frac{1}{\omega _q}}\{[n(\omega _q)n(\omega _qm_T)]\rho (m_T\omega _q,q)`$ (82)
$`+[n(\omega _q)n(\omega _q+m_T)]\rho (m_T+\omega _q,q)\}.`$
The integral (82) is much simpler than the analogous expression for $`k0`$, since the angular integration is trivial. Nevertheless, a complete evalution of (82) requires a numerical integration. We refer to figures 8 and 9 for a numerical evalutation of the dimensionless ratio $`N\mathrm{\Gamma }_0(m_T,T)/T`$ in the intermediate ($`\lambda =1.0`$) and small coupling regimes ($`\lambda =0.2`$).
Just as in the critical case in the near-critical regime with $`m_T\lambda T/(4\pi )`$ the integral (82) is dominated by loop momenta of order $`qq^{}=\frac{\lambda T}{4\pi }\mathrm{ln}\frac{\lambda T}{4\pi m_T}\frac{\lambda T}{4\pi }m_T`$ hence we can approximate $`\mathrm{\Pi }_I(q_0,q)`$ and $`\mathrm{\Pi }_R(q_0,q)`$ at $`q_0=m_T\pm \omega _qm_T\pm q`$ as follows
$`\mathrm{\Pi }_I(m_T\pm q,q)`$ $`=`$ $`\pm {\displaystyle \frac{\lambda T}{4\pi q}}\left[\mathrm{ln}{\displaystyle \frac{2q}{m_T}}1+𝒪\left({\displaystyle \frac{m_T^2}{q^2}}\right)\right]`$ (83)
$`\mathrm{\Pi }_R(m_T\pm q,q)`$ $`=`$ $`{\displaystyle \frac{\lambda T}{4q}}\left\{\mathrm{\Theta }(qq_0)\mathrm{\Theta }(qq_0){\displaystyle \frac{2m_T}{\pi ^2q}}\left[\mathrm{ln}{\displaystyle \frac{q^2}{m_T^2}}+1\right]+𝒪\left({\displaystyle \frac{m_T^2}{q^2}}\right)\right\}`$ (85)
where we have used eqs.(78) and (80). Thus as anticipated by the discussion of the effective static scattering amplitude in the near critical case (70) we see that the real part of the polarization is indeed similar to the one in the critical case for the relevant loop momenta up to corrections of order $`m_T/qm_T/\lambda T`$ in the region of semisoft loop momenta.
Therefore neglecting terms of order $`m_T/\lambda T`$ which are negligible in the region of interest, we find for the spectral densities expressions similar to these in the critical case :
$$\rho (m_T+\omega _q,q)=\frac{4q}{\lambda T}\frac{L(q)}{\pi ^2\left[\frac{4q}{\lambda T}+1\right]^2+L^2(q)},\rho (m_T\omega _q,q)=\frac{4q}{\lambda T}\frac{L(q)}{\left[\frac{4q\pi }{\lambda T}\right]^2+L^2(q)},$$
where we have introduced
$$L(q)=\mathrm{log}\frac{2q}{em_T}.$$
An analysis similar to that in the critical case reveals that soft loop momenta $`q\lambda T`$ are effectively screened by the renormalization of the scattering amplitude in the near critical region, leading to a contribution of order $`m_T`$ to the damping rate. Semisoft and hard loop momenta $`q>\lambda T`$ are not screened by the resummation of the scattering amplitude and determine the leading contributions to the damping rate.
As argued above, the dominat loop momenta in the integral (82) are of order $`qq^{}=\frac{\lambda T}{4\pi }\mathrm{ln}\frac{\lambda T}{4\pi m_T}\frac{\lambda T}{4\pi }m_T`$ hence we can approximate
$$n(\omega _q)n(\omega _qm_T)n(\omega _q+m_T)n(\omega _q)n^{}(q)m_T,\omega _qq$$
allowing an analytical estimate the damping rate from the approximate expression
$$\mathrm{\Gamma }_0(m_T,T)\stackrel{\lambda Tm_T}{=}\frac{4}{NT}_a^{\mathrm{}}\frac{dq}{2\pi ^2}q^2n^{}(q)\left\{\frac{L(q)}{\pi ^2\left[\frac{4q}{\lambda T}+1\right]^2+L(q)^2}+\frac{L(q)}{\left[\frac{4q\pi }{\lambda T}\right]^2+L(q)^2}\right\}.$$
(86)
The resulting integrals are infrared finite but having recognized that the leading contribution arises from the semisoft loop momenta $`q\lambda T`$ we have introduced an explicit infrared cutoff $`a=𝒞m_T`$ with $`𝒞1`$. Since the integral is dominated by semisoft and hard loop momenta $`q^{}\lambda T`$ we can approximate further $`\frac{4q}{\lambda T}+1\frac{4q}{\lambda T}`$ whence the two contributions to the damping rate coincide. Moreover, the dependence in $`a`$ is negligible in the critical limit. This is confirmed by our numerical analysis which uses the exact expression (82), the results of which are displayed in figures 8 and 9. The integrals in (86) again reveal a crossover scale $`q^{}`$ at which $`4\pi q^{}/\lambda TL(q^{})`$. For $`qq^{}`$ the logarithmic term $`L^2(q)`$ dominates in the denominators and for $`qq^{}`$ the term $`(4\pi q^{}/\lambda T)^2`$ dominates and the integrand falls off just as in the perturbative two loops case.
We now distinguish between the following two possibilities:
* $`\lambda T/(4\pi )e^{4\pi /\lambda }m_T\lambda T/(4\pi )`$, to which we refer as the semicritical regime. In this case $`q^{}T`$ is soft and the classical approximation to the Bose-Einstein distribution functions applies. Furthermore, we can expand in inverse powers of the logarithm $`\mathrm{log}\lambda _3(m_T)=\mathrm{log}\frac{\lambda T}{4m_T}`$ and a straightforward analysis along the lines presented for the critical case reveals that the damping rate is approximatively constant and given by
$$\mathrm{\Gamma }_0(m_T,T)=\frac{\lambda T}{2\pi N}\left[1+𝒪\left(\frac{1}{\mathrm{ln}\lambda _3(m_T)}\right)\right].$$
This is the same result as in the case $`T=T_c,k0`$, Eq.(58). This result is of course expected, in the semisoft region of loop momentum $`Tq\lambda T`$ and for $`\lambda T/(4\pi )e^{4\pi /\lambda }m_T\lambda T/(4\pi )`$ the screening of the scattering amplitude is ineffective and the term $`(q/\lambda T)`$ dominates over the logarithms, therefore the dependence on the mass is negligible in this region.
* $`m_T\lambda T/(4\pi )e^{4\pi /\lambda }`$, which we refer to as the ultracritical regime. In this case $`q^{}T`$ is hard and the classical approximation breaks down and the full Bose-Einstein occupation factors must be kept. In this regime the logarithm gives the dominant contribution in the denominators and the asymptotic damping rate is given by
$$\mathrm{\Gamma }_0(m_T,T)=\frac{4}{\pi NT}_a^{\mathrm{}}𝑑qq^2\frac{n^{}(q)}{\mathrm{log}\frac{2q}{em_T}}=\frac{4\pi }{3N}\frac{T}{\mathrm{log}(T/m_T)}\left[1+𝒪\left(\frac{1}{\mathrm{ln}\frac{T}{m_T}}\right)\right],$$
(87)
which is exactly the same result as in equation (69) with the momentum scale $`k`$ replaced by the thermal mass $`m_T`$. It must be noticed that in this regime there is no dependence on the coupling $`\lambda `$.
Thus we conclude that the relaxation rate near the critical point $`|TT_c|\lambda T_c`$ and $`\stackrel{}{k}=0`$ has the same features as the critical rate for $`k0`$ and $`T=T_c`$, provided we exchange the infrared scales $`m_T`$ and $`k`$.
Discussion of the results:
The two loops calculation in perturbation theory at the critical point revealed the importance of the different scales of loop momentum. The loop integrals are dominated by the contribution of the soft momentum scales $`q\lambda T`$. The contribution from semisoft loop momenta $`Tq\lambda T`$ is subdominant by a factor $`(k/\lambda T)\mathrm{ln}[\lambda T/k]1`$ in the long-wavelength limit $`k\lambda T`$ and the contribution from hard loop momentum modes $`qT`$ is suppressed even further in the weak coupling limit by an extra power of the coupling $`\lambda `$.
The large $`N`$ limit leads to a non-perturbative resummation and results in an infrared renormalization of the static scattering amplitude as a consequence of dimensional reduction and crossover to an effective three dimensional theory for momenta $`q\lambda T`$.
The effective three dimensional coupling that emerges from this analysis of the static scattering amplitude at the critical point is $`\lambda _{3,eff}(q)=\lambda T/(4q+\lambda T)`$ which is driven to the Wilson-Fisher three dimensional fixed point as $`q0`$. Thus soft loop momenta $`q\lambda T`$ are effectively screened by this infrared renormalization of the coupling but semisoft loop momenta $`q>\lambda T`$ are coupled with the three dimensional coupling $`\lambda _3(q)=\lambda T/4q`$ and infrared screening is ineffective for these. The importance of this effective coupling for the damping rate can be understood intuitively from figures 2, 3 and 4: the resummation of bubbles that leads to the effective scattering amplitude also renormalizes the spectral density that determines the self-energy, as shown in figure 4.
This is precisely the most important mechanism that leads to our results in the large $`N`$ limit. Whereas the lowest order perturbative calculation was dominated by the contribution of soft loop momenta $`q\lambda T`$ these are effectively screened by the infrared renormalization of the coupling which is near the three dimensional fixed point. The dominant contribution now arises from the semisoft $`q\lambda T`$ and hard $`qT`$ loop momentum. In the perturbative computation at two loops these scales provided subleading contributions of order $`\lambda T`$ and $`\lambda ^2T`$ respectively to the damping rate.
Clearly the resummation via the large $`N`$ approximation incorporates the screening of the soft loop momentum scales but also reveals the emergence of an ultrasoft scale $`k_{us}\frac{\lambda T}{4\pi }e^{\frac{4\pi }{\lambda }}`$. At the critical temperature for $`kk_{us}`$ or for homogeneous fluctuations near the critical point $`T_cm_Tk_{us}`$ the damping rate is determined by the contribution of the classical, semisoft loop momentum scales $`Tq\lambda T`$, the soft scales $`\lambda Tq`$ being screened by the infrared renormalization of the coupling and the crossover to a three dimensional effective theory. For long-wavelength fluctuations of the order parameter at the critical temperature with $`kk_{us}`$ or for homogeneous fluctuations near criticality in the ultracritical region $`m_Tk_{us}`$ the classical approximation breaks down and the damping rate is completely determined by the hard loop momenta $`qT`$. Critical slowing down, i.e. the vanishing of the relaxation rate as $`k0`$ or for $`k=0`$ as $`TT_c^+`$ only emerges in this ultrasoft limit as shown by equations (69) and (87).
Thus, whereas the large $`N`$ expansion has provided a consistent resummation and the important ingredient of screening of the couplings for the soft loop momentum modes and leads to critical slowing down of long wavelength fluctuations important limitations of the results obtained here remain. As we argued in the beginning sections a quasiparticle interpretation of the long-wavelength collective excitations of the order parameter requires that the resonance parameter $`\mathrm{\Gamma }(k,T)/\omega _p(k)1`$ with $`\omega _p(k)`$ being the position of the quasiparticle pole or effectively the microscopic time scale of oscillations of these fluctuations. To leading order in the large $`N`$ limit $`\omega _p(k)=\sqrt{k^2+m_T^2}`$ in the calculation of the damping rate and the resonance parameter. Although critical slowing down emerges from the large $`N`$ limit, we see that our results to this order indicate that $`\mathrm{\Gamma }(k,T)/\omega _p(k)1`$ for $`k,m_T0`$.
There are several possible alternatives: either the quasiparticle picture is not appropriate to describe the collective fluctuations of the order parameter at or near the critical point or further resummations and or other contributions must be taken into account to obtain a description of critical slowing down of collective fluctuations that can be understood within a quasiparticle picture. In particular an assessment of i) vertex corrections, and ii) wave function renormalization must be pursued which, however, are beyond the leading order in the large $`N`$ studied here and thus outside the scope and goals of this article. We are currently studying these contributions and expect to report our conclusions in a forthcoming article.
## VI Conclusions and further questions
In this article we have began the program of studying transport and relaxation at and near the critical point in second order phase transitions. The focus here is to provide a systematic study of critical slowing down from first principles in a phenomenologically motivated quantum field theory. The ultimate goal of this program is to assess the potential experimental signatures associated with the critical slowing down of long wavelength fluctuations at the chiral phase transition. Obtaining a robust understanding of such phenomena will have important implications in the QGP and or chiral phase transitions in early universe cosmology and ultrarelativistic heavy ion collisions.
Our study reveals novel phenomena that require a non-perturbative framework for their consistent and sistematic treatment which is different from the hard thermal loop resummation program used in gauge theories.
Whereas critical slowing down has been studied thoroughly in classical critical phenomena and these results were used for preliminary estimates of the correlation length at freezeout in heavy ion collisions, we are aware of only one prior attempt to study critical slowing down in a full relativistic quantum field theory and a similar recent analysis. In references the Wilson renormalization group was used to explore the relaxation of the $`k=0`$ mode of the order parameter slightly away from criticality working at $`N=1`$. The final results of are that for $`t=(TT_c)/T_c`$ approaching the critical limit the damping rate for homogenous configurations vanishes as $`t^\nu \mathrm{log}t`$ with $`\nu 0.50.6`$. The analysis of, relies on a truncation of the exact renormalization group equations and their numerical evolution. In our opinion there are two very important limitations in this approach: i) in refs. the absorptive parts were treated in a rather simplistic manner and associated with the scattering vertex rather than the self-energy, but more importantly, this simplified treatment does not include consistently the Landau damping and multiple particle thresholds that are the important ingredients in a consistent and sistematic description of damping and relaxation. ii) the method relies on the numerical evolution of a set of equations that had been truncated without a clear control or understanding of the errors induced by this truncation. In particular, it is not clear if the result $`\mathrm{\Gamma }(t)t^\nu \mathrm{log}t`$ is stable with respect to different truncations.
In our study we have systematically focused on the important aspects associated with Landau damping, many-particle threshold effects and a consistent study of real-time phenomena at finite temperature. As is evident in our study of absorptive parts of the self-energy in sections II and III a simplified treatment that does not include consistently these can hardly reveal the rich hierarchy of scales and the different physics associated with these: the soft scale $`k_s\lambda T/(4\pi )`$ and the ultrasoft scale $`k_{us}\lambda T\mathrm{exp}[4\pi /(\lambda T)]`$. Consistently to next to leading order in the large $`N`$ limit we see that slowing down of relaxation of long-wavelength fluctuations only begins to emerge at the ultra-soft scale and in contrast to the results obtained in, we obtain $`\mathrm{\Gamma }T/[N\mathrm{log}t]`$. Thus, although there is agreement on the statement that the relaxation rate vanishes at criticality the consistent large $`N`$ resummation leads to a very different behavior of the relaxation rate.
Our main results can be summarized as follows: a consistent treatment of critical slowing down and of transport phenomena at or near a critical point requires a non-perturbative framework to resum the contributions from soft loop momenta which is different from the hard thermal loop program of abelian and non-abelian gauge plasmas. A perturbative two loop calculation reveals clearly the emergence of a hierarchy of loop momentum scales from hard $`qT`$ to semisoft $`Tq\lambda T`$ and soft $`\lambda Tq`$, for weak coupling $`\lambda 1`$ the scales in this hierarchy are widely separated and the semisoft and soft scales are classical. Recognizing the shortcomings of a perturbative treatment for long-wavelength fluctuations, we implemented a non-perturbative resummation via the large $`N`$ limit to next to leading order. The large $`N`$ limit provides a consistent non-perturbative framework for resummation of infrared contributions. It clearly displays the infrared renormalization of the scattering amplitude in the static limit at or near the critical point and the crossover to three dimensional physics for soft loop momentum. The resummation of the scattering amplitude leads to an effective three dimensional coupling that interpolates between the bare coupling for loop momenta $`q\lambda T`$ and the three dimensional Wilson-Fisher fixed point for $`q\lambda T`$.
The infrared renormalization of the effective coupling screens the contribution from soft loop momentum to the self-energy and the relaxation rate, which is now dominated by the contribution of semisoft and hard loop momenta $`Tq\lambda T`$. Furthermore a new ultrasoft (in the weak coupling limit) non-perturbative scale emerges $`k_{us}\frac{\lambda T}{4\pi }e^{\frac{4\pi }{\lambda }}`$ that signals the breakdown of the classical approximation and the dominance of hard loop momentum modes.
For $`k,m_Tk_{us}`$ the damping rate is dominated by the classical semisoft scales and given by $`\mathrm{\Gamma }(k,T)=\frac{\lambda T}{2N\pi }`$ whereas for $`kk_{us}`$ the hard loop momenta region dominates and leads to the damping rate $`\mathrm{\Gamma }(k,T)=\frac{4\pi T}{3N\mathrm{ln}\frac{T}{k}}`$ at criticality or $`\mathrm{\Gamma }_0(m_T,T)=\frac{4\pi T}{3N\mathrm{ln}\frac{T}{m_T}}`$ near criticality for homogeneous fluctuations, which reveal the slowing down of relaxation of critical ultrasoft fluctuations with a damping rate that is independent of the coupling.
As discussed above, however, these results and those found in seem to indicate a breadown of the quasiparticle picture of collective excitations of the order parameter because the resonance parameter $`\mathrm{\Gamma }(k,T)/\omega _p(k)1`$ in the long-wavelength limit and the excitation decays on time scales much shorter than the natural oscillation time $`\omega _p^1(k)`$. At this stage it is not clear if this feature is a true physical manifestation of relaxation of collective excitations at or near the critical point or that further resummation and other contributions that are beyond the leading order in the large $`N`$ must be accounted for. We are currently studying this possibility by introducing the renormalization group at finite temperature and analyzing in detail the contribution from vertex and wave function renormalizations and expect to report on further understanding on these issues in a forthcoming article. At this stage our study has revealed a wealth of new phenomena and a hierarchy of scales which will require a deeper understanding for a complete and consistent treatment of transport and eventually hydrodynamics near or at the critical point. Only a thorough understanding of these phenomena can lead to an unambiguous assessment of the phenomenological implications of critical fluctuations either in the formation of cosmological relics in the early universe or in experimental observables in ultrarelativistic heavy ion collisions thus motivating and justifying their study.
Acknowledgements
D. B. thanks the N.S.F for partial support through grant awards: PHY-9605186 and INT-9815064 and LPTHE (University of Paris VI and VII) for warm hospitality and partial support, he also thanks S. Raja for interesting conversations. H. J. de Vega thanks the Dept. of Physics at the Univ. of Pittsburgh for hospitality. We thank the CNRS-NSF cooperation programme for partial support. The early stages of the work of M.S. were supported by a grant from Padova University.
## A Equations of motion and spectral densities in the large $`N`$ limit
We study the relaxation of the order parameter via the real time description of non-equilibrium quantum field theory. This formulation requires the time evolved density matrix and is cast in terms of a path integral along a contour in complex time, a forward branch corresponds to the forward time evolution of the density matrix via the unitary time evolution operator and the backward branch represents the inverse unitary time evolution that post-multiplies the density matrix. Consequently there are four propagators: corresponding to fields on either branch. For a more complete description of this formulation the reader is referred to and references therein. The main ingredient in this program are the free field Wightmann and Green’s functions for the bosonic field $`\stackrel{}{\mathrm{\Phi }}`$. In terms of the spatial Fourier transform of the bosonic fields $`\stackrel{}{\mathrm{\Phi }}_\stackrel{}{k}`$ these are given by
$`\mathrm{\Phi }_{\stackrel{}{k},a}(t)\mathrm{\Phi }_{\stackrel{}{k},b}(t^{})_o\mathrm{\Phi }_{\stackrel{}{k},a}^{}(t)\mathrm{\Phi }_{\stackrel{}{k},b}^+(t^{})_o=i\delta _{a,b}G_\stackrel{}{k}^>(t,t^{}),`$ (A1)
(A2)
$`\mathrm{\Phi }_{\stackrel{}{k},a}(t^{})\mathrm{\Phi }_{\stackrel{}{k},b}(t)_o\mathrm{\Phi }_{\stackrel{}{k},a}^{}(t^{})\mathrm{\Phi }_{\stackrel{}{k},b}^+(t)_o=i\delta _{a,b}G_\stackrel{}{k}^<(t,t^{}),`$ (A3)
(A4)
$`\mathrm{\Phi }_{\stackrel{}{k},a}^+(t)\mathrm{\Phi }_{\stackrel{}{k},b}^+(t^{})_o=i\delta _{a,b}G_\stackrel{}{k}^{++}(t,t^{}),`$ (A5)
(A6)
$`\mathrm{\Phi }_{\stackrel{}{k},a}^{}(t)\mathrm{\Phi }_{\stackrel{}{k},b}^{}(t^{})_o=i\delta _{a,b}G_\stackrel{}{k}^{}(t,t^{}),`$ (A7)
$`G_\stackrel{}{k}^{++}(t,t^{})=G_\stackrel{}{k}^>(t,t^{})\mathrm{\Theta }(tt^{})+G_\stackrel{}{k}^<(t,t^{})\mathrm{\Theta }(t^{}t),`$ (A8)
(A9)
$`G_\stackrel{}{k}^{}(t,t^{})=G_\stackrel{}{k}^>(t,t^{})\mathrm{\Theta }(t^{}t)+G_\stackrel{}{k}^<(t,t^{})\mathrm{\Theta }(tt^{}),`$ (A10)
(A11)
$`G_\stackrel{}{k}^+(t,t^{})=G_\stackrel{}{k}^<(t,t^{}),`$ (A12)
(A13)
$`G_\stackrel{}{k}^+(t,t^{})=G_\stackrel{}{k}^>(t,t^{}),`$ (A14)
(A15)
$`G_\stackrel{}{k}^>(t,t^{})=G_\stackrel{}{k}^<(t^{},t).`$ (A16)
(A17)
where $`A(t)B(t^{})=\text{Tr}[A(t)B(t^{})\rho (0)]`$ denotes the expectation value of Heisenberg field operators with respect to the initial normalized density density matrix which is taken to describe a thermal state and the subscript $`o`$ refers to free fields. It is clear that these real time propagators satisfy the identity:
$$G_\stackrel{}{k}^{++}(t,t^{})+G_\stackrel{}{k}^{}(t,t^{})G_\stackrel{}{k}^+(t,t^{})G_\stackrel{}{k}^+(t,t^{})=0.$$
The retarded and advanced propagators are defined as
$`G_{\mathrm{R},\stackrel{}{k}}(t,t^{})`$ $`=`$ $`G_\stackrel{}{k}^{++}(t,t^{})G_\stackrel{}{k}^+(t,t^{})=\left[G_\stackrel{}{k}^>(t,t^{})G_\stackrel{}{k}^<(t,t^{})\right]\mathrm{\Theta }(tt^{}),`$
$`G_{\mathrm{A},\stackrel{}{k}}(t,t^{})`$ $`=`$ $`G_\stackrel{}{k}^{++}(t,t^{})G_\stackrel{}{k}^+(t,t^{})=\left[G_\stackrel{}{k}^<(t,t^{})G_\stackrel{}{k}^>(t,t^{})\right]\mathrm{\Theta }(t^{}t),`$
Where for the cases under consideration with fields in thermal equilibrium
$`G_\stackrel{}{k}^>(t,t^{})`$ $`=`$ $`{\displaystyle \frac{i}{2k}}\left[[1+n_\stackrel{}{k}]e^{ik(tt^{})}+n_\stackrel{}{k}e^{ik(tt^{})}\right],`$ (A19)
$`G_\stackrel{}{k}^<(t,t^{})`$ $`=`$ $`{\displaystyle \frac{i}{2k}}\left[n_\stackrel{}{k}e^{ik(tt^{})}+[1+n_\stackrel{}{k}]e^{ik(tt^{})}\right],`$ (A20)
$`n_\stackrel{}{k}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{exp}(\beta \omega _\stackrel{}{k})1}}.`$ (A21)
From the Lagrangian density in terms of the auxiliary fields given by (30) it is straightforward to find the free field real time correlation functions for the auxiliary fields. In terms of the spatial Fourier transform of the auxiliary field $`\chi _\stackrel{}{k}`$ these are given by
$`\chi _\stackrel{}{k}^+(t)\chi _\stackrel{}{k}^+(t^{})_o=i\delta (tt^{})`$ (A22)
(A23)
$`\chi _\stackrel{}{k}^{}(t)\chi _\stackrel{}{k}^{}(t^{})_o=i\delta (tt^{})`$ (A24)
(A25)
$`\chi _\stackrel{}{k}^+(t)\chi _\stackrel{}{k}^{}(t^{})_o=\chi _\stackrel{}{k}^{}(t)\chi _\stackrel{}{k}^+(t^{})_o=0`$ (A26)
Figures 10(a) and 10(b) depict the series of Feynman diagrams for the full $`\chi ^+\chi ^+`$ propagator and for the full plus-plus component of the propagator of the composite field $`(\stackrel{}{\mathrm{\Phi }})^2=\stackrel{}{\mathrm{\Phi }}\stackrel{}{\mathrm{\Phi }}`$, i.e. $`(\stackrel{}{\mathrm{\Phi }}^+)^2(\stackrel{}{\mathrm{\Phi }}^+)^2`$. Figures 11(a), 11(b) and 12(a), 12(b) depict similar relations for the $`\chi ^+\chi ^{}`$ and $`\chi ^{}\chi ^{}`$ propagators. Thus using the free field propagators for the auxiliary field given above we find the following relations to all orders for the full propagators
$`\chi _\stackrel{}{k}^+(t)\chi _\stackrel{}{k}^+(t^{})=i\delta (tt^{})+\lambda \mathrm{\Phi }_\stackrel{}{k}^{+,2}(t)\mathrm{\Phi }_\stackrel{}{k}^{+,2}(t^{})`$ (A27)
$`\chi _\stackrel{}{k}^+(t)\chi _\stackrel{}{k}^{}(t^{})=\lambda \mathrm{\Phi }_\stackrel{}{k}^{+,2}(t)\mathrm{\Phi }_\stackrel{}{k}^{,2}(t^{})`$ (A28)
$`\chi _\stackrel{}{k}^{}(t)\chi _\stackrel{}{k}^{}(t^{})=i\delta (tt^{})+\lambda \mathrm{\Phi }_\stackrel{}{k}^{,2}(t)\mathrm{\Phi }_\stackrel{}{k}^{,2}(t^{})`$ (A29)
with the definition
$$\mathrm{\Phi }_\stackrel{}{k}^{\pm ,2}(t)d^3xe^{i\stackrel{}{k}\stackrel{}{x}}\stackrel{}{\mathrm{\Phi }}^\pm (\stackrel{}{x},t)\stackrel{}{\mathrm{\Phi }}^\pm (\stackrel{}{x},t)$$
The correlation functions of the bilinear composite operator can be written in terms of spectral densities in the following manner
$`\mathrm{\Phi }_\stackrel{}{k}^{+,2}(t)\mathrm{\Phi }_\stackrel{}{k}^{+,2}(t^{})={\displaystyle 𝑑\omega \left[\rho _{\varphi ^2\varphi ^2}^>(\omega ;k)\theta (tt^{})+\rho _{\varphi ^2\varphi ^2}^<(\omega ;k)\theta (t^{}t)\right]e^{i\omega (tt^{})}}`$ (A30)
$`\mathrm{\Phi }_\stackrel{}{k}^{+,2}(t)\mathrm{\Phi }_\stackrel{}{k}^{,2}(t^{})={\displaystyle 𝑑\omega \rho _{\varphi ^2\varphi ^2}^<(\omega ;k)e^{i\omega (tt^{})}}`$ (A31)
Familiar manipulations introducing a complete set of energy eigenstates in the trace lead to the KMS condition
$$\rho _{\varphi ^2\varphi ^2}^<(\omega ;k)=e^{\beta \omega }\rho _{\varphi ^2\varphi ^2}^>(\omega ;k)$$
(A32)
The main reason for presenting these formal steps is that the auxiliary field itself does not have a KMS relationship for its spectral functions because it is not a canonical field but a Lagrange multiplier. However the relations (A27-A29) which hold to all orders relate the correlators of the auxiliary field to those of the bilinear composite operator for which the spectral functions associated with their correlators do obey the KMS condition.
Writing the retarded correlator for the auxiliary fields as a spectral representation
$$\chi _\stackrel{}{k}^+(t)\chi _\stackrel{}{k}^+(t^{})\chi _\stackrel{}{k}^+(t)\chi _\stackrel{}{k}^{}(t^{})=i\frac{dq_0}{2\pi }\rho _\chi (q_0,k)e^{iq_0(tt^{})},$$
using the spectral representations (A30)-(A31) and the representation for $`\mathrm{\Theta }(tt^{})`$ given by (19) we obtain the relation between the spectral representation for the retarded correlator of the auxiliary field and that for the bilinear composite in the following form
$$\rho _\chi (q_0,k)=1+\lambda 𝑑\omega \rho _{\varphi ^2\varphi ^2}^>(\omega ;k)\frac{1e^{\beta \omega }}{q_0\omega +iϵ}$$
(A33)
where we have used the KMS condition (A32). The next step of the program is to obtain the spectral density $`\rho _\chi (q_0,k)`$ to leading order in the large $`N`$ limit. This is achieved through linear response analysis for the expectation value of the auxiliary field.
### 1 Linear response for the auxiliary field
The real time expectation value of the auxiliary field is obtained by coupling an external source to the auxiliary field in the original Lagrangian (30) $`+J_\chi \chi `$ with the same external source for the two time branches. Assuming the addition of counterterms in the Lagrangian to ensure that the expectation value of the auxiliary field vanishes for vanishing external source we have
$$\delta _\stackrel{}{k}(t)\chi _\stackrel{}{k}^+(t)=i𝑑t^{}J_{\chi ,\stackrel{}{k}}(t^{})\left[\chi _\stackrel{}{k}^+(t)\chi _\stackrel{}{k}^+(t^{})\chi _\stackrel{}{k}^+(t)\chi _\stackrel{}{k}^{}(t^{})\right].$$
Introducing the Fourier transforms
$$\delta _\stackrel{}{k}(t)=\frac{dq_0}{2\pi }\delta (q_0,\stackrel{}{k})e^{iq_0t};J_{\chi ,\stackrel{}{k}}(t)=\frac{dq_0}{2\pi }J(q_0,\stackrel{}{k})e^{iq_0t}$$
we find
$$\delta (\stackrel{}{k},q_0)=J_\chi (\stackrel{}{k},q_0)\rho _\chi (k,q_0)$$
We now use the tadpole method to obtain the equation of motion for the expectation value of the auxiliary field in leading order in the large $`N`$ limit, thereby obtaining an explicit expression for $`\rho _\chi (q_0,\stackrel{}{k})`$ to this order. The implementation of the tadpole method begins by shifting the auxiliary field
$$\chi (\stackrel{}{x},t)=\delta (\stackrel{}{x},t)+\stackrel{~}{\chi }(\stackrel{}{x},t);\stackrel{~}{\chi }(\stackrel{}{x},t)=0$$
and requiring that $`\stackrel{~}{\chi }(\stackrel{}{x},t)=0`$ to all orders in perturbation theory. A counterterm is added to the Lagrangian to cancel the tadpole contributions so as to make the expectation value of the auxiliary field to vanish in the absence of the source term thus allowing to extract the spectral density straightforwardly.
To leading order in the large $`N`$ limit we obtain the equation of motion (after the cancellation of the tadpole term) to be given by
$$\delta (\stackrel{}{x},t)+d^3x^{}𝑑t^{}\mathrm{\Pi }_r(\stackrel{}{x}\stackrel{}{x}^{},tt^{})\delta (\stackrel{}{x}^{},t)=J_\chi (\stackrel{}{x},t)$$
with the retarded polarization given by
$`\mathrm{\Pi }_r(\stackrel{}{x}\stackrel{}{x}^{},tt^{})`$ $`=`$ $`2i{\displaystyle \frac{\lambda }{N}}{\displaystyle \underset{a,b}{}}[\mathrm{\Phi }_a^+(\stackrel{}{x},t)\mathrm{\Phi }_b^+(\stackrel{}{x}^{},t)\mathrm{\Phi }_a^+(\stackrel{}{x},t)\mathrm{\Phi }_b^+(\stackrel{}{x}^{},t)`$ (A35)
$`\mathrm{\Phi }_a^+(\stackrel{}{x},t)\mathrm{\Phi }_b^{}(\stackrel{}{x}^{},t)\mathrm{\Phi }_a^+(\stackrel{}{x},t)\mathrm{\Phi }_b^{}(\stackrel{}{x}^{},t)]`$
In terms of the spatial Fourier transform the equation of motion becomes
$$\delta _\stackrel{}{k}(t)+_{\mathrm{}}^{\mathrm{}}𝑑t^{}\mathrm{\Pi }_{k,r}(tt^{})\delta _\stackrel{}{k}(t^{})=J_{\chi ,\stackrel{}{k}}(t)$$
and the retarded polarization kernel simplifies to
$`\mathrm{\Pi }_{k,r}(tt^{})`$ $`=`$ $`2i\lambda {\displaystyle \frac{d^3q}{(2\pi )^3}\left[(iG_\stackrel{}{q}^>(tt^{}))(iG_{\stackrel{}{q}+\stackrel{}{k}}^>(tt^{}))(iG_\stackrel{}{q}^<(tt^{}))(iG_{\stackrel{}{q}+\stackrel{}{k}}^<(tt^{}))\right]\mathrm{\Theta }(tt^{})}`$ (A36)
$`=`$ $`4\lambda {\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}{\displaystyle \frac{1}{4q|\stackrel{}{k}+\stackrel{}{q}|}}\{(1+n_\stackrel{}{q}+n_{\stackrel{}{q}+\stackrel{}{k}})\mathrm{sin}\left[(q+|\stackrel{}{k}+\stackrel{}{q}|)(tt^{})\right]`$ (A38)
$`+(n_\stackrel{}{q}n_{\stackrel{}{q}+\stackrel{}{k}})\mathrm{sin}\left[(|\stackrel{}{k}+\stackrel{}{q}|q)(tt^{})\right]\}\mathrm{\Theta }(tt^{}),`$
using the representation of the theta function given by (19) we find the time-Fourier representation of the retarded polarization to be given by
$$\mathrm{\Pi }_{k,r}(tt^{})=\frac{dq_0}{2\pi }\mathrm{\Pi }(q_0,q)e^{iq_0(tt^{})}$$
The Fourier transform of the polarization is now written as a dispersion integral in terms of the spectral density as
$$\mathrm{\Pi }(q_0,q)=\frac{1}{\pi }𝑑\omega \frac{\mathrm{\Pi }_I(q,\omega )}{q_0\omega +iϵ}$$
where
$`\mathrm{\Pi }_I(q,\omega )=2\lambda \pi {\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{4p|\stackrel{}{p}+\stackrel{}{q}|}}`$ $`\{[1+n_{\stackrel{}{q}+\stackrel{}{p}}+n_\stackrel{}{p}][\delta (\omega |\stackrel{}{p}+\stackrel{}{q}|p)\delta (\omega +|\stackrel{}{p}+\stackrel{}{q}|+p)]`$ (A40)
$`+[n_\stackrel{}{p}n_{\stackrel{}{q}+\stackrel{}{p}}][\delta (\omega |\stackrel{}{p}+\stackrel{}{q}|+p)\delta (\omega +|\stackrel{}{p}+\stackrel{}{q}|p)]\}`$
Finally, in terms of the time Fourier transform the equation of motion for the expectation value of the auxiliary field is given by
$$\delta (q_0,\stackrel{}{q})\left[1+\mathrm{\Pi }(q_0,q)\right]=J_\chi (q_0,\stackrel{}{q})$$
and we can read off the propagators for the auxiliary field in Fourier space
$$G_\chi (q_0,q)=\frac{1}{1+\mathrm{\Pi }(q_0,q)}$$
(A41)
The series of diagrams that are being summed leading to the propagator for the auxiliary field is shown in fig. (2).
We now have all of the elements necessary to obtain $`\rho _{\varphi ^2\varphi ^2}^>(q_0,q),\rho _{\varphi ^2\varphi ^2}^<(q_0,q)`$, writing $`\mathrm{\Pi }(q_0,q)=\mathrm{\Pi }_R(q_0,q)+i\mathrm{\Pi }_I(q_0,q)`$ and comparing the imaginary parts of(A33) and (A41) and using the KMS condition (A32) we finally find
$`\lambda \rho _{\varphi ^2\varphi ^2}^>(q_0,q)={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{\Pi }_I(q_0,q)[1+n(q_0)]}{\left[1+\mathrm{\Pi }_R(q_0,q)\right]^2+\mathrm{\Pi }_I^2(q_0,q)}}`$ (A42)
$`\lambda \rho _{\varphi ^2\varphi ^2}^<(q_0,q)={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{\Pi }_I(q_0,q)n(q_0)}{\left[1+\mathrm{\Pi }_R(q_0,q)\right]^2+\mathrm{\Pi }_I^2(q_0,q)}}`$ (A43)
$`n(q_0)={\displaystyle \frac{1}{e^{\beta q_0}1}}`$ (A44)
We postpone the evaluation of $`\mathrm{\Pi }(q_0,q)`$ to Appendix B and now focus on obtaining the resummed self energy for the order parameter.
### 2 Equation of motion for the order parameter
We now obtain the equation of motion for the order parameter to $`𝒪(1/N)`$ in the linearized approximation again via the tadpole method and recognize the self-energy to this order. To this effect we write the field as in (5) with
$$\mathrm{\Phi }^i(\stackrel{}{x},t)=\phi (\stackrel{}{x},t)\delta _{i,1};\eta ^i(\stackrel{}{x},t)=0$$
where we chose the particular direction “1” by choosing explicitly the external source in (4) as $`J^i(\stackrel{}{x},t)=J(\stackrel{}{x},t)\delta ^{i,1}`$ to give the field an expectation value solely in this direction. The equation of motion for $`\phi (\stackrel{}{x},t)`$ is obtained by imposing that $`\eta ^i(\stackrel{}{x},t)=0`$ consistently in the perturbative expansion. In terms of the spatial Fourier transform of the order parameter $`\phi (t)`$ we find
$$\ddot{\phi }_k(t)+[k^2+\delta M^2(T)+M_{tad}^2(T)]\phi _k(t)+_{\mathrm{}}^{\mathrm{}}\mathrm{\Sigma }_{ret,k}(tt^{})\phi _k(t^{})𝑑t^{}=J_k(t)$$
where $`J_k(t)`$ is the external source that generates the initial value problem and $`M_{tad}^2(T)(\stackrel{}{\mathrm{\Phi }}(\stackrel{}{x},t)^2𝒪(N)`$ is the tadpole contribution which is the leading order in the large $`N`$ limit. The $`𝒪(1/N)`$ contribution to the self-energy is calculated in terms of the auxiliary field and is given by
$$\mathrm{\Sigma }_{ret,k}(tt^{})=\frac{4i\lambda }{N}\frac{d^3q}{(2\pi )^3}\left[(iG_{\stackrel{}{k}+\stackrel{}{q}}^{++}(tt^{}))\chi _\stackrel{}{q}^+(t)\chi _\stackrel{}{q}^+(t^{})(iG_{\stackrel{}{k}+\stackrel{}{q}}^+(tt^{}))\chi _\stackrel{}{q}^+(t)\chi _\stackrel{}{q}^{}(t^{})\right]$$
with $`\chi _\stackrel{}{q}^+(t)\chi _\stackrel{}{q}^+(t^{})`$ and $`\chi _\stackrel{}{q}^+(t)\chi _\stackrel{}{q}^{}(t^{})`$ the full propagators up to $`𝒪(1/N)`$ given by (A27-A28) in terms of the spectral representations given by (A30-A31) with the spectral densities given in terms of the self-energy of the auxiliary field by (A42-A43). The contribution to the propagator of the scalar field up to order $`𝒪(1/N)`$ is depicted in fig. 4. The contribution to the auxiliary field propagators from the delta functions $`\pm i\delta (tt^{})`$ gives a local tadpole which is cancelled along with the leading order $`𝒪(1)`$ tadpole contribution by the counterterm to set the theory at the critical point up to this order in the large $`N`$ expansion. Using the spectral representation for the propagators of the auxiliary field and the free field propagators for the bosonic fields given by (A19)-(A20) and after some straightforward algebra using the relation $`1+n(q_0)=n(q_0)`$ we finally obtain
$$\mathrm{\Sigma }_{ret,k}(tt^{})=_{\mathrm{}}^{\mathrm{}}\stackrel{~}{\rho }(\omega ,k)\mathrm{sin}\left[\omega (tt^{})\right]𝑑\omega $$
with $`\stackrel{~}{\rho }(\omega ,k)`$ given by eqs.(37)-(41).
## B The retarded polarization of the auxiliary field
The spectral density (A40) is the same as (18) up to the factor $`(N+2)/N`$, a relatively straightforward calculation with the Bose-Einstein distribution functions for massless particles then leads to
$$\mathrm{\Pi }_I(q_0,q)=\frac{\lambda }{8\pi }\left\{\mathrm{\Theta }(|q_0|q)\text{sign}(q_0)+\frac{2T}{q}\mathrm{ln}\left[\frac{1e^{\frac{|q_0+q|}{2T}}}{1e^{\frac{|q_0q|}{2T}}}\right]\right\}$$
(B1)
The real part must be obtained via the dispersive integral (41). We are only interested in the finite temperature contribution to both the real and imaginary part therefore we only consider the second term in (B1). It proves convenient to write the polarization as a dispersion relation
$$\mathrm{\Pi }(q_0i0,q)=\frac{1}{\pi }_{\mathrm{}}^+\mathrm{}𝑑\omega \frac{\mathrm{\Pi }_I(\omega ,q)}{\omega q_0+i0}$$
and to analytically continue $`0+iq_0=s`$. Using the fact that $`\mathrm{\Pi }_I(\omega ,q)`$ is an odd function of $`\omega `$ we obtain the dispersion relation
$$\mathrm{\Pi }(s,q)=\frac{\lambda T}{2\pi ^2q}_0^{\mathrm{}}𝑑\omega \frac{\omega }{\omega ^2+s^2}\mathrm{ln}\left[\frac{1e^{\frac{|\omega +q|}{2T}}}{1e^{\frac{|\omega q|}{2T}}}\right]\frac{\lambda }{8\pi ^2}\mathrm{ln}\left[\frac{q^2+s^2}{\mu ^2}\right].$$
(B2)
where $`\mu ^2`$ is a subtraction point.
We compute this integral using the sine-Fourier transform as follows. The integrand of eq.(B2) is the product of two odd functions of $`\omega `$ :
$$f_1(\omega )=\frac{\omega }{\omega ^2+s^2}\text{and}f_2(\omega )=\mathrm{ln}\left[\frac{1e^{\frac{|\omega +q|}{2T}}}{1e^{\frac{|\omega q|}{2T}}}\right].$$
We can then apply the Plancherel formula
$$_0^{\mathrm{}}𝑑\omega f_1(\omega )f_2(\omega )=_0^{\mathrm{}}𝑑x\stackrel{~}{f}_1(x)\stackrel{~}{f}_2(x)$$
where $`\stackrel{~}{f}_1(x)`$ and $`\stackrel{~}{f}_2(x)`$ are the sine-Fourier transforms of $`f_1(\omega )`$ and $`f_2(\omega )`$, respectively. That is,
$$\stackrel{~}{f}_i(x)=\sqrt{\frac{2}{\pi }}_0^{\mathrm{}}𝑑\omega f_i(\omega )\mathrm{sin}\omega x$$
where $`i=1,2`$. We find
$$\stackrel{~}{f}_1(x)=\sqrt{\frac{\pi }{2}}e^{sx}$$
and
$$\stackrel{~}{f}_2(x)=\sqrt{\frac{2}{\pi }}\frac{\mathrm{sin}qx}{2Tx^2}\left[2\pi Tx\mathrm{coth}\left(2\pi Tx\right)1\right].$$
We have now that
$$\mathrm{\Pi }(s,q)=\frac{\lambda }{(2\pi )^2q}_0^{\mathrm{}}𝑑x\frac{e^{sx}}{x^2}\mathrm{sin}qx\left[2\pi Tx\mathrm{coth}\left(2\pi Tx\right)1\right]\frac{\lambda }{8\pi ^2}\mathrm{ln}\frac{q^2+s^2}{\mu ^2}.$$
It is convenient to split this integral into two terms,
$`\mathrm{\Pi }(s,q)`$ $`=`$ $`{\displaystyle \frac{\lambda }{(2\pi )^2q}}\{{\displaystyle _0^{\mathrm{}}}dx{\displaystyle \frac{e^{sx}}{x}}[1{\displaystyle \frac{\mathrm{sin}qx}{x}}]`$ (B3)
$`+`$ $`{\displaystyle _0^{\mathrm{}}}dx{\displaystyle \frac{e^{sx}}{x}}[2\pi T\mathrm{coth}\left(2\pi Tx\right)\mathrm{sin}qx1]\}{\displaystyle \frac{\lambda }{8\pi ^2}}\mathrm{ln}{\displaystyle \frac{q^2+s^2}{\mu ^2}}.`$ (B5)
We carried out the integration explicitly with the result
$$\mathrm{\Pi }(s,q)=\frac{\lambda }{(2\pi )^2}\left\{1+\mathrm{ln}\frac{\mu }{4\pi T}+\frac{s}{q}\text{arctg}\frac{q}{s}\frac{i\pi T}{q}\mathrm{ln}\left[\frac{\mathrm{\Gamma }\left(\frac{is+q}{4\pi iT}\right)\mathrm{\Gamma }\left(1+\frac{is+q}{4\pi iT}\right)}{\mathrm{\Gamma }\left(\frac{isq}{4\pi iT}\right)\mathrm{\Gamma }\left(1+\frac{isq}{4\pi iT}\right)}\right]\right\}$$
where we used Malmsten formula for the Gamma functions.
Back in real frequencies we have,
$$\mathrm{\Pi }(q_0\pm i0,q)=\mathrm{\Pi }_R(q_0,q)\pm i\mathrm{\Pi }_I(q_0,q)$$
where $`\mathrm{\Pi }_I(q_0,q)`$ is given by eq.(B1) and
$`\mathrm{\Pi }_R(q_0,q)`$ $`=`$ $`{\displaystyle \frac{\lambda }{(2\pi )^2}}\{{\displaystyle \frac{\pi ^2T}{q}}[\theta (qq_0)\theta (qq_0)]+\mathrm{ln}{\displaystyle \frac{\mu }{4\pi T}}`$ (B6)
$`+`$ $`{\displaystyle \frac{q_0}{2q}}\mathrm{ln}\left|{\displaystyle \frac{q+q_0}{qq_0}}\right|+{\displaystyle \frac{2\pi T}{q}}\text{Im}\mathrm{ln}\left[\mathrm{\Gamma }(1+{\displaystyle \frac{qq_0}{4\pi iT}})\mathrm{\Gamma }(1+{\displaystyle \frac{q+q_0}{4\pi iT}})\right]\}.`$ (B8)
The limit $`T/q>1`$ can be taken in a straightforward manner and we obtain the high temperature limit of the polarization to be given by
$`\mathrm{\Pi }(q_0,q+iϵ)`$ $`=`$ $`{\displaystyle \frac{i\lambda T}{4\pi q}}\mathrm{ln}\left[{\displaystyle \frac{q_0+iϵ+q}{q_0+iϵq}}\right]`$ (B9)
$`+`$ $`{\displaystyle \frac{\lambda }{(2\pi )^2}}\left[\mathrm{ln}{\displaystyle \frac{\mu }{4\pi T}}+{\displaystyle \frac{q_0}{2q}}\mathrm{ln}\left({\displaystyle \frac{q_0+iϵ+q}{q_0+iϵq}}\right)+2\gamma \right]+𝒪\left({\displaystyle \frac{1}{T}}\right)`$ (B10)
where $`\gamma `$ is the Euler-Mascheroni constant.
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# COLOR SUPERCONDUCTIVITY IN COLD, DENSE QUARK MATTER
## 1 Cold, dense fermionic matter
Consider a degenerate, non-interacting fermionic system, where all momentum states up to the Fermi surface are occupied. The fermion occupation number is a step function, $`n_0(ϵ_0)=\theta (ϵ_0)`$, where $`ϵ_0\omega \mu `$ is the energy of the fermions with respect to the Fermi energy, $`\mu `$, and $`\omega `$ is their kinetic energy. For non-relativistic particles with 3-momentum $`𝐤`$ and mass $`m`$, $`\omega =𝐤^2/(2m)`$, while for massless (ultrarelativistic) particles, $`\omega =k|𝐤|`$. (The units are $`\mathrm{}=k_B=c=1`$.) The excitation spectrum of the fermions consists of a particle branch, $`ϵ_0^\mathrm{p}ϵ_0`$, and a hole branch, $`ϵ_0^\mathrm{h}ϵ_0`$. The formation of particle-hole excitations at the Fermi surface costs no energy, $`ϵ_0^\mathrm{p}+ϵ_0^\mathrm{h}0`$.
Now switch on an interaction, for the sake of simplicity a point-like four-fermion interaction with interaction strength $`G^2`$. Let the sign of $`G^2`$ be defined such that $`G^2>0`$ in the case of an attractive interaction, and $`G^2<0`$, if the interaction is repulsive. Due to Pauli’s exclusion principle and energy conservation, scattering of fermions can occur exclusively at the Fermi surface. In other words, physical scattering processes are possible only for $`ϵ_00`$. The amplitude for fermion-fermion scattering is $`^\mathrm{?}`$
$$\mathrm{\Gamma }(ϵ_0)\frac{G^2}{1G^2\mathrm{ln}(\mu /ϵ_0)},$$
(1)
where we suppress all constant factors that are irrelevant for the qualitative arguments presented here. As $`ϵ_00`$, in the case of an attractive interaction the scattering amplitude develops a singularity at an energy scale
$$ϵ_0^{}\mu \mathrm{exp}(1/G^2).$$
(2)
For repulsive interactions, no such singularity occurs. Obviously, the singularity (2) occurs even when the attractive interaction is arbitrarily weak, $`G^20^+`$, all that changes is the energy scale $`ϵ_0^{}`$ of the singularity.
This singularity is the famous Cooper instability $`^\mathrm{?}`$. It is cured by the formation of Cooper pairs which, as bosons, condense in the true, energetically favored ground state of the system. Macroscopically, this leads to the phenomenon of superfluidity, or, if the fermions carry charge, superconductivity.
The excitation spectrum of the system also changes. A gap $`\varphi _0`$ forms at the Fermi surface, separating the branch for quasiparticle excitations, $`ϵ^\mathrm{p}=ϵ`$, $`ϵ\sqrt{ϵ_0^2+\varphi _0^2}`$, from the one for quasihole excitations, $`ϵ^\mathrm{h}+ϵ`$. Now exciting a quasiparticle–quasihole pair costs at least an energy $`ϵ^\mathrm{p}+ϵ^\mathrm{h}2\varphi _0`$. The quasiparticle occupation number, $`n(ϵ_0)`$, is “smeared” around the (original) Fermi surface, $`n(ϵ_0)=(ϵϵ_0)/(2ϵ)`$.
To compute the gap, one has to solve a gap equation $`^\mathrm{?}`$. In the case of a point-like four-fermion interaction, this equation takes the form
$$\varphi _0=G^2_\mu ^0\frac{\mathrm{d}ϵ_0}{ϵ}\varphi _0.$$
(3)
(Again, irrelevant constant factors are omitted.) Besides the trivial (energetically disfavored) solution $`\varphi _0=0`$, this equation has always a non-trivial (energetically favored) solution, $`\varphi _00`$. Since the gap is a constant for point-like four-fermion interactions, we can divide both sides of Eq. (3) by $`\varphi _0`$. The remaining integral can be solved exactly, with the result
$$\varphi _02\mu \mathrm{exp}(1/G^2).$$
(4)
Apparently, $`\varphi _0>ϵ_0^{}`$, cf. Eq. (2). In other words, since there are no quasiparticle states with energy $`ϵ^\mathrm{p}\varphi _0`$ (or quasihole states with $`ϵ^\mathrm{h}\varphi _0`$), the gap has just the right order of magnitude to prevent the scattering amplitude (1) from developing a singularity.
The assumption of a point-like four-fermion interaction can be relaxed. Assume that the interaction is mediated by a scalar boson of mass $`M`$. For the sake of definiteness, assume that the boson mass is generated by many-body interactions at nonzero density, $`Mg\mu `$. The boson-fermion coupling is denoted by $`g`$, the boson propagator is $`\mathrm{\Delta }(P)=1/(M^2P^2)`$, $`PP^\mu =(p_0,𝐩)`$, $`P^2p_0^2𝐩^2`$. In this case, the gap equation (3) becomes $`^\mathrm{?}`$
$$\varphi _0(k)=g^2_\mu ^0\frac{\mathrm{d}ϵ_0}{ϵ}\varphi _0(q)\frac{q}{k}\mathrm{ln}\left[\frac{M^2+(k+q)^2}{M^2+(kq)^2}\right].$$
(5)
Here, $`qϵ_0+\mu `$. The frequency dependence of the boson propagator and, consequently, that of the gap function has been neglected, based on the argument that in weak coupling, $`g1`$, $`p_0ϵ\varphi _0\mu \mathrm{exp}(1/g^2)Mg\mu `$. The logarithm arises from the integration over the angle between the boson 3-momentum $`𝐩𝐤𝐪`$ and the 3-momentum of the fermion in the condensate, $`𝐪`$, for details see Ref. $`^\mathrm{?}`$. Note that this factor enhances contributions from the region of momenta $`𝐪𝐤`$ (collinear enhancement).
In the case of a boson-mediated interaction, the gap is no longer constant, but a function of momentum. Consequently, the gap equation (5) is no longer a simple fix-point equation for $`\varphi _0`$, but an integral equation which has to be solved numerically. However, to estimate the order of magnitude of the gap at the Fermi surface, $`k\mu `$ (for massless fermions), one may make the following approximation. First note that, due to the factor $`1/ϵ`$, the integrand peaks at the Fermi surface. It is then sufficient to approximate the slowly varying logarithm and the gap function with their values for $`q=k=\mu `$. This leads to the estimate
$$\varphi _02\mu \mathrm{exp}\left[\frac{1}{g^2\mathrm{ln}(1+2\mu ^2/M^2)}\right]$$
(6)
for the value of the gap function at the Fermi surface, $`\varphi _0\varphi _0(\mu )`$, which should be compared with Eq. (4). All that changed is that the coupling constant is effectively increased by the logarithm originating from collinear enhancement. For $`Mg\mu `$, the logarithm becomes $`\mathrm{ln}(1/g)`$ in weak coupling.
## 2 Cold, dense quark matter
The density in cold quark matter increases $`\mu ^3`$. Asymptotic freedom then implies that single-gluon exchange becomes the dominant interaction between quarks. Single-gluon exchange is attractive in the color-antitriplet channel, and therefore leads to color superconductivity in cold, dense quark matter $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. Recently, considerable activity was generated <sup>6-39</sup> by the work of Refs. $`^{\mathrm{?},\mathrm{?}}`$, which suggested that the zero-temperature color-superconducting gap $`\varphi _0`$ could be as large as 100 MeV. This order of magnitude is quite surprising, because earlier work by Bailin and Love $`^\mathrm{?}`$ estimated the gap to be $`\varphi _01`$ MeV. While gaps of order 100 MeV could also be relevant for the physics of nuclear collisions (see below), gaps of about 1 MeV allow at most for the possibility that neutron star cores, if consisting of quark matter, could be color superconductors.
The authors of Refs. $`^{\mathrm{?},\mathrm{?}}`$ based their arguments on a simple model where quarks interact via a point-like four-fermion interaction, giving rise to a gap equation of the type (3). The coupling strength $`G^2`$ was adjusted such that the model reproduced the order of magnitude of the chiral transition at nonzero temperature and zero quark chemical potential, $`T_\chi 150`$ MeV. The earlier work of Bailin and Love $`^\mathrm{?}`$ was already more sophisticated in the sense that they used one-gluon exchange, employing gluon propagators with electric and magnetic screening masses. This gives rise to a gap equation of the form (5), with $`g`$ being the strong coupling constant.
Unfortunately, both approaches fail to capture an essential property of single-gluon exchange: at zero temperature, due to the absence of magnetic screening, magnetic interactions are truly long-range. Surprisingly, this fact was already known to Barrois $`^\mathrm{?}`$, but apparently never made it into the published literature.
Long-range magnetic interactions have the important consequence that one can no longer neglect the frequency dependence of the boson propagator, as done in the derivation of (5). For massless boson exchange, the gap equation assumes the (approximate) form
$$\varphi _0(k_0)g^2_\mu ^0\frac{\mathrm{d}ϵ_0}{ϵ}\varphi _0(ϵ)\frac{1}{2}\mathrm{ln}\left(\frac{\mu ^2}{ϵ^2k_0^2}\right).$$
(7)
Again, as in (5) there is an additional logarithm, representing collinear enhancement. In QCD it arises from the exchange of ultrasoft, magnetic gluons. While the collinear enhancement in Eq. (5) is cut off by the mass of the scalar boson, $`M`$, here it is cut off by the energy of the magnetic gluon, $`p_0=ϵ\pm k_0`$. For $`ϵk_0\varphi _0`$, $`p_0`$ is of the order of $`\varphi _0`$, too, while in weak coupling, $`Mg\mu \varphi _0`$. For magnetic gluon exchange in QCD, the contribution of the collinear region, $`𝐪𝐤`$, to the gap integral is therefore much larger than in the case of massive boson exchange.
To estimate the effect of the logarithm on the solution of the gap equation, let us neglect the energy dependence of the gap function in the integrand and consider its value at the Fermi surface, $`k=\mu `$, $`k_0=\varphi _0`$. One may also make the approximation $`^\mathrm{?}`$ $`\mathrm{ln}[\mu ^2/(ϵ^2\varphi _0^2)]2\mathrm{ln}(\mu /ϵ)`$. Then, the integral is again exactly solvable, with the (order of magnitude) result
$$\varphi _02\mu \mathrm{exp}(1/g).$$
(8)
Due to the explicit $`ϵ`$ dependence of the logarithm in (7), the power of $`g`$ in the exponent is reduced as compared to the BCS result (4). The case of massive scalar boson exchange, Eq. (6), interpolates between these two cases, as $`g^2g^2\mathrm{ln}(1/g)g`$ for $`g1`$. This reflects the fact mentioned earlier that, while the collinear singularity $`𝐪=𝐤`$ in (5) is cut off by the mass of the scalar boson, $`Mg\mu `$, the singularity in (7) is cut off by the gluon energy, $`p_0\varphi _0`$, which is, in weak coupling, much smaller than $`M`$. In the literature, the parametric dependence on $`g`$ of the solution (8) to the gap equation (7) was first discussed in Refs. $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$.
As of today, various refinements of the solution (8) have been discussed $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$. The value of the gap function at the Fermi surface, $`k=\mu `$, $`k_0=\varphi _0`$, is
$$\varphi _0=b\mu g^5\mathrm{exp}(\frac{c}{g})[1+O(g)],c=\frac{3\pi ^2}{\sqrt{2}}.$$
(9)
Furthermore, the gap function has a non-trivial (on-shell) energy dependence,
$$\varphi _0(ϵ)=\varphi _0\mathrm{sin}\left[\frac{\pi }{2}\frac{g}{c}\mathrm{ln}\left(\frac{b\mu }{f(ϵ)}\right)\right].$$
(10)
The constant $`c`$ was first computed by Son $`^\mathrm{?}`$. To obtain the correct numerical value for $`c`$, one has to account for the modifications of the gluon propagator in the presence of a dense medium $`^\mathrm{?}`$. Then, what dominates the gap equation (7) and determines $`c`$ is the contribution from nearly static, Landau-damped magnetic gluons. Furthermore, Son showed that in computing $`c`$, it is essential to retain the energy dependence (10) of the gap function. To the level of accuracy considered by Son $`^\mathrm{?}`$, $`b=1`$ and $`f(ϵ)ϵ`$. The prefactor $`g^5`$ arises from subleading contributions of static electric and non-static magnetic gluons.
The constant $`b`$ collects constant factors in these subleading contributions. It was first computed by Schäfer and Wilczek $`^\mathrm{?}`$ and the present authors $`^\mathrm{?}`$,
$$b=512\pi ^4\left(\frac{2}{N_f}\right)^{5/2}b^{},$$
(11)
with an undetermined constant $`b^{}`$. Schäfer and Wilczek $`^\mathrm{?}`$ obtained $`b^{}=1/2`$ and, as before, $`f(ϵ)ϵ`$. We showed $`^\mathrm{?}`$ that actually $`b^{}=1`$ and $`f(ϵ)ϵϵ_0`$. The additional factor 2 is the same that occurs in Eqs. (4), (6), and (8). It arises from the measure of integration in the gap equation, $`\mathrm{d}ϵ_0/ϵ\mathrm{d}\mathrm{ln}(ϵϵ_0)`$, and $`\mathrm{ln}(ϵϵ_0)|_\mu ^0\mathrm{ln}(2\mu /\varphi _0)`$. This also explains the modification of $`f(ϵ)`$. In addition, we pointed out $`^\mathrm{?}`$ that it is important to compute the gap function on the correct quasiparticle mass shell, and we considered for the first time the case of non-zero temperature. As will be discussed in more detail below, the temperature $`T_c`$ where the color-superconducting condensate melts is related to the zero-temperature gap in the same way as in BCS theory, $`T_c/\varphi _0=e^\gamma /\pi 0.567`$, where $`\gamma 0.577`$ is the Euler–Mascheroni constant.
Brown, Liu, and Ren $`^\mathrm{?}`$ computed $`T_c`$ from the quark-quark scattering amplitude. Using the result $`T_c/\varphi _0=e^\gamma /\pi `$ of Ref. $`^\mathrm{?}`$, they obtain
$$b^{}=\mathrm{exp}\left(\frac{\pi ^2+4}{8}\right)0.176,$$
(12)
where the correction to the previous estimate $`b^{}=1`$ arises from a finite, $`\mu `$ dependent contribution to the wavefunction renormalization for quarks in a dense medium. The authors of Ref. $`^\mathrm{?}`$ also assert that there are no further corrections to $`b^{}`$ at this order in $`g`$. If correct, this is a remarkable result, because from previous calculations $`^{\mathrm{?},\mathrm{?}}`$ it appeared that computing $`b^{}`$ exactly to leading order in $`g`$ would be a formidable task.
Unlike calculations of the free energy or the Debye mass, where the perturbative expansion in powers of $`g^2`$ appears to be well-behaved $`^\mathrm{?}`$ even when extrapolated down to moderate values of $`\mu `$, this result for the wavefunction renormalization (which is equivalent to non-Fermi liquid behavior) indicates that perturbation theory is not such a good approximation, at least for quarks near the Fermi surface. The confirmation of the results of Ref. $`^\mathrm{?}`$ is clearly an outstanding problem for the field.
Figure 1 shows the value of the zero-temperature gap $`\varphi _0`$ at the Fermi surface according to Eq. (9) as a function of $`\mu `$ for $`N_f=2`$ and $`3`$ massless quark flavors with $`b^{}=1`$ (full and dotted lines), and for $`N_f=3`$ flavors with $`b^{}=\mathrm{exp}[(\pi ^2+4)/8]`$ (dashed line). The running of the coupling $`g(\mu )`$ with the chemical potential $`\mu `$ was computed from the 3-loop QCD $`\beta `$ function $`^\mathrm{?}`$, however, not for 6 but only for 3 flavors of massless quarks. Therefore, the QCD scale $`\mathrm{\Lambda }=364`$ MeV is chosen somewhat larger than the standard value, to give the value $`\alpha _s(\mu =2\mathrm{GeV})0.309`$. Although an extrapolation of the weak-coupling result (9) to large $`g`$ (small $`\mu `$) appears audacious, it is interesting to note that the maximum value of the gap for $`b^{}=1`$ is of the order of 100 MeV, quite in agreement with the earlier estimates of Refs. $`^{\mathrm{?},\mathrm{?}}`$. However, taking into account the quark wavefunction renormalization $`^\mathrm{?}`$ reduces the gap to values of a few MeV. These values are of the order of typical superfluid gaps in ordinary hadronic matter. This lends credibility to the conjecture that quark and hadronic matter are continuously connected $`^\mathrm{?}`$, although symmetry arguments $`^{\mathrm{?},\mathrm{?}}`$ suggest that, at zero temperature, there is a first order phase transition between these two phases of nuclear matter.
## 3 Not so cold, dense quark matter
To understand how the color-superconducting gap changes with temperature, it is instructive to first consider the simpler BCS case. At nonzero temperature $`T`$, the gap equation (3) becomes
$$\varphi =G^2_\mu ^0\frac{\mathrm{d}ϵ_0}{ϵ(\varphi )}\varphi \mathrm{tanh}\left[\frac{ϵ(\varphi )}{2T}\right].$$
(13)
Here, $`\varphi `$ is the value of the gap at temperature $`T`$, $`\varphi \varphi (T)`$, and as before, $`\varphi _0\varphi (0)`$ denotes the zero-temperature gap in the following. The $`\varphi `$ dependence of the quasiparticle excitation energy $`ϵ(\varphi )\sqrt{ϵ_0^2+\varphi ^2}`$ has been made explicit, to distinguish $`ϵ(\varphi )`$ from $`ϵϵ(\varphi _0)`$ used previously.
Again, $`\varphi `$ is constant and can be divided out on both sides of Eq. (13), with the result
$$1=G^2_\mu ^0\frac{\mathrm{d}ϵ_0}{ϵ(\varphi )}\mathrm{tanh}\left[\frac{ϵ(\varphi )}{2T}\right].$$
(14)
The effect of the $`\mathrm{tanh}`$ is to reduce the value of the integrand, such that $`\varphi `$ in the factor $`1/ϵ(\varphi )`$ has to decrease in order to balance the 1 on the left-hand side. At some critical temperature $`T_c`$, this balance can no longer be achieved and $`\varphi =0`$ is the only solution of Eq. (13). $`T_c`$ is the temperature where the superconducting condensate melts. Physically, the random thermal energy of the fermions exceeds their binding energy in a Cooper pair. Thus, $`T_c`$ must be of the same order as $`\varphi _0`$.
There is an easy way to compute the change of $`\varphi `$ with $`T`$, which was to our knowledge first suggested in Ref. $`^\mathrm{?}`$. Note that $`\mathrm{tanh}[ϵ(\varphi )/2T]1`$ far from the Fermi surface, $`ϵ(\varphi )|ϵ_0|\varphi _0T`$. A nonzero temperature influences the integrand only in the region close to the Fermi surface, $`|ϵ_0|\varphi _0T`$. Let us therefore divide the range of integration into two parts, $`0ϵ_0\kappa \varphi _0`$, and $`\kappa \varphi _0ϵ_0\mu `$, where $`\kappa 1`$. Then, the $`\mathrm{tanh}`$ need only be kept in the first region, and the integral over the second region can be performed similarly as at zero temperature,
$$1G^2_{\kappa \varphi _0}^0\frac{\mathrm{d}ϵ_0}{ϵ(\varphi )}\mathrm{tanh}\left[\frac{ϵ(\varphi )}{2T}\right]+G^2\mathrm{ln}\left(\frac{\mu }{\kappa \varphi _0}\right).$$
(15)
Using the solution (4) for the zero-temperature gap $`\varphi _0`$, the second term becomes $`1G^2\mathrm{ln}(2\kappa )`$. Apparently, the 1 on the left-hand side is almost completely saturated by this second term. Cancelling the 1 and writing $`\mathrm{ln}(2\kappa )_{\kappa \varphi _0}^0dϵ_0/ϵ(\varphi _0)`$, one obtains the condition
$$G^2_{\kappa \varphi _0}^0dϵ_0\left\{\frac{1}{ϵ(\varphi )}\mathrm{tanh}\left[\frac{ϵ(\varphi )}{2T}\right]\frac{1}{ϵ(\varphi _0)}\right\}=0.$$
(16)
The dependence on $`\kappa `$ is spurious: one might as well send $`\kappa \mathrm{}`$ $`^\mathrm{?}`$, because the integrand vanishes when $`\kappa 1`$. Equation (16) determines $`\varphi (T)/\varphi _0`$ as a function of $`T`$. In particular, at $`T_c`$, where $`\varphi =0`$, one derives the well-known result
$$\frac{T_c}{\varphi _0}=\frac{e^\gamma }{\pi }0.567$$
(17)
mentioned above.
In QCD, it turns out $`^\mathrm{?}`$ that the effect of temperature on the gap equation is essentially identical to that in BCS theory: the integrand in (7) is multiplied with $`\mathrm{tanh}[ϵ(\varphi )/2T]`$. For the same reasons as in the BCS case, this factor is negligible far away from the Fermi surface. One may again divide the range of integration into two parts, and neglect the effects of temperature in the one far from the Fermi surface. The integral over this region can be computed as for $`T=0`$. Quite similarly to the treatment in the BCS case, it is found to saturate the left-hand side of the gap equation up to corrections of order $`O(g)`$.
Therefore, the integral over the region close to the Fermi surface must also be of order $`O(g)`$, in order to cancel these corrections. To see this, it is permissible to compute this integral to leading order in $`g`$. One then derives the same condition (16) as in the BCS case, except that $`G^2`$ is replaced by $`g`$, see Ref. $`^\mathrm{?}`$ for details. Consequently, to leading order in $`g`$, the $`T`$ dependence of the gap at the Fermi surface, normalized to the zero-temperature gap, $`\varphi (T)/\varphi _0`$, is the same as in BCS theory. In particular, the ratio $`T_c/\varphi _0`$ is again given by Eq. (17). In retrospect, this is not surprising: the prefactor $`b`$ of the zero-temperature gap, Eq. (11), was seen to be determined by subleading terms in the gap equation \[terms of order $`O(g)`$ relative to the leading terms due to Landau-damped magnetic gluons\]. As explained above, temperature affects the gap equation at the same subleading order.
An immediate consequence is that when multiplying the ordinate of Fig. 1 with $`0.567`$, one obtains the location of the phase transition to the color-superconducting phase in the $`T\mu `$ phase diagram of nuclear matter. In the 2-flavor case without wavefunction corrections, the transition temperature is of order $`100`$ MeV. The color-superconducting phase could then be accessible in heavy-ion collisions at BNL–AGS or GSI–SIS energies, which explore the range of moderate temperatures and high (net) baryon density in the nuclear matter phase diagram. However, in the 3-flavor case, including the effects of the quark wavefunction renormalization, the transition temperature is at most $`6`$ MeV. For such small temperatures, color superconductivity occurs at best in neutron star cores, if they consist of quark matter.
## Acknowledgements
D.H.R. would like to thank the organizers of QCD 2000 for the invitation to give this talk. The authors thank T. Schäfer and D.T. Son for many enlightening discussions. D.H.R. thanks RIKEN, BNL and the U.S. Dept. of Energy for providing the facilities essential for the completion of this work, and to Columbia University’s Nuclear Theory Group for continuing access to their computing facilities.
## References
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# NT@UW-00-07,KRL MAP-266 Chiral Three-Nucleon Forces from 𝑝-wave Pion Production
## Abstract
Production of $`p`$-wave pions in nucleon-nucleon collisions is studied according to an improved power counting that embodies the constraints of chiral symmetry. Contributions from the first two non-vanishing orders are calculated. We find reasonable convergence and agreement with data for a spin-triplet cross section in $`pppp\pi ^0`$, with no free parameters. Agreement with existing data for a spin-singlet cross section in $`pppn\pi ^+`$ constrains a short-range operator shown recently to contribute significantly to the three-nucleon potential.
The use of (approximate) chiral symmetry of QCD to determine the form of the low-energy effective Lagrangian has proven to be a powerful aid to the understanding of strong interaction physics . It has long been known that the use of chiral symmetry for pion-nucleon ($`\pi N`$) scattering leads to a qualitative understanding of the pion-ranged part of the three-nucleon force, believed to produce important effects in nucleon-deuteron ($`Nd`$) scattering and few-nucleon bound states . Yet, discrepancies between theory and experiment (for $`A_y`$ in $`Nd`$ scattering and for excited levels in bound states ) remain which have been widely attributed to unknown three-nucleon forces. A novel three-nucleon force, expected on the basis of power counting arguments, involves the exchange of a pion between one nucleon and two others interacting via short-ranged forces . This force can indeed affect $`Nd`$ scattering at a currently observable level, and thus potentially resolve the remaining discrepancies . It depends on a pion-two-nucleon interaction of a form determined by chiral symmetry, but strength determined by parameters, $`d_i`$ of Eq. (7), not fixed by symmetry. We argue here that the production of $`p`$-wave pions in nucleon-nucleon ($`NN`$) scattering offers a unique opportunity to determine $`d_i`$ .
In the last few years, the various $`NNNN\pi `$ reactions have been studied both experimentally and theoretically , with a focus on near-threshold energies. The first high-quality data concerned the total cross section, and most theoretical analyses have concentrated on $`\eta `$ $`\stackrel{<}{}`$ $`0.4`$, a region dominated by the $`Ss`$ state. (Final states are labeled by $`Ll`$ with $`L`$ and $`l`$ being the relative angular momentum of the nucleon pair and the pion with respect to the two-nucleon center of mass, respectively; $`\eta `$ is the maximum pion momentum in units of the pion mass, $`m_\pi `$). Many different mechanisms are expected at these kinematics: heavy meson exchanges , (off-shell) pion rescattering , excitations of baryon resonances , and pion emission from exchanged mesons .
The pion dynamics are largely controlled by chiral symmetry constraints, and the hope that the use of Chiral Perturbation Theory ($`\chi `$PT) would yield insights led to the use of tree-level $`\chi `$PT to calculate the cross sections close to threshold . Ref. emphasized that the diverse contributions to the $`Ss`$ final states can be ordered in powers of $`\sqrt{m_\pi /M_{QCD}}`$, where $`M_{QCD}1`$ GeV is the typical QCD mass scale. The implication of this relatively large parameter is that loop diagrams enter at next-to-leading order in $`s`$-wave pion production. Thus a test of the convergence of the series is hindered. We shall show that such difficulties are not present for the case of $`\pi `$ production in $`p`$-waves ($`\eta 1`$), because the production proceeds through leading-order operators better determined from other processes.
Our arguments rely on the use of symmetries. One may obtain the results of QCD by using the most general Lagrangian involving the low-energy degrees of freedom (pion $`𝝅`$, nucleon $`N`$, and delta isobar $`\mathrm{\Delta }`$) which has the same symmetries as QCD. These are approximate chiral symmetry, parity and time-reversal invariance. Chiral symmetry plays a crucial role in low-energy processes because it demands that, in the chiral limit where the quark masses go to zero, the pion interactions contain derivatives, which are weak at small momenta, $`Q`$. Because the quark masses are small, any non-derivative pion interactions are also weak. Although the nucleon mass $`m_N`$ is not small, it plays no dynamical role at low energies. The delta isobar can be excited, but its mass difference to the nucleon, $`\delta mm_\mathrm{\Delta }m_N`$, is not large. For processes in which $`Qm_\pi `$ it is convenient to introduce the “chiral index” of an interaction $`\mathrm{\Delta }=d+\frac{f}{2}2`$, where $`d`$ is the number of small-scale factors, that is, derivatives, $`m_\pi `$, and $`\delta m`$; and $`f`$ is the number of fermion field operators. Chiral symmetry implies that $`\mathrm{\Delta }0`$ . Our interaction Lagrangian is given, using an appropriate choice of fields, by the expressions
$`_{\mathrm{int}}^{(0)}`$ $`=`$ $`{\displaystyle \frac{1}{4f_\pi ^2}}N^{}𝝉(𝝅\times \dot{𝝅})N+{\displaystyle \frac{g_A}{2f_\pi }}N^{}(𝝉\stackrel{}{\sigma }\stackrel{}{}𝝅)N`$ (2)
$`+{\displaystyle \frac{h_A}{2f_\pi }}[N^{}(𝑻\stackrel{}{S}\stackrel{}{}𝝅)\mathrm{\Delta }]+\mathrm{}`$
and
$`_{\mathrm{int}}^{(1)}={\displaystyle \frac{i}{8m_Nf_\pi ^2}}N^{}𝝉(𝝅\times \stackrel{}{}𝝅)\stackrel{}{}N{\displaystyle \frac{c_3}{f_\pi ^2}}N^{}(\stackrel{}{}𝝅)^2N`$ (3)
$`N^{}{\displaystyle \frac{\overline{c}_4}{2f_\pi ^2}}\stackrel{}{\sigma }\stackrel{}{}𝝅\times \stackrel{}{}𝝅𝝉N{\displaystyle \frac{ig_A}{4m_Nf_\pi }}N^{}𝝉\stackrel{}{\sigma }\dot{𝝅}\stackrel{}{}N`$ (4)
(5)
$`{\displaystyle \frac{h_A}{2m_Nf_\pi }}[iN^{}𝑻\dot{𝝅}\stackrel{}{S}\stackrel{}{}\mathrm{\Delta }]{\displaystyle \frac{d_1}{f_\pi }}N^{}𝝉\stackrel{}{\sigma }\stackrel{}{}𝝅NN^{}N`$ (6)
$`{\displaystyle \frac{d_2}{2f_\pi }}\stackrel{}{}𝝅\times N^{}\stackrel{}{\sigma }𝝉NN^{}\stackrel{}{\sigma }𝝉N+\mathrm{},`$ (7)
where $`\overline{c}_4=c_4+\frac{1}{4m_N}`$. The terms denoted by $`\mathrm{}`$ include Hermitian conjugates, $`s`$-wave $`\pi N`$ scattering terms and terms of higher powers in pion fields. Our principle aim is to determine the parameters $`d_i=𝒪(1/f_\pi ^2M_{QCD})`$, which determine the desired three-nucleon force. Here we fix signs by taking $`g_A=+1.26`$ in the chiral limit. The $`c_i`$ have been determined from $`\pi N`$ scattering at tree level ($`c_3^{(tree)}=3.90\text{GeV}^1`$ and $`c_4^{(tree)}=2.25\text{GeV}^1`$ ) as well as to one-loop order ($`c_3^{(loop)}=5.29\text{GeV}^1`$ and $`c_4^{(loop)}=3.63\text{GeV}^1`$ ). (Since we treat the delta isobar explicitly, we subtract its contribution from these values of $`c_i`$ .) As will be established below, up to next-to-leading order, the $`d_i`$, which support only $`SSp`$ transitions, are the only undetermined parameters in $`p`$-wave pion production.
The next step is to extend the power counting of Ref. to the region $`\eta 1`$, where the outgoing pion has energy $`\omega =𝒪(m_\pi )`$ and momentum $`|\stackrel{}{q}|=𝒪(m_\pi )`$, and the two nucleons in the final state have momentum $`|\stackrel{}{p^{}}|=𝒪(m_\pi )`$ and total energy $`p^0=𝒪(m_\pi ^2/m_N)`$. The unique difficulty of using $`\chi `$PT for pion production is that the entire pion energy is supplied by the relatively large momentum of the initial nucleons, $`|\stackrel{}{p}|=𝒪(\sqrt{m_Nm_\pi })`$. Note that the non-relativistic approximation holds, as $`p^4/8m_N^3m_\pi ^2/m_Nm_\pi p^2/2m_N`$.
The scales of momenta and energy are not the same, so it is simpler to count powers of the small scales in time-ordered perturbation theory. Equivalently, one first integrates over the time component of loop momenta in covariant diagrams. In this case, an intermediate state is associated with an energy denominator $`1/E`$, a loop with a $`Q^3/(4\pi )^2`$, a spatial (time) derivative with $`Q`$ ($`E`$), and a virtual pion vertex with $`1/E^{1/2}`$ from wave function normalization. For $`N`$, $`EQ^2/m_N`$, for $`\mathrm{\Delta }`$, $`EQ^2/m_N+\delta m`$, and for $`𝝅`$, $`E\sqrt{Q^2+m_\pi ^2}`$.
Final-state interactions (FSI) are those which occur after the emission of the real pion. In this case, the nucleons have typical $`Qm_\pi `$. The energies of intermediate states containing a $`𝝅`$ or $`\mathrm{\Delta }`$ can be $`Em_\pi `$, but otherwise $`Em_\pi ^2/m_N`$. The sum of “irreducible” sub-diagrams where all energies are $`𝒪(m_\pi )`$ is by definition the $`NN`$ potential, which is then amenable to a $`\chi `$PT expansion. The sum of “reducible” sub-diagrams produces the final-state wave function $`|\psi _f`$.
In contrast, all intermediate states occurring before the radiation of the real pion are characterized by loop momenta $`\sqrt{m_Nm_\pi }`$. For these kinematics we find that any additional loop requires at least i) one more interaction —pion exchange or shorter range— with an associated factor no larger than $`1/f_\pi ^2`$; ii) a volume integral with an associated factor of $`(\sqrt{m_\pi m_N})^3/(4\pi )^2`$; and iii) an additional time slice. If the additional time slice cuts a pion line, a factor of $`1/\sqrt{m_\pi m_N}`$ comes in, and the overall extra loop factor is at least $`\frac{1}{f_\pi ^2}\frac{(\sqrt{m_\pi m_N})^3}{(4\pi )^2}\frac{1}{\sqrt{m_\pi m_N}}=\frac{m_\pi }{m_N},`$ that is, a suppression by two powers of the expansion parameter. If the additional time slice does not cut a pion line, a factor of $`1/m_\pi `$ appears, and there is a relative enhancement of $`\sqrt{m_N/m_\pi }`$. Integrals over two-nucleon states typically also have enhancements by factors of $`\pi `$ from the unitarity cut. Thus we resum those diagrams that differ by the addition of interactions between the initial nucleons (ISI), and the effects are contained in an initial state wave function $`|\psi _i`$.
These considerations yield a pion production amplitude $`T=\psi _f|K|\psi _i.`$ Both the kernel $`K`$ and $`|\psi _{i,f}`$ can be obtained from the chiral expansion, but the currently available $`|\psi _i`$ do not yield an accurate fit to the measured $`NN`$ scattering phase shifts. Therefore we use a phenomenological coupled-channel ($`NN`$, $`N\mathrm{\Delta }`$, $`\mathrm{\Delta }\mathrm{\Delta }`$) model, CCF of Ref. , fitted to $`NN`$ scattering.
The leading contributions to $`p`$-wave production are displayed in Fig. 1. At lowest order ($`𝒪(1)`$, apart from overall factors) there are contributions from the direct production off the nucleon and off the delta, where all vertices are from $`^{(0)}`$ (Fig. 1i, ii). At next-to-leading non-vanishing order ($`𝒪(m_\pi /m_N)`$) there are four types of contributions. First, there is a recoil correction to the direct production. Second, there are rescattering diagrams that proceed through the seagull vertices in $`^{(1)}`$ proportional to $`1/4f_\pi ^2`$ (Galilean correction to the Weinberg-Tomozawa term), $`c_3`$, and $`c_4`$ (Fig. 1iii). Third, there is a rescattering through the Weinberg-Tomozawa term, where the primary production vertex is proportional to the external pion momentum. Fourth, there are short-range $`\pi (N^{}N)^2`$ interactions proportional to $`d_1`$ and $`d_2`$ (Fig. 1iv). Diagram iv) and most of the rescattering diagrams contribute to charged-pion production only.
With our theory in place, we consider the available pion production database. This has been enriched recently by very accurate determinations of spin observables at $`0.5`$ $`\stackrel{<}{}`$ $`\eta `$ $`\stackrel{<}{}`$ $`1`$ for $`pppp\pi ^0`$ , $`pppn\pi ^+`$ , and $`ppd\pi ^+`$ . It is useful to describe the total cross section in terms of components $`{}_{}{}^{2S+1}\sigma _{m}^{}`$, where $`S`$ is the initial $`NN`$ spin with projection $`m`$ along the direction of the incoming momentum. The $`{}_{}{}^{2S+1}\sigma _{m}^{}`$ can be expressed as linear combinations of the total cross section and the double polarization observables $`\mathrm{\Delta }\sigma _T`$ and $`\mathrm{\Delta }\sigma _L`$ .
In order to test convergence for the $`p`$-wave production, we need an observable where the lowest contributing partial wave is $`p`$ and the initial and final nucleons are not both in $`S`$ states. Such an observable exists, namely the $`{}_{}{}^{3}\sigma _{1}^{}`$ cross section in neutral-pion production with $`Pp`$ as the lowest partial waves contributing. While the ratios between double polarization observables and the total cross section, $`\mathrm{\Delta }\sigma _T/\sigma _{\mathrm{tot}}`$ and $`\mathrm{\Delta }\sigma _L/\sigma _{\mathrm{tot}}`$, have recently been accurately measured at IUCF , the total cross section is known to a much lesser accuracy (see the compilation of Ref. ). To determine the error of the total cross section, we simply take the total spread of the data as the error band. We defer a more detailed analysis until it can benefit from the soon-to-be-available much better data.
In $`p`$-wave production the lowest-order loop contribution enters one order higher, at $`𝒪((m_\pi /m_N)^{\frac{3}{2}})`$, than the rescattering terms of Fig. 1, and are ignored. Besides the coupling constants of the pion to the baryon fields, the only parameter that enters is $`c_3`$. We will use both values given above to get an estimate of loop effects on the final result.
The predictions of $`\chi `$PT are compared to the data in Fig. 2. Up to values of $`\eta 0.7`$ the data is well described. Deviations at higher energies might be due to higher partial waves entering, and/or to higher-order $`p`$-wave contributions. In any case, we see that sub-leading corrections are smaller than leading contributions throughout the range $`\eta `$ $`\stackrel{<}{}`$ $`1`$.
We next consider the amplitude for the $`{}_{}{}^{1}S_{0}^{}(^3S_1^3D_1)p`$ transition, denoted $`a_0`$, which has recently been extracted from the reaction $`pppn\pi ^+\text{[34]}`$. The loop corrections are again expected to be small, but the number of rescattering diagrams is larger, since isospin-odd operators (the recoil correction to the Weinberg-Tomozawa term as well as the $`c_4`$ term of Eq. (7)) enter. The striking feature of $`a_0`$ is that interactions proportional to the $`d_i`$’s also contribute. Because there seems to be reasonable convergence in the $`p`$ waves, we assume that they can be reliably computed and that we can attribute any deviation between theory and experiment to the effects of the terms involving the coefficients $`d_i`$. The contact interactions enter as the linear combination $`d_1+4d_2`$. Thus there is one unknown parameter to be fixed by the data. On the basis of dimensional analysis we expect $`d\frac{1}{5}(d_1+4d_2)=\frac{\delta }{f_\pi ^2M_{QCD}}`$ with $`\delta =𝒪(1)`$.
Our result for $`a_0`$ is shown in Fig. 3. We find a destructive interference between direct nucleon and delta contributions that makes $`a_0`$ small and more sensitive to sub-leading terms. For the $`c_i`$ parameters we employ the values extracted from the tree level fit to $`\pi N`$ scattering ($`c_i^{(tree)}`$). We use dipole form factors; to make contact with Ref. , we employ cutoff parameters $`\mathrm{\Lambda }=1`$ GeV for diagrams containing pion exchange and $`\mathrm{\Lambda }=m_\omega `$ for the contact interactions. The result for $`\delta =0`$ is not in disagreement with data, whereas a value of $`\delta =1`$ leads to a serious disagreement with experiment. In Ref. $`\delta =0.2`$ was shown to yield an important contribution to $`A_y`$ in $`Nd`$ scattering at energies of a few MeV. Using $`\delta =0.2`$ here is also consistent with the pion production data.
In contrast to $`{}_{}{}^{3}\sigma _{1}^{}`$, the result for $`a_0`$ is quite sensitive to the cutoff parameter used in the rescattering contribution, because the momentum range scanned by the $`c_4`$ term is quite large. For example, our results for $`a_0`$ can vary up to a factor of 2 if the corresponding cutoff parameter is increased to 2 GeV. The cutoff sensitivity is not a serious difficulty because it also occurs in calculations of three-nucleon forces. From the viewpoint of an effective field theory this can be simply understood: the large momentum pieces of the loop integrals involved in the evaluation of the $`c_4`$ contribution can be absorbed by a counterterm, namely $`d_2`$. Thus, the cutoff dependence of $`c_4`$ directly translates into a scale dependence of $`d_2`$. A reasonable phenomenological estimate should follow from using the same cutoff and parameter set in both calculations. On the experimental side, it is clear from Fig. 3 that a reduction of the uncertainty in the data would allow a stronger constraint on $`\delta `$. We find this a strong motivation to the continuation of the existing program on pion production.
We have shown that there is convergence in $`p`$-wave pion production, and that data on this reaction can be used to extract information about the three-nucleon force. It is clear that more accurate data would be very useful. In particular, the parameter $`d`$ could be extracted and the calculation of Ref. repeated to predict three-nucleon observables. We find it very gratifying that chiral symmetry provides a direct connection between pion production at energies $`350`$ MeV (IUCF) and $`Nd`$ scattering at energies $`10`$ MeV (Madison, TUNL).
We would like to thank Jim Friar, Brad Keister, and Daniel Phillips for useful discussions. This research was supported in part by the U.S. DOE, the NSF, and the Humboldt Foundation.
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# Pulling Pinned Polymers and Unzipping DNA
\[
## Abstract
We study a class of micromanipulation experiments, exemplified by the pulling apart of the two strands of double-stranded DNA (dsDNA). When the pulling force is increased to a critical value, an “unzipping” transition occurs. For random DNA sequences with short-ranged correlations, we obtain exact results for the number of monomers liberated and the specific heat, including the critical behavior at the transition. Related systems include a random heteropolymer pulled away from an adsorbing surface and a vortex line in a type II superconductor tilted away from a fragmented columnar defect.
\]
Recent years have seen an explosion in the use of single molecule techniques to probe biological and other “soft” materials. It now possible, for example, to monitor the breaking of individual “lock and key” bonds and the unfolding of individual proteins ; the mechanical properties of single DNA molecules and the behavior of single molecular motors have been characterized in great detail. In contrast to more traditional experiments, these new approaches give access to fluctuations on the scale of individual molecules, without the requirement for averaging over a macroscopic sample. One can, moreover, now push or pull directly on a micron-sized object, and watch how it responds. The potentially profound implications both for complex fluids and for biological physics—where single molecule techniques can often more closely mimic conditions in the cell than conventional assays—are only beginning to be explored.
Despite a number of notable contributions, theory has often been out-paced by these rapid experimental advances. Certainly, the tools available to analyze single-molecule experiments have not yet reached the level of sophistication and generality of theories of mesoscopic quantum systems. This is especially true when it comes to the role of quenched randomness, which, though often present, is typically neglected in initial theories of a given system. This Letter seeks to fill some of this gap. We study a class of micromanipulation experiments in which a polymer or other line-like object is pulled away from a confining potential well. An example of such a situation is the pulling apart of the two strands of dsDNA (fig. 1). Formally, the distance between the two strands may be viewed as the coordinate of a single polymer, and the base-pairing interactions between complementary strands as a potential well. At a critical value of the pulling force, a novel phase transition occurs in which the two strands are pulled completely apart. Aspects of this transition for a homopolymer (or, equivalently, DNA with all base pairs the same) have been studied in a related model of a flux line in a type II superconductor . Here, we show that the transition is markedly different for random heteropolymers. In particular, the number of monomers liberated as the transition is approached diverges much more strongly for heteropolymers than for homopolymers; similar differences appear in the specific heat. We calculate exact critical exponents and crossover functions for the random case.
Figure 1 sketches the DNA-opening experiment: One of the two single strands of a dsDNA molecule is attached to a glass slide, while a constant force $`𝑭`$ directed away from the slide acts on the end of the other strand. Methods for exerting a constant force on the piconewton scale have been developed by several groups . Under the influence of the force $`𝑭`$, the DNA partially “unzips” at one end, breaking $`m`$ base pairs. In thermal equilibrium, the degree of opening $`m`$ is of course a fluctuating quantity. Because the base sequence of protein-coding DNA appears to many statistical tests to be random and uncorrelated along the backbone (at least up to a length
scale set by the sequence’s mosaic structure), the free energy landscape in which $`m`$ fluctuates can be taken to have a quenched random component. Bockelmann, Essevaz-Roulet, and Heslot have performed an elegant series of experiments in a different statistical ensemble, measuring the average force required to hold the positions of both single strands fixed . However, because of subtleties associated with the thermodynamic limit in a single molecule system (see below), the two ensembles are not equivalent.
In the remainder of this paper, we first introduce a coarse-grained model for the interaction of the two single strands of the dsDNA. By focusing on the unzipping transition induced by pulling on the single strands, we can avoid treating most of the degrees of freedom explicitly, obtaining a problem that can be solved exactly by a mapping to a Markov process. Although for concreteness we will focus primarily on the DNA-opening realization of our model, our results also apply to a number of related physical systems, some of which will be described in the conclusion. Throughout, we set $`k_\text{B}=1`$.
The bulk melting transition of dsDNA (see inset to Fig. 1) can be described by a Peyrard-Bishop-like model . One views the two single strands as Gaussian polymers whose $`n^{\text{th}}`$ monomers have positions $`𝒓_1(n)`$ and $`𝒓_2(n)`$. Below the melting transition, it should be possible to neglect non-native base pairings. The interactions between the two strands, coarse-grained over a number of bases, can then be described by a phenomenological potential energy term $`V_n[𝒓(n)]=[1+\stackrel{~}{\eta }(n)]h[𝒓(n)]`$. Here $`h`$ is a short-ranged attractive potential whose strength is temperature-dependent, $`𝒓𝒓_1𝒓_2`$, and the variation with base sequence of the strength of the attraction between strands is described by $`\stackrel{~}{\eta }(n)`$, which we take to be a random variable with short-ranged correlations. The effective Hamiltonian then becomes, up to an uninteresting center of mass term,
$$_{\text{melt}}=_0^N𝑑n\left\{\frac{Td}{2b^2}\left(\frac{d𝒓}{dn}\right)^2+V_n[𝒓(n)]\right\}$$
(1)
where $`d`$ is the spatial dimension and $`b`$ is $`\sqrt{2}`$ times the Kuhn length of single-stranded DNA. Standard arguments show that the partition function $`Z_{\text{melt}}𝒟[𝒓(n)]\mathrm{exp}(_{\text{melt}}/T)`$ obeys an (imaginary) time-dependent Schrödinger equation.
The unzipping transition may be studied by adding to $`_{\text{melt}}`$ the term $`_{\text{pull}}=𝑭𝒓(0)=𝑭𝒓(N)_0^N𝑑n𝑭𝑑𝒓/𝑑n`$. If the term proportional to $`𝒓(N)`$, which should not affect the opening near $`n=0`$ for a sufficiently long polymer, is dropped, a time-dependent version of the “non-Hermitian quantum mechanics” studied in results. In the time-independent case (corresponding to pulling apart homopolymeric dsDNA) there is a sharp first order unzipping transition at a critical value of the pulling force $`F_\text{c}`$ satisfying $`ϵ_0(T)=F_\text{c}^2b^2/(2dT)`$, where $`ϵ_0<0`$ is the ground state energy of the Hermitian quantum mechanics problem obtained by setting $`F=0`$. In general, $`F^2b^2/(2dT)`$ is the free energy per monomer of the unzipped monomers aligned with the pulling force. The free energy per monomer of the dsDNA that has remained zipped is $`ϵ_0`$, independent of $`F`$. The physical interpretation of the unzipping transition is thus clear: For $`F<F_\text{c}(T)`$, the DNA minimizes its free energy by remaining in the double-stranded form, while for $`F>F_\text{c}`$ it is advantageous to pull apart as many bases as possible. As $`FF_\text{c}^{}`$, the free energy difference between the bound, double-stranded phase and the pulled out, single-stranded phase becomes very small, and thermal fluctuations unbind a large number of monomers near the end of the DNA. As the transition is approached, the equilibrium ensemble average of the number $`m`$ of monomers that are pulled out diverges like
$$m(F_\text{c}F)^1.\text{(homopolymer)}$$
(2)
Similarly, $`(mm)^2(F_\text{c}F)^2`$ near the transition. The thermal fluctuations about $`m`$ are thus comparable to $`m`$ itself. The divergence in (2) is analogous to the divergence in interface height near a wetting transition.
We now determine how results such as (2) are modified for a random DNA sequence. Sequence randomness is at worst a marginal perturbation at the ($`F=0`$) melting transition in 3 dimensions . The application of a Harris-like criterion , however, shows that the same cannot be true for the unzipping transition: The typical variation per monomer in the base-pairing energy of a pulled out section of length $`<m>`$ scales like $`<m>^{1/2}\sqrt{F_\text{c}F}`$, which vanishes more slowly as the transition is approached than the relevant energy difference $`|ϵ_0|F^2b^2/(2dT)F_\text{c}F`$. To determine the correct critical behavior, we focus on the free energy cost of pulling out a given monomer. Define $`(m)`$ to be the free energy of a dsDNA molecule, subject to an applied force $`F`$, of which exactly the first $`m`$ monomers are unzipped. The change in $``$ from pulling out one additional monomer should have the form
$$\frac{d}{dm}=f+\eta (m),$$
(3)
which may be integrated twice to obtain the partition function $`Z=_0^{\mathrm{}}𝑑m\mathrm{exp}[(m)/T]`$ . Here $`f`$ is the average free energy difference between an unzipped and a bound pair of complementary monomers. It vanishes like $`F_\text{c}F`$ near the transition, and reduces to the familiar $`|ϵ_0|F^2b^2/(2dT)`$ in the absence of sequence randomness. The additional term $`\eta (m)`$ takes account of sequence-dependent deviations from the average; it reflects the bare sequence (described by $`\stackrel{~}{\eta }(m)`$), dressed by thermal fluctuations. As long as $`\stackrel{~}{\eta }(m)`$ is a random variable with only short-ranged correlations, it is reasonable to expect that $`\eta (m)`$ should also be short-range-correlated, with a correlation length on the order of the typical size $`\xi `$ of the regions of local melting of the dsDNA strand . On long enough scales, we can then take $`\eta `$ to be Gaussian white noise, with correlator $`\overline{\eta (m)\eta (m^{})}=\mathrm{\Delta }\delta (mm^{})`$, where the overbar indicates a “disorder average” over the possible realizations of the quenched random base sequence. The parameters $`f`$ and $`\mathrm{\Delta }`$ may be calculated from the $`F=0`$ partition function $`Z_{\text{melt}}`$, for example in a low temperature expansion. The model summarized by Eq. (3) can also be derived from a discrete, Ising-like description of the dsDNA , and it still holds both when the single strands that have been liberated are characterized as freely-jointed or worm-like chains and when there are significant excluded volume interactions .
The study of the unzipping transition can thus be reduced to that of a single coordinate $`m`$ in the random potential $`(m)`$. One immediate consequence is that there is no large parameter that defines a thermodynamic limit, and thus no equivalence between the ensemble considered here and the conjugate ensemble in which $`m`$ is held fixed. Below the unzipping transition, $`f>0`$, and $`(m)`$ diverges with probability unity as $`m\mathrm{}`$. In the ensemble studied here, the unzipping fork is thus always confined to the vicinity of $`m=0`$. In the absence of randomness, the probability of unzipping $`m`$ monomers is $`(f/T)\mathrm{exp}(mf/T)`$, and one recovers, e.g. (2). If there is sequence randomness, the typical random contribution to $`(m)`$ is of order $`\sqrt{\mathrm{\Delta }m}`$; the random part thus exceeds the average contribution $`fm`$, which is responsible for the confinement, for $`m\mathrm{\Delta }/f^2`$. This length scale diverges faster than $`1/f`$ as $`f0`$, suggesting that a typical value of $`m`$ might show a $`1/f^2`$ divergence instead of the non-random $`1/f`$, with a crossover at $`f\mathrm{\Delta }/T`$.
Because of the confinement near $`m=0`$, the unzipping transition does not exhibit self-averaging (see below). Nonetheless, one can still calculate averaged quantities and distributions over the possible realizations of randomness. To do this, one wishes to study the disorder-averaged free energy $`T\overline{\mathrm{ln}Z}`$. The partition function $`\stackrel{~}{Z}(m)_0^m𝑑m^{}\mathrm{exp}[(m^{})/T]`$ of a finite-sized system of $`m`$ (bound or liberated) monomers satisfies
$$\frac{d\stackrel{~}{Z}}{dm}=e^{(m)/T}\text{and}\stackrel{~}{Z}(0)=0;$$
(4)
$`Z`$ follows simply by taking the limit of an infinitely long polymer: $`Zlim_m\mathrm{}\stackrel{~}{Z}(m)`$. Together, (3) and (4) form a system of coupled Langevin equations, analogous, for example, to those describing the Brownian motion of a massive particle, with $``$ playing the role of momentum and $`\stackrel{~}{Z}`$ that of position. The associated Fokker-Planck equation for the joint distribution $`P(,\stackrel{~}{Z};m)`$ of $``$ and $`\stackrel{~}{Z}`$ at “time” $`m`$ follows in the usual manner :
$$\frac{P}{m}=\left[\frac{\mathrm{\Delta }}{2}\frac{^2}{^2}f\frac{}{}e^{/T}\frac{}{\stackrel{~}{Z}}\right]P.$$
(5)
By Laplace transforming with respect to $`\stackrel{~}{Z}`$ and to $`m`$, one can solve (5) and obtain an exact expression for the partial distribution $`𝑑P(,\stackrel{~}{Z};m\mathrm{})`$ in terms of modified Bessel functions of order $`2fT/\mathrm{\Delta }`$; $`T\overline{\mathrm{ln}Z}`$ and other thermodynamic quantities then follow by integration.
The device of treating the quenched randomness as a Langevin noise thus leads to a number of exact results. In particular, the average degree of opening $`\overline{m}=T\overline{\mathrm{ln}Z}/f`$ satisfies
$`\overline{m}`$ $`=`$ $`{\displaystyle \frac{2T^2}{\mathrm{\Delta }\mathrm{\Gamma }(2fT/\mathrm{\Delta })}}{\displaystyle _0^{\mathrm{}}}𝑑yy^{2fT/\mathrm{\Delta }1}(\mathrm{ln}y)^2e^y`$ (7)
$`{\displaystyle \frac{2T^2\mathrm{\Gamma }^{}(2fT/\mathrm{\Delta })^2}{\mathrm{\Delta }\mathrm{\Gamma }(2fT/\mathrm{\Delta })^2}},`$
where $`\mathrm{\Gamma }^{}(z)=d\mathrm{\Gamma }(z)/dz`$. The small $`f`$ behavior of Eq. (7) is given by
$$\overline{m}\frac{\mathrm{\Delta }}{2f^2}(F_\text{c}F)^2,\text{(random heteropolymer)}$$
(8)
confirming our expectations for a crossover from a $`1/f`$ to a $`1/f^2`$ power law when $`f\mathrm{\Delta }/T`$. Similarly, the singular part of the heat capacity associated with the unzipping transition, $`C^2\overline{\mathrm{ln}Z}/T^2`$, crosses over from a $`1/f^2`$ to a $`1/f^3`$ divergence. One can also compute the disorder-averaged values of higher cumulants of $`m`$. For small $`f`$, $`\overline{m^2m^2}=T^2\overline{\mathrm{ln}Z}/f^21/f^3`$. Unlike in the non-random case, the square root of this quantity diverges more slowly than $`\overline{m}1/f^2`$, indicating that thermal fluctuations in $`m`$ for a given realization of the quenched randomness (and for thus a given heteropolymer) typically become small compared to the mean value as the transition is approached.
A real space renormalization group approach to the model of equations (3) and (4) due to Le Doussal, Monthus, and Fisher gives further insight into the unzipping transition. This technique gives leading order results in the limit $`f0`$, where the authors have argued that it should be exact. In this limit, it allows one to calculate the distribution $`Q(m)`$ of thermal average values over different realizations of randomness. This takes the form of a scaling function of $`mf^2/\mathrm{\Delta }`$:
$$Q(m)=\frac{f^2}{\pi \mathrm{\Delta }}e^{\frac{mf^2}{2\mathrm{\Delta }}}_0^{\mathrm{}}𝑑we^{\frac{wmf^2}{2\mathrm{\Delta }}}\frac{\sqrt{w}}{w+1}.$$
(9)
This distribution yields the same asymptotic behavior of $`\overline{m}`$ near the unzipping transition as the Fokker-Planck approach. It also predicts that $`\overline{(m\overline{m})^2}1/f^4`$; $`m`$ for a polymer with a given random sequence of base pairs can thus deviate significantly from the disorder average, and this system is not self-averaging.
For the randomness-dominated critical properties reported here to be observable, the variance $`\mathrm{\Delta }`$ in the base-pairing energy must be sufficiently large. Then at the crossover from non-random to random behavior, $`f`$ will also be large, and the typical value of $`mT^2/\mathrm{\Delta }`$ will be small enough that finite size effects do not become an issue. In this respect, dsDNA appears to be a very good candidate system. Under physiological conditions, the difference in free energy of binding between polymers with only A-T base pairs and those with only G-C base pairs is of order $`T`$, meaning that $`m`$ is only a few monomers when the crossover from pure to random behavior occurs.
In sum, we have described a randomness-dominated unzipping transition of dsDNA, obtaining exact expressions for the critical behavior and for the crossover from random to non-random scaling. Most notably, we find that when the base sequence is random and has only short-ranged correlations, the average degree of opening $`\overline{m}`$ diverges like $`1/(F_\text{c}F)^2`$ as the pulling force $`F`$ approaches a critical value $`F_\text{c}`$, in marked contrast to the $`1/(F_\text{c}F)`$ divergence found when all of the base pairs are identical. It should be possible to arrive at analogous results for the case of DNA whose base sequence has long-ranged correlations (as may be the case for non-coding DNA ). If a typical variation about the average energy grows like $`m^\beta `$, then balancing this against $`mf`$ suggests $`\overline{m}1/f^{1/(1\beta )}`$; the short-range-correlated case is recovered when $`\beta =1/2`$. The biological significance of our results remains to be determined: Processes such as DNA replication and recombination often involve unzipping of the dsDNA. Usually, however, this is accomplished by a molecular motor relying on an outside energy source, so non-equilibrium effects must be considered. More generally, the dynamics of the unzipping transition is an open question.
We have focussed on the case of unzipping DNA, but our results hold equally well for a number of more conventional condensed matter systems described by the Hamiltonian $`_{\text{melt}}+_{\text{pull}}`$ . The pulling of a Gaussian random heteropolymer away from an adsorbing surface is a natural extension of recent work on homopolymers . Other examples include a heteropolymer under tension pinned to a bulk defect , a magnetic flux line in a type II superconductor confined to a fragmented columnar pin and subject to a transverse field , and a simplified model of the corner wetting transition in two dimensions . Related models are likely to be relevant to the transverse surface magnetization and surface specific heat near the Bose glass transition of a bulk superconductor and to adhesion in a random environment .
After this work was submitted for publication, we learned that related results had been obtained, in a different physical context, for a discrete version of Eqns. (3) and (4. It is a pleasure to thank D. Branton, D.S. Fisher, and T. Hwa for helpful conversations and T. Hwa for bringing to our attention. This research was supported by the NSF through grant DMR97-14725 and through the Harvard MRSEC via grant DMR98-09363.
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# 1 Introduction
## 1 Introduction
Most of the applications of lattice gauge theories are based on and are employing their manifest gauge invariance. However, in order to get a better understanding of the structure of the lattice theory itself and to interprete correctly results obtained in Monte Carlo simulations, it is instructive to compare also gauge variant quantities such as gauge and fermion field propagators with corresponding analytic perturbative results. In this respect, compact lattice QED within the Coulomb phase serves as a very useful ‘test ground’. In the weak coupling limit this theory is supposed to describe non-interacting massless photons.
In order to fix the gauge the Lorentz (or Landau) gauge condition is normally applied. For non-Abelian gauge theories there is no unique solution, i.e. so-called Gribov copies occur . Within continuum QED such a problem arises, too, if the theory is defined on a torus . The lattice discretization may cause additional problems. Indeed, various lattice studies \[3 – 11\] have revealed nontrivial effects. The standard Lorentz (or Landau) gauge fixing procedure leads to a $`\tau `$–dependence of the non-zero-momentum transverse photon correlator inconsistent with the expected zero-mass behavior . Numerical and analytical studies have shown that there is a connection between ‘bad’ gauge (or Gribov) copies and the appearance of periodically closed double Dirac sheets (DDS). The removal of DDS by appropriate gauge transformations restores the correct perturbative behavior of the photon correlator at non-zero momentum, but it does not completely resolve the Gribov ambiguity problem. Gribov copies related to other local extrema of the gauge functional and connected with zero–momentum modes (ZMM) of the gauge fields still appear. They ‘damage’ gauge dependent observables such as the zero–momentum gauge field correlator and the fermion propagator , respectively.
There is a special Lorentz gauge, for which both the double Dirac sheets and the zero-momentum modes can be removed from the gauge fields . We call it zero–momentum Lorentz gauge (ZML) . It allows to reach the global extremum of the Lorentz gauge functional in almost $`100\%`$ of the cases. In comparison with the standard Lorentz gauge procedure (LG) it demonstrates very clearly the strong effects caused by the zero–momentum modes.
In the given talk we are going to review the results of with special emphasis on the question, how Gribov copies influence the Wilson–fermion propagator within the Coulomb phase of quenched QED. We want to show that a reliable estimate of the (renormalized) fermion mass requires either the removal of the zero-momentum modes or their proper perturbative treatment, when comparing the numerical results with analytic expressions. By employing the ZML-gauge we shall estimate the fermion mass in agreement with standard perturbation theory.
## 2 The Action and the Observables
We consider 4d compact QED in the quenched approximation on a finite lattice ($`V=N_s^3\times N_t`$). The pure gauge part of the standard Wilson action reads
$$S_G=\beta \underset{x,\mu <\nu }{}\left(1\mathrm{cos}\theta _{x,\mu \nu }\right),$$
(1)
with the plaquette angle $`\theta _{x,\mu \nu }=\theta _{x,\mu }+\theta _{x+\widehat{\mu },\nu }\theta _{x+\widehat{\nu },\mu }\theta _{x,\nu }`$ related to the link variables $`\theta _{x,\mu }(\pi ,\pi ].`$ $`\beta =1/e_0^2`$ is the inverse bare coupling. The lattice spacing is put $`a=1`$, for simplicity.
The fermion part is given by
$$S_F=\underset{x,y}{}\overline{\psi }_x𝐌_{xy}(\theta )\psi _y,𝐌=\mathrm{𝟏}\kappa 𝐃,$$
(2)
$$𝐃_{xy}=\underset{\mu =1}{\overset{4}{}}\left\{U_{x,\mu }P_\mu ^{}\delta _{y,x+\widehat{\mu }}+U_{x\widehat{\mu },\mu }^{}P_\mu ^+\delta _{y,x\widehat{\mu }}\right\},$$
where $`P_\mu ^\pm =\widehat{1}\pm \gamma _\mu `$ and $`U_{x,\mu }=\mathrm{e}^{\mathrm{i}\theta _{x,\mu }}.`$ The hopping-parameter $`\kappa `$ is related to the bare mass $`m_0`$ by $`\kappa =1/(8+2m_0).`$
In quenched QED the observables have to be averaged with respect to the gauge field $`\{\theta \}`$ with the weight $`\mathrm{exp}(S_G)`$. We imply periodic boundary conditions (b.c.) except for the fermion fields, which we choose to be anti-periodic in the $`x_4\tau `$-direction.
The first gauge variant observable we are going to discuss is the transverse photon correlator at non-zero momentum
$`\mathrm{\Gamma }_T^{\mathrm{ph}}(\stackrel{}{p};\tau )=\mathrm{\Phi }(\stackrel{}{p};\tau )\mathrm{\Phi }^{}(\stackrel{}{p};0),`$ (3)
$`\mathrm{\Phi }(\stackrel{}{p};\tau )={\displaystyle \underset{\stackrel{}{x}}{}}\mathrm{exp}(i\stackrel{}{p}\stackrel{}{x}+{\displaystyle \frac{i}{2}}p_\mu )\mathrm{sin}\theta _{\stackrel{}{x}\tau ,\mu }`$
with $`(\mu =1,3,\stackrel{}{p}=(0,p,0)).`$ The second one is the fermion propagator. For a given gauge field $`\{\theta \}`$ we have
$$\mathrm{\Gamma }(\tau )=\frac{1}{V}\underset{\stackrel{}{x},x_4}{}\underset{\stackrel{}{y}}{}𝐌_{\stackrel{}{x},x_4;\stackrel{}{y},x_4+\tau }^1(\theta ).$$
(4)
In the following we shall restrict ourselves to the vectorial part
$$\mathrm{\Gamma }_V(\tau )=\frac{1}{4}\mathrm{Re}\mathrm{Tr}\left(\gamma _4\mathrm{\Gamma }(\tau )\right),$$
(5)
with the trace taken with respect to the spinor indices. For the b.c. mentioned above, $`\mathrm{\Gamma }_V(\tau )`$ is an even function of $`\tau N_t/2.`$
Lateron, we shall compare the expectation value $`\mathrm{\Gamma }_V`$ with the result of a simple approximation, which takes only constant gauge field modes into account. The correlator in a uniform background $`\theta _{x,\mu }\varphi _\mu ,\pi <\varphi _\mu \pi ,`$ $`\mu =1,\mathrm{},4`$ can be represented as
$`\mathrm{\Gamma }_V(\tau ;\varphi )`$ $`=`$ $`{\displaystyle \frac{1\delta _{\tau ,0}}{2(1+)}}\times `$
$`\times `$ $`{\displaystyle \frac{[^\tau +^{2N_t\tau }]\mathrm{cos}(\varphi _4\tau )+[^{N_t+\tau }+^{N_t\tau }]\mathrm{cos}[\varphi _4(N_t\tau )]}{1+^{2N_t}+2^{N_t}\mathrm{cos}(\varphi _4N_t)}},`$
where
$$=1+\frac{^2+𝒦^2}{2(1+)}\frac{\sqrt{^2+𝒦^2}\sqrt{(+2)^2+𝒦^2}}{2(1+)};$$
$$=m_0+\underset{l=1}{\overset{3}{}}\left(1\mathrm{cos}\varphi _l\right),𝒦=\sqrt{\underset{l=1}{\overset{3}{}}\mathrm{sin}^2\varphi _l},m_0>0.$$
For $`\varphi _\mu =0,\mu =1,\mathrm{},4`$ the free fermion correlator for finite lattice size is reproduced.
## 3 Lorentz Gauge Fixing
In numerical simulations the Lorentz gauge is fixed by iteratively maximizing the gauge functional
$$F(\theta )=\frac{1}{V_4}\underset{x}{}F_x(\theta );F_x(\theta )=\frac{1}{8}\underset{\mu =1}{\overset{4}{}}\left[\mathrm{cos}\theta _{x\mu }+\mathrm{cos}\theta _{x\widehat{\mu };\mu }\right]$$
(7)
with respect to the (local) gauge transformations
$$U_{x\mu }\mathrm{\Lambda }_xU_{x\mu }\mathrm{\Lambda }_{x+\widehat{\mu }}^{};\mathrm{\Lambda }_x=\mathrm{exp}\{i\mathrm{\Omega }_x\}U(1).$$
(8)
The algorithm is called standard Lorentz gauge fixing (LG), if it consists only of local maximization and overrelaxation steps with respect to gauge transformations periodic in space-time. The standard procedure gets normally stuck into local maxima of the gauge functional (7) (gauge copies).
It has been argued that the Gribov problem has to be solved by searching for the global maximum providing the best gauge copy . In we have shown that in order to reach the global maximum we have necessarily to suppress both the double Dirac sheets (DDS) and the zero-momentum modes (ZMM) in the gauge fields. Let us explain this in more detail.
DDS can be identified as follows. The plaquette angle $`\theta _{x,\mu \nu }`$ is decomposed into the gauge invariant (electro-) magnetic flux $`\overline{\theta }_{x,\mu \nu }(\pi ,\pi ]`$ and the discrete gauge-dependent contribution $`2\pi n_{x,\mu \nu },n_{x,\mu \nu }=0,\pm 1,\pm 2`$ . The latter represents a Dirac string passing through the given plaquette if $`n_{x,\mu \nu }=\pm 1`$ (Dirac plaquette). A set of Dirac plaquettes providing a world sheet of a Dirac string on the dual lattice is called Dirac sheet. DDS consist of two sheets with opposite flux orientation extending over the whole lattice and closing themselves by the periodic b.c. They can easily be identified by counting the total number of Dirac plaquettes $`N_{DP}^{(\mu \nu )}`$ for each choice $`(\mu ;\nu )`$. The necessary condition for the occurence of DDS is that at least for one of the six possibilities $`(\mu ;\nu )`$ holds
$$N_{DP}^{(\mu \nu )}2\frac{V}{N_\mu N_\nu }.$$
(9)
DDS can be removed by periodic gauge transformations.
The ZMM of the gauge field
$$\varphi _\mu =\frac{1}{V}\underset{x}{}\theta _{x,\mu }$$
(10)
do not contribute to the pure gauge field action either. For gauge configurations representing small fluctuations around constant modes it is easy to see, that the global maximum of the functional (7) requires $`\varphi _\mu 0`$. The latter condition can be achieved by non-periodic gauge transformations
$$\theta _{x,\mu }\theta _{x,\mu }^c=c_\mu +\theta _{x,\mu }\mathrm{mod}2\pi ,c_\mu (\pi ,\pi ].$$
(11)
We realize a proper gauge fixing procedure as proposed in . Successive Lorentz gauge iteration steps are always followed by non-periodic gauge transformations suppressing the ZMM. Additionally we check, whether the gauge fields contain yet DDS. The latter can be excluded by repeating the procedure with initial random gauges. We call the combined procedure zero-momentum Lorentz gauge (ZML gauge). It yields the global maximum of the gauge functional with very high accuracy.
In Figures 1 and 2 we show, how the achieved values of the gauge functional (7) are correlated with the occurence of DDS visible as sharp peaks in the number of Dirac plaquettes. Whereas for LG strong fluctuations occur, they disappear after ZML gauge. The few DDS seen in Fig. 2 are easily removed by restarting the procedure with random initial gauges. Random gauges can also be used in order to convince oneself that the ZML gauge prescription leads to the global maximum of the gauge functional in more than $`99\%`$ of the cases.
## 4 Results
First let us convince ourselves that the removal of the above mentioned gauge copies leads to the correct behaviour of the transverse photon propagator.
In Fig. 3 we show the normalized correlator $`\mathrm{\Gamma }_T^{ph}(\stackrel{}{p};\tau )/\mathrm{\Gamma }_T^{ph}(\stackrel{}{p};0)`$ for lowest non-vanishing momentum and for different Lorentz gauge prescriptions. For the standard one (LG) we see a clear deviation from the expected perturbative zero-mass result. We show also the result obtained with an axial Lorentz gauge (ALG) using an initial maximal-tree axial gauge condition , which provides a ’unique’ prescription. The latter gauge fixing prescription turns out to be even worse! On the other hand the ZML gauge provides an excellent agreement with the perturbative result. In fact, as we convinced ourselves earlier, it is already sufficient to remove the DDS gauge copies in order to reach this agreement . The given observations do not change, when $`\beta `$ and/or the lattice size are increased considerably .
In the following we want to concentrate on the pure effect of the ZMM. Therefore, we compare the ZML gauge with a version of the standard Lorentz gauge, where the DDS are removed and the ZMM are left. We shall abbreviate the latter version also by LG.
For both the LG and ZML gauges we have computed the averaged fermion correlator employing the conjugate gradient method and point-like sources. In the upper part of Fig. 4 we have plotted $`\mathrm{\Gamma }_V(\tau )`$ (normalized to unity at $`\tau =1`$). The situation seen is typical for a wide range of parameter values within the Coulomb phase. Obviously, there is a strong dependence of the fermion propagator on the gauge fixing procedure resulting in the presence or absence of ZMM. The masses to be extracted seem to have different values. Let us determine the effective mass $`m_{eff}(\tau )`$ in accordance with
$$\frac{\mathrm{\Gamma }(\tau +1;\theta )_\theta }{\mathrm{\Gamma }(\tau ;\theta )_\theta }=\frac{\mathrm{cosh}[E(\tau )(N_t/2\tau 1)]}{\mathrm{cosh}[E(\tau )(N_t/2\tau )]}$$
(12)
where $`E(\tau )=\mathrm{ln}(m_{eff}(\tau )+1)`$. See the lower part of Fig. 4. In the LG case no plateau is visible, whereas the ZML case provides a very stable one. Thus, only the ZML gauge yields a reliable mass estimate, whereas the standard method to fix the Lorentz gauge obviously fails.
To get deeper insight into the effect of ZMM for the LG case (with DDS suppressed) we measure the probability distributions $`P(\varphi )`$ for the space- and time-like components of ZMM according to Eq. (10). The distributions turn out to be flat up to an effective cutoff at $`|\varphi _\mu |\pi /N_\mu `$ and to be widely independent of $`\beta `$. In accordance with Eq. (2) we compute the fermion propagator for constant modes in the LG case and average
$$\mathrm{\Gamma }_V(\tau ;\varphi )_\varphi =[\mathrm{d}\varphi ]P(\varphi )\mathrm{\Gamma }_V(\tau ;\varphi ).$$
(13)
The results for several parameter sets are presented in Fig. 5 together with the corresponding free (i.e. zero-background) propagator. We see clearly that the constant mode contributions strongly change the behavior of the fermion propagator and, naively speaking, produce a larger mass.
Finally, in Fig. 6 we present the fermion mass extracted from the vector fermion propagator within the ZML gauge for $`\beta =2.0`$ and various $`\kappa `$-values. We see a nice linear behaviour from which by extrapolating to zero mass (solid line) we estimate the critical value $`\kappa _c=0.1307\pm 0.0001`$.
## 5 Conclusions
We have studied the effect of different gauge copies of the gauge field on gauge dependent correlators, in particular on the Wilson fermion propagator. We have convinced ourselves that the standard Lorentz gauge fixing prescription to maximize the functional (7) provides gauge copies with DDS and ZMM. These modes disturb the photon and the fermion correlator in comparison with perturbation theory and consequently spoil the (effective) mass estimate. A Lorentz gauge employing non-periodic gauge transformations in order to suppress the ZMM – additionally to DDS – (the ZML gauge) allows to reach the global maximum of the Lorentz gauge functional. Furthermore, it provides a reliable fermion mass determination, at least, if $`\kappa `$ is chosen not too close to the chiral critical line $`\kappa _c(\beta )`$. A computation of the fermion propagator with constant background gauge fields taken from the ZMM of the quantum fields demonstrates the disturbing effect of these modes very clearly. Moreover, it shows the effect to be independent of the bare coupling and not to disappear for large volumes.
So far, we have studied the quenched approximation of U(1) lattice gauge theory. The gauge action (1) is invariant under non-periodic gauge transformations (11). Thus, we are allowed to use the ZML gauge for evaluating gauge dependent objects. Contrary to the gauge action, the fermionic part (2) does depend on the ZMM because of the (anti-) periodic boundary conditions. In this case another way of dealing with the Gribov problem has to be searched for.
The problems we have discussed here for compact QED show that gauge fixing has to be carried out and to be interpreted with care. This lesson has to be taken into account also in lattice QCD when extracting masses from gauge variant gauge and fermion correlators, respectively.
## Acknowledgements
The work has been supported by the grant INTAS-96-370, the RFRB grant 99-01-01230 and the JINR Dubna Heisenberg-Landau program.
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# Critical behavior of a non-equilibrium interacting particle system driven by an oscillatory field.
## Abstract
First- and second-order temperature driven transitions are studied, in a lattice gas driven by an oscillatory field. The short time dynamics study provides upper and lower bounds for the first-order transition points obtained using standard simulations. The difference between upper and lower bounds is a measure for the strength of the first-order transition and becomes negligible small for densities close to one half. In addition, we give strong evidence on the existence of multicritical points and a critical temperature gap, the latter induced by the anisotropy introduced by the driving field.
Pacs 64.60.Cn {Order-disorder transitions. Statistical Mechanics of model systems.}
Pacs 82.20.Mj {Nonequilibrium kinetics.}
Pacs 66.30.Hs {Self-diffusion and ionic conduction in nometals.}
Far from equilibrium systems (FFES) are ubiquitous in nature and their theoretical understanding will contribute to the progress of scientific areas in physics, chemistry, biology, ecology, economy, etc. Since the theoretical development of non-equilibrium statistical mechanics is still in its infancy, a useful approach to FFES is to study simple models by means of various techniques such as numerical simulations, mean-field approximations, phenomenological scaling, field-theoretical developments, etc. Within the broad context of FFES, driven diffusive systems (DDS) have very recently received growing attention ; for reviews see e.g. . The classical model for DDS was proposed by Katz et al. (KLS) and is based on the equilibrium Ising model . Using the lattice gas language, the KLS model introduces an external driving field to the Ising model. However, due to this modification, the system now evolves towards a non-equilibrium stationary-state (NESS). In spite of considerable effort devoted to the study of the KLS model, there are still controversies on the understanding of numerical data and its theoretical description is the subject of an ongoing debate .
In this work we study a DDS subjected to the action of an oscillating driving field. One of the motivations for this approach is that the periodical field can be realized in numerous practical applications such as charged colloids between the plates of a capacitor , electrophoresis experiments in pulsed fields , gas condensation in the presence of ultrasonic waves , segregation of granular materials in vibrating containers, etc.
The aim of this work is to perform an extensive simulation study of the dependence of the temperature-driven transitions of the model on both the density of particles and the magnitude of the field. Measurements of stationary properties combined to an study of the short time dynamics allow us to drawn a detailed phase diagram of the model that lead us to the discovery of a multicritical point. Furthermore, we developed a coupled mean field approach that yields results in agreement with the simulations.
The model is defined on the square lattice assuming a rectangular geometry $`L_x,L_y`$, using “brick wall” (periodic) boundary conditions across (along) the $`y`$ ($`x`$)axis where the oscillatory field is applied, respectively. A lattice configuration $`\eta `$ is specified by the set of occupation numbers $`n_{i,j}=\{0,1\}`$, corresponding to each site of coordinates $`(i,j)`$, i.e. $`\eta =\{n_{i,j}\}`$. Nearest-neighbor attraction with a coupling constant $`J>0`$, is considered. So, in the absence of a field the Hamiltonian $``$ is given by
$$=4J\underset{<ij,i^{}j^{}>}{}n_{i,j}n_{i^{},j^{}},$$
(1)
where the summation is over nearest-neighbor sites only. The driving oscillatory field $`E`$ acts along the $`\pm y`$direction with half-period $`\tau `$. The coupling to a thermal bath at temperature $`T`$ and the action of the field are considered through Metropolis jump rates, namely $`min[1,exp(\{\mathrm{\Delta }ϵE(\tau )\}/k_BT)`$, where $`k_B`$ is the Boltzmann constant, $`\mathrm{\Delta }`$ is the change in $``$ after the exchange, and $`ϵ=(1,0,1)`$ for a particle attempting to hop (against, orthogonal, along) the driving field, respectively. For $`E=0`$ and half-filled lattices, the model reduces to the Ising model in absence of magnetic field. In the thermodynamic limit the Ising model exhibits a second-order phase transition at a temperature $`T_c^I=2.2692..J/k_B`$.
Monte Carlo simulations are performed on lattices of aspect ratios $`L_x/L_y=2`$ and $`1`$, with $`30L_y480`$. $`T`$ is reported in units of $`J/k_B`$ and $`E`$ is given in units of $`J`$. The starting configuration is obtained by randomly filling the sample with probability $`\rho _o`$, which is also the density of particles that remains constant. One Monte Carlo time step (mcs) involves $`L_xLy`$ trials. Data are obtained disregarding $`10^6`$ mcs in order to allow the system to reach a NESS, and averages are taken over the subsequent $`10^6`$ mcs. Using this procedure a single data point, as e.g. shown in figure 2, requires $`1`$ day of CPU time in an AMD 700 MHz processor.
The model has also been studied by means of a coupled mean-field (CMF) approach. In order to write down the CMF equations the local density of particles $`\rho _{i,j}`$ at site $`(i,j)`$ is defined which is the probability of finding a particle in this site. Due to normalization, one has $`\rho _{i,j}+h_{i,j}=1`$, where $`h_{i,j}`$ is the probability for the site $`(i,j)`$ to be empty. Then, one has to consider all events that may cause $`\rho _{i,j}`$ to change. $`\rho _{i,j}`$ may increase by the arrival of particles due to unbiased (biased) diffusion perpendicular (parallel) to the driving field, respectively. Similarly, the density may decrease by an outgoing flux of particles to neighboring sites. The implementation of the CMF leads to a set of $`L_xLy`$ coupled non-linear differential equations. Here, we will only sketch out the form of such equations for the sake of space. Let $`\eta ^{}[(i,j);(i^{},j^{})]`$ be the configuration obtained from $`\eta `$ by interchanging the content of site $`(i,j)`$ with that of a neighboring site $`(i^{},j^{})`$. Then, the Metropolis rates are functions $`F`$ of $`(\eta ^{})(\eta )ϵE(\tau )=\mathrm{\Delta }[(i,j);(i^{},j^{})]ϵE(\tau )`$. So, $`\rho _{ij}`$ evolves in time according to:
$`{\displaystyle \frac{d\rho _{i,j}}{dt}}=h_{i,j}\{\rho _{i+1,j}F\{\mathrm{\Delta }[(i,j);(i+1,j)],T\}+\rho _{i1,j}F\{\mathrm{\Delta }[(i,j);(i1,j)],T\}+`$ (2)
$`\rho _{i,j+1}F\{\mathrm{\Delta }[(i,j);(i,j+1)],T,E(\tau )\}+\rho _{i,j1}F\{\mathrm{\Delta }[(i,j);(i,j1)],T,E(\tau )\}\}`$
$`\rho _{i,j}\{h_{i+1,j}F\{\mathrm{\Delta }[(i,j);(i+1,j)],T\}+h_{i1,j}F\{\mathrm{\Delta }[(i,j);(i1,j)],T\}+`$
$`h_{i,j+1}F\{\mathrm{\Delta }[(i,j);(i,j+1)],T,E(\tau )\}+h_{i,j1}F\{\mathrm{\Delta }[(i,j);(i,j1)],T,E(\tau )\}\}.`$
The set of equations (3) is solved numerically starting from a random initial distribution of particles and using an integration time step of $`\mathrm{\Delta }t=0.25`$, in arbitrary units. Numerical integrations are performed until $`t=25000`$ and averages are taken for $`t20000`$. In the CMF approach the excluded volume interaction is taken into account in a probabilistic way and stochastic fluctuations are disregarded, in contrast to the Monte Carlo method which has intrinsic fluctuations and excluded volume is deterministically satisfied. However, the CMF approach is derived directly form the microscopics, so it contains the same symmetries than the lattice gas model. One advantage of the CMF method is that one can obtain the spatial mass distribution. In fact, figure 1 corresponds to a NESS where a multi stripped pattern is observed.
An intriguing feature of driven dissipative systems is the occurrence of highly ordered and complex patterns as shown in figure 1. Since the system constantly gains (loss) energy from (to) the external field (thermal bath), respectively, the observed stationary states are by no means equilibrium states. In fact, they are truly non-equilibrium steady states.
In order to perform a quantitative investigation, the longitudinal order parameter ($`OP_x`$) is defined as the excess density, namely
$$OP_x(RL_x)^1\underset{i=1}{\overset{L_x}{}}|P(i)\rho _o|,$$
(3)
where $`P(i)=(L_y)^1_{j=1}^{L_y}n_{ij}`$ is the density profile along the $`x`$direction and $`R=(2\rho _o(1\rho _o))`$ is a normalization constant. Similarly, $`OP_y`$ can also be defined.
The dependence of the nature of the ordered phase on the period of the applied field has been investigated . For temperatures below criticality, it is found that for short periods (say $`\tau <4L_y`$) the system exhibits NESS with stripped patterns such as that shown in figure 1. However, for larger periods (say $`\tau >4L_y`$) the system alternates between almost equilibrium states (AES) such as those corresponding to molecules in a gravitational field. The crossover from NESS to AES has a characteristic time of the order of $`\tau 4L_y`$. In this work, we are interested in the critical behavior of NESS so we have restricted ourselves to the case $`\tau =10`$ mcs, without loosing generality because the same behavior will be valid for periods such as $`\tau <4L_y`$ for finite lattices and all periods in the thermodynamic limit. So, $`\tau `$ plays an important role in this model. In fact, for the case treated in this work, namely $`\tau <4L_y`$, $`OP_x`$ is a well defined quantity independent of time $`t`$. However, for $`\tau >4L_y`$, $`OP_x`$ and $`OP_y`$ are functions of time $`t`$, since the system alternates between AES as mentioned above. So, the half-period changes the nature of the problem and a crossover from NESS to AES is observed . In addition, since the oscillatory field causes the current of the driven gas averaged over long times to vanish, the symmetries of the model are different from those of the KLS model. From the theoretical point of view, this fact is essential to establish the universality class of the model, as will be discussed below.
Figure 2(a) shows results obtained for $`E=1`$. For low densities ($`\rho _o0.15`$) the observed transitions are abrupt and exhibit strong metastability, so they are first-order. In contrast, for $`\rho _o0.40`$ one observes second-order or very weak first-order like behavior. Notice that for $`\rho _o=0.20`$ and $`\rho _o=0.40`$ we have also included data which demonstrate the particle-hole exchange-invariance of the results. The existence of both first- and second-order transitions can also be observed by using the CMF approach. These results are in excellent agreement with Monte Carlo data, as shown in the inset of figure 2(a). Figure 2(b) and 2(c) show that for low densities, $`T_c`$ depends on the amplitude of the field, so that the higher the field the lower the $`T_c`$. Furthermore, these figures also reveal that weaker first-order transitions are obtained for smaller amplitudes of the field. Remarkably, results obtained by means of the CMF approach exhibit the same decreasing trend than the Monte Carlo data.
Figure 3 shows the phase diagrams obtained for a fixed lattice size ($`L_x/L_y=2,L_y=120`$) and two values of the driving field, namely $`E=1`$ and $`E=50\mathrm{}`$.
Using a method recently proposed for the study of the short time dynamics of weak first-order transitions it is possible to determine both lower and upper bounds for $`T_c(\rho _o)`$ valid in the thermodynamic limit and further generalize the phase diagram for $`E>1`$. The idea behind the proposed method is based on the existence of two pseudo critical points at $`T^{}`$ and $`T^{}`$ near the weak first-order transition point $`T_c`$ with $`T^{}<T_c<T^{}`$. These points can be obtained accurately from two short time dynamical processes starting from fully disordered and zero temperature states, respectively. In second-order transitions $`T^{}`$ and $`T^{}`$ overlap with the transition point $`T_c`$, so the difference between $`T^{}`$ and $`T^{}`$ also gives a criterion for the weakness of the first-order transition . Consider a system at $`T<T_c(\rho _o)`$ and the evolution process from a fully disordered state. Due to the geometrical constrained $`L_x/L_y^\varphi 1`$ ($`\varphi 0.2`$) configurations at short times exhibit multi-stripped patterns that are long lived, only relaxing to the single stripe state after a time of the order $`tL_x^3L_y`$ . Even in the case of square geometry, both the present model and the KLS model display multi-stripped configurations up to $`t10^5`$ mcs . It is then clear that the short time dynamics must be studied using an order parameter which takes into account multi-stripped configurations as that given by eq. (3).
Our results for the short time dynamical behavior have been summarized in figure 4.
For the used density ($`\rho _o=0.16`$) power laws have been obtained for $`T^{}=2.76`$ (figure 4(a)) and $`T^{}=2.40`$ (figure 4(b)) starting from ordered and fully disordered states, respectively. Also, figure 4(c) shows that the lower bound given by the short time dynamics is independent of the lattice size. Notice that the curves obtained for different aspect ratios are shifted but the power law behavior is obtained at the same temperature. The same results have been obtained for the upper bound, pointing out that the bounds drawn in the phase diagram (figure 3) are independent of the lattice size and consequently also valid in the thermodynamic limit. The transition points estimated using a finite lattice (figure 2) satisfy $`T^{}<T_c<T^{}`$ as would also do the true transition points in the $`L_x,L_y\mathrm{}`$ limit. Also, the difference $`\mathrm{\Delta }T=T^{}T^{}`$ depends on the strength of the first-order transition while $`\mathrm{\Delta }T0`$ at the second-order transition point for $`\rho _o=1/2`$.
Coming back to the phase diagram, it is found that for $`\rho _o0.30`$, $`T_c(E)`$ steadily increases with the strength of the field, reaching a saturation value at $`T_c(E=\mathrm{})1.41T_c(E=0)`$ for $`\rho _o=1/2`$, in excellent agreement with results for the KLS model . However, for lower densities (e.g. for $`\rho _o<0.1`$, in figure 3) $`T_c(E)`$ steadily decreases when increasing the magnitude of the field. So, $`T_c(E)`$ exhibits opposite trends depending on the density and consequently, it is expected that for some characteristic density $`\rho _o^M`$ ($`0.20\rho _o^M0.15`$) the critical temperature will be the same for all magnitudes of the driving field. Therefore, the point $`(\rho _o^M,T_c(E,\rho _o^M))`$ is a multicritical point, in the sense that for these special values of density and temperature this point is a critical point for all values of the amplitude of the field. Due to the observed symmetry, $`(1\rho _o^M,T_c(E,\rho _o^M))`$ is also a multicritical point.
Assuming that the critical curves have the simplest form allowed by the symmetry of the system, we propose the following expression for the critical temperature:
$`T_c(E,\rho _o)`$ $`=T_c(\mathrm{},1/2)k_{\mathrm{}}f(E)({\displaystyle \frac{1}{2}}\pm \rho _o^M)^{1/\beta }k_{\mathrm{}}(1f(E))({\displaystyle \frac{1}{2}}\pm \rho _o)^{1/\beta },E>0,`$ (4)
where for $`E\mathrm{}`$, $`k_{\mathrm{}}`$ is the coefficient of the higher order term and $`f(E)0`$, respectively. Equation (4) can be thought as the first approximation to the phase coexistence curve, valid close to $`\rho _o=1/2`$, so that $`\beta `$ is the order parameter critical exponent of the second-order transition. In order to fit eq. (4) to the data we will first summarize the symmetries present in our model. The model exhibits full translational and reflexion invariance as the Ising model, but the rotational symmetry is broken due to the anisotropy introduced by the field. If we consider short time scales, the up-down symmetry is also broken by the field. However, a renormalization group study will consider the system at a coarse-grained level. Then, we expect that the up-down symmetry will be restored at long time scales. Consequently, the present model displays the same symmetries than the randomly driven lattice gas with $`\beta =\frac{1}{3}`$ . Taking this value for $`\beta `$, the critical curve for $`E=\mathrm{}`$ can be fitted using a single parameter, yielding $`k_{\mathrm{}}=15\pm 3`$. Assuming that $`f(E)=\mathrm{exp}(E)`$, $`\rho _o^M`$ is the only parameter left to be fitted, yielding $`\rho _o^M=0.160\pm 0.005`$ for $`E=1`$ (see figure 4). Discrepancies between the fit and the data for densities far from $`\rho _o=1/2`$ are expected since the expansion given by eq. (4) holds close to that point only. Notice that equation (4) satisfies that $`(\rho _o^M=0.160\pm 0.005,T_c(\rho _o^M)=2.59\pm 0.01))`$ is a multicritical point. This value is in agreement with the estimation performed using the short time dynamics study that gives $`T^{}=2.76>T_c(\rho _o^M)>T^{}=2.40`$ (see figure 4). The existence of the multicritical point can also be confirmed by means of a short time dynamics simulations. In fact, figure 4(d) shows that plots of $`OP_x`$ versus $`t`$ obtained for different fields ($`1E\mathrm{}`$) yield the same lower bound for $`T_c(\rho _o^M)`$ given by $`T^{}=2.40`$, independently of the strength of the field. This behavior is characteristic of the multicritical point, as observed in figure 3. It should be noticed that fits of the phase diagram assuming $`\beta =1/2`$, as theoretically expected for the KLS model , are far from being satisfactory. Also, an excellent fit of the curve can be obtained assuming $`\beta =1/4`$ (yielding $`(\rho _o^M=0.168\pm 0.005,T_c(\rho _o^M)=2.57\pm 0.01)`$), but this value of the exponent is not supported by the symmetry considerations above mentioned.
For the sake of comparison we have included in the phase diagram the critical temperature of the Ising model $`T_c^I`$ as well as the prediction of eq. (4) in the $`E0`$ limit. The latter is in excellent agreement with the lower bound estimate given by the short time dynamics method for $`E=0.01`$. Notice that these estimations for the driven system are consistent with the location of the multicritical point that should also hold for $`E0`$. These results show that, for $`\rho =1/2`$, there is a gap in the critical temperature between the case $`E=0`$ (Ising model) and the limit $`E0`$ of the present model. Such a gap is expected to be even greater for $`\rho 1/2`$ because in this case the coexistence temperature of the Ising model is lower than $`T_c^I`$ while the coexistence temperature of the driven diffusive system has a lower bound given by the multicritical temperature. The existence of these temperature gaps dramatically reflects the anisotropy introduced by the driving field and the non-equilibrium nature of the studied model. A physical explanation of this observation remains as an open question.
In summary, the phase diagram of a DDS in the presence of an oscillatory driving field is determined for $`E=1`$ and $`E=\mathrm{}`$. We give strong evidence on the existence of a multicritical point and a critical temperature gap separating the cases $`E=0`$ from $`E0`$. To our best knowledge, these features have never been reported in the field of DDS.
Acknowledgments: This work was supported by CONICET, UNLP, ANPCyT and Fundación Antorchas (Argentina). We acknowledge useful discussions with B. Schmittmann.
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# Untitled Document
YITP-00-13 cond-mat/0004250
Thermodynamic limit of the Six-Vertex Model with Domain Wall Boundary Conditions
V. Korepin and P. Zinn-Justin
C.N. Yang Institute for Theoretical Physics
State University of New York at Stony Brook
Stony Brook, NY 11794–3840, USA
We address the question of the dependence of the bulk free energy on boundary conditions for the six vertex model. Here we compare the bulk free energy for periodic and domain wall boundary conditions. Using a determinant representation for the partition function with domain wall boundary conditions, we derive Toda differential equations and solve them asymptotically in order to extract the bulk free energy. We find that it is different and bears no simple relation with the free energy for periodic boundary conditions. The six vertex model with domain wall boundary conditions is closely related to algebraic combinatorics (alternating sign matrices). This implies new results for the weighted counting for large size alternating sign matrices. Finally we comment on the interpretation of our results, in particular in connection with domino tilings (dimers on a square lattice).
04/2000
1. Introduction
The six vertex model is an important model of classical statistical mechanics in two dimensions. The prototypical model is the ice model, which was solved by Lieb in 1967 by means of Bethe Ansatz, followed by several generalizations , , . The solution of the most general six vertex model was given by Sutherland in 1967. The bulk free energy was calculated in these papers for periodic boundary conditions (PBC). A detailed classification of the phases of the model can be found for example in the book (see also the more recent work on anti-periodic boundary conditions).
Earlier, in 1961 Kasteleyn, while studying dimer arrangements on a quadratic lattice, expressed doubts on the independence of the bulk free energy on boundary conditions . For more on dimer arrangements, see , and . Interest in this subject was renewed with recent work on domino tilings (which are equivalent to dimers on a square lattice) of an Aztec diamond , , demonstrating a strong effect of the boundary on a typical domino configuration (see also ). Dimers (or domino tilings) can be considered as a particular case of the six vertex model, and therefore a natural question is to investigate the effect of boundary conditions on the thermodynamic limit of the six vertex model.
Independently of this, new boundary conditions of the six-vertex model, the so-called domain wall boundary conditions (DWBC), were first introduced in 1982 (we shall define them in detail below). An important recursion relation for the partition function was discovered in this paper. Later these recursion relations helped to find a determinant representation for the partition function of the six vertex model with DWBC , . The determinant representation simplifies somewhat in the homogeneous case. In this case the partition function satisfy Toda differential equation . In this paper we use this differential equation in order to calculate the bulk free energy for domain wall boundary conditions.
Let us all mention that there is a one to one correspondence between arrow configurations in the six vertex model with DWBC and Alternating Sign Matrices (ASM) . This mapping was used in order to count the number of ASM. More on ASM can be found in and .
The plan of the paper is as follows. In Section 2 we define the six-vertex model with domain wall boundary conditions, and derive the determinant representation for the partition function. In Section 3 we derive Toda differential equation for the partition function. In Section 4 we consider the thermodynamic limit; we derive the explicit expression of the bulk free energy in the ferroelectric and disordered phases, and compare it with PBC. Finally, in Section 5 we comment on the connection of our results with other subjects (ASM, domino tilings, height model) and conclude this discussion in Section 6.
2. Determinant representation of the partition function of the six-vertex model
In this section we shall define the inhomogeneous six-vertex model with domain wall boundary conditions, and rewrite its partition function as a determinant. We will then particularize our formula to the homogeneous case.
Fig. 1: A configuration of the inhomogeneous six-vertex model with domain wall boundary conditions.
First we define the configurations of the model. They are given by assigning arrows to each edge of a $`N\times N`$ square lattice (see Fig. 1). The “domain wall” boundary conditions correspond to fixing the horizontal external arrows to be outgoing and the vertical external arrows to be incoming. The partition function is then obtained by summing over all possible configurations:
$$Z=\underset{\text{arrow configurations}}{}\underset{i,k=1}{\overset{N}{}}w_{ik}$$
where the statistical weights $`w_{ik}`$ are assigned to each vertex of the lattice. Since we are considering an inhomogeneous model, we need two sets of spectral parameters $`\{\lambda _i\}`$ and $`\{\mu _k\}`$ which are associated with the horizontal and vertical lines. The weight $`w_{ik}`$ depends on the arrow configuration around the vertex $`(i,k)`$ and is given by
$$w_{ik}=\{\begin{array}{cc}a(\lambda _i,\mu _k)\text{}\hfill & \\ b(\lambda _i,\mu _k)\text{}\hfill & \\ c(\lambda _i,\mu _k)\text{}\hfill & \end{array}$$
(all other weights are zero) where the functions $`a`$, $`b`$, $`c`$ are chosen as follows:
$$\begin{array}{cc}\hfill a(\lambda ,\mu )& =\mathrm{sinh}(\lambda \mu \gamma )\hfill \\ \hfill b(\lambda ,\mu )& =\mathrm{sinh}(\lambda \mu +\gamma )\hfill \\ \hfill c(\lambda ,\mu )& =\mathrm{sinh}(2\gamma )\hfill \end{array}$$
Here $`\gamma `$ is an anisotropy parameter which does not depend on the lattice site. The partition function is therefore a function of the $`2N`$ spectral parameters and we shall denote it by $`Z_N(\{\lambda _i\},\{\mu _k\})`$.
The model thus defined satisfies the following essential property (Yang–Baxter equation) shown on Fig. 2. The vertex with diagonal edges is assigned weights (the so-called $`R`$ matrix) which are the same as the usual weights, up to a shift of the difference of the spectral parameter. Here we shall not need the explicit expression of the $`R`$ matrix.
Fig. 2: Yang–Baxter equation. Summation over arrows of the internal edges is implied, whereas external arrows are fixed.
We shall now list the following four properties which determine entirely $`Z_N(\{\lambda _i\},\{\mu _k\})`$ and sketch their proof (for a detailed algebraic proof the reader is referred to ):
a) $`Z_1=\mathrm{sinh}(2\gamma )`$.
This is by definition.
b) $`Z_N(\{\lambda _i\},\{\mu _k\})`$ is a symmetric function of the $`\{\lambda _i\}`$ and of the $`\{\mu _k\}`$.
It is sufficient to prove that exchange of $`\mu _i`$ and $`\mu _{i+1}`$ (for any $`i`$) leaves the partition function unchanged. This can be obtained by repeated use of the Yang–Baxter property:
$$\begin{array}{cc}\hfill R_{}(\mu _i\mu _{i+1})& Z_N(\{\mathrm{}\mu _i,\mu _{i+1}\mathrm{}\})=\text{}=\text{}\hfill \\ \hfill =\mathrm{}& =\text{}=R_{}(\mu _i\mu _{i+1})Z_N(\{\mathrm{}\mu _{i+1},\mu _i\mathrm{}\})\hfill \end{array}$$
where $`R_{}=R_{}`$ is the appropriate entry of the $`R`$ matrix; and similarly for the $`\{\lambda _i\}`$.
c) $`Z_N(\{\lambda _i\},\{\mu _k\})=\mathrm{e}^{(N1)\lambda _i}P_{N1}(\mathrm{e}^{2\lambda _i})`$ where $`P_{N1}`$ is a polynomial of degree $`N1`$, and similarly for the $`\mu _k`$.
Let us choose one configuration. Then the only weights which depend on $`\lambda _i`$ are the $`N`$ weights on row $`i`$. Since the outgoing arrows are in opposite directions, at least one of the weights must be $`c`$. Therefore there are at most $`N1`$ weights $`a`$ and $`b`$, and the product of all weights is of the form $`\mathrm{e}^{(N1)\lambda _i}P_{N1}(\mathrm{e}^{2\lambda _i})`$. This property remains of course valid when we sum over all configurations.
d) $`Z_N(\{\lambda _i\},\{\mu _k\})`$ obeys the following recursion relation:
$$\begin{array}{cc}\hfill Z_N(\{\lambda _i\},\{\mu _k\})_{|\lambda _j\mu _l=\gamma }=& \mathrm{sinh}(2\gamma )\underset{\genfrac{}{}{0pt}{}{1kN}{kl}}{}\mathrm{sinh}(\lambda _j\mu _k+\gamma )\hfill \\ & \underset{\genfrac{}{}{0pt}{}{1iN}{ij}}{}\mathrm{sinh}(\lambda _i\mu _l+\gamma )Z_{N1}(\{\lambda _i\}_{ij},\{\mu _k\}_{kl})\hfill \end{array}$$
Because of property b), we can assume that $`j=l=1`$. Since $`\lambda _k\mu _l=\gamma `$ implies $`a(\lambda _j\mu _l)=0`$, by inspection all configurations with non-zero weights are of the form shown on Fig. 3. This immediately proves Eq. (2.1).
Fig. 3: Graphical proof of the recursion relation.
It is easy to see that the four properties a), b), c) and d) characterize entirely $`Z_N(\{\lambda _i\},\{\mu _k\})`$. This is enough to prove that $`Z_N(\{\lambda _i\},\{\mu _k\})`$ has the following determinant representation :
$$\begin{array}{cc}\hfill Z_N(\{\lambda _i\},\{\mu _k\})=& \frac{\underset{1i,kN}{}\mathrm{sinh}(\lambda _i\mu _k+\gamma )\mathrm{sinh}(\lambda _i\mu _k\gamma )}{_{1i<jN}\mathrm{sinh}(\lambda _i\lambda _j)_{1k<lN}\mathrm{sinh}(\mu _k\mu _l)}\hfill \\ & \underset{1i,kN}{det}\left[\frac{\mathrm{sinh}(2\gamma )}{\mathrm{sinh}(\lambda _i\mu _k+\gamma )\mathrm{sinh}(\lambda _i\mu _k\gamma )}\right]\hfill \end{array}$$
Indeed, one can check that this expression satisfies the four properties listed above.
The expression (2.1) might seem singular when two spectral parameters $`\lambda _i`$ and $`\lambda _j`$ coincide (and similarly for the $`\mu _k`$); but in fact the pole created by the factor $`\mathrm{sinh}(\lambda _i\lambda _j)`$ is compensated by the zero of the determinant due to the fact that two rows are identical. Therefore, particular care must be taken when considering the homogeneous limit where all the $`\lambda _i`$ are equal (and all the $`\mu _k`$). This limit was studied in detail in , and we shall simply summarize the result of the calculation. Let us call $`t`$ the common value of $`\lambda _i\mu _k`$ for all $`i`$ and $`k`$. When the $`\lambda _i`$ are close to one another one must Taylor expand the function
$$\varphi (t)\frac{\mathrm{sinh}(2\gamma )}{\mathrm{sinh}(t+\gamma )\mathrm{sinh}(t\gamma )}$$
which appears in the determinant. This leads to the following expression:
$$Z_N(t)=\frac{(\mathrm{sinh}(t+\gamma )\mathrm{sinh}(t\gamma ))^{N^2}}{\left(_{n=0}^{N1}n!\right)^2}\underset{1i,kN}{det}\left[\frac{\mathrm{d}^{i+k2}}{\mathrm{d}t^{i+k2}}\varphi (t)\right]$$
3. Determinant representation and Toda chain hierarchy
We shall now investigate the properties of the determinant which appears in Eq. (2.1), and for which we introduce the notation
$$\tau _N(t)=\underset{1i,kN}{det}\left[m_{i+k2}\right]$$
with
$$m_n=\frac{\mathrm{d}^n}{\mathrm{d}t^n}\varphi (t)$$
Let us write down the bilinear Hirota equation satisfied by the $`\tau _N`$. For completeness, we recall that they are a consequence of Jacobi’s determinant identity:
$$\text{}\text{}=\text{}\text{}\text{}\text{}$$
The large squares represent a given matrix, and the shaded regions are the sub-matrices whose determinants one must consider. Applying it to $`\tau _{N+1}`$ (up to a re-shuffling of the rows and columns), we find :
$$\tau _N\tau _N^{\prime \prime }\tau _N^{}{}_{}{}^{2}=\tau _{N+1}\tau _{N1}N1$$
where primes denote differentiation with respect to $`t`$. This is supplemented by the initial data: $`\tau _0=1`$ and $`\tau _1=\varphi `$. Equivalently, we have:
$$(\mathrm{log}\tau _N)^{\prime \prime }=\frac{\tau _{N+1}\tau _{N1}}{\tau _N^2}N1$$
which is the form of the equation that we shall use.
Note that if we introduce the combinations $`\mathrm{e}^{\phi _N}=\tau _N/\tau _{N1}`$, $`N1`$, Eq. (3.1) implies for the $`\phi _N`$:
$$\phi _N^{\prime \prime }=\mathrm{e}^{\phi _{N+1}\phi _N}\mathrm{e}^{\phi _N\phi _{N1}}N2$$
and $`\phi _1^{\prime \prime }=\mathrm{e}^{\phi _2\phi _1}`$. These are the usual Toda (semi-infinite) chain equations , , . Another possible form is
$$\psi _N^{\prime \prime }=\underset{M}{}C_{MN}\mathrm{e}^{\psi _M}N1$$
with $`\psi _N=\phi _{N+1}\phi _N`$ and $`C_{MN}`$ ($`M`$, $`N1`$) the Cartan matrix of the semi-infinite diagram $`A_{\mathrm{}}`$.
This suggests a connection with the Toda chain hierarchy , , , , , . Indeed, let us mention that given a Hänkel matrix $`(m_{i+k2})`$ – i.e. whose entries only depend on $`i+k`$ – the $`m_n`$ can be made to depend on a set of parameters $`\{t_q\}_{q1}`$ in such a way that the determinants $`\tau _N`$ become $`\tau `$-functions of the whole Toda (semi-infinite) chain hierarchy . Namely, one must choose
$$m_n(\{t_q\})=d\rho (x)x^n\mathrm{e}^{_{q1}t_qx^q}$$
where $`\mathrm{d}\rho (x)`$ is an arbitrary measure<sup>1</sup> This must be considered as a formal expression; e.g. the measure may not necessarily be smooth or positive. (in the matrix model context , the $`t_q`$ are the coefficients of the polynomial potential). Here, we are in the simplest situation where only one parameter $`t_1t`$ is allowed to evolve. We immediately check that Eq. (3.1) implies that $`m_n(t)=\frac{\mathrm{d}^n}{\mathrm{d}t^n}m_0(t)`$, which is consistent with Eq. (3.1).
4. The thermodynamic limit
We shall now consider the thermodynamic (i.e. large $`N`$) limit of the expression (2.1) in the various regimes of the six-vertex model. For that we shall use the Hirota equation in its form (3.1).
When $`N\mathrm{}`$ it is expected that the partition function behaves in the following way:
$$\mathrm{log}Z_N(t)=N^2F(t)+O(N)$$
where $`F(t)`$ is the bulk free energy (we shall always set the temperature $`k_BT=1`$). Our main goal is to compute explicitly $`F(t)`$.
Comparing the expected asymptotic (4.1) with the exact formula (2.1), we find that the determinant $`\tau _N`$ must be of the form
$$\tau _N=\left(\underset{n=0}{\overset{N1}{}}n!\right)^2\mathrm{e}^{N^2f(t)+O(N)}$$
where
$$f(t)=F(t)\mathrm{log}(\mathrm{sinh}(t+\gamma )\mathrm{sinh}(t\gamma ))$$
We now want to substitute the expansion (4.1) into the equation (3.1). For that we need to assume that the sub-dominant corrections to the bulk free energy vary slowly as a function of $`N`$; we shall discuss the validity of this assumption below. We then find that the expansion is consistent since both left and right hand sides of (3.1) turn out to be of order $`N^2`$. The resulting equation for $`f`$ is:
$$f^{\prime \prime }=\mathrm{e}^{2f}$$
This is an ordinary second order differential equation, which can be readily solved. The general solution depends on two parameters $`\alpha `$ and $`t_0`$:
$$\mathrm{e}^{f(t)}=\frac{\alpha }{\mathrm{sinh}(\alpha (tt_0))}$$
If the weights are chosen to be real, then the free energy should be real and this implies that $`\alpha `$ must be real or purely imaginary.
So far everything we have done was independent of the particular form of the function $`\varphi (t)`$ and therefore independent of $`\gamma `$. In order to fix the two constants in (4.1), we must now discuss separately the different regimes of the six-vertex model. Let us recall that the latter are usually distinguished by the value of the parameter (cf Eq. (8.3.21) of )
$$\mathrm{\Delta }=\frac{a^2+b^2c^2}{2ab}$$
The weights $`a`$, $`b`$, $`c`$ were defined in Eq. (2.1) (with $`\lambda \mu t`$). In this parameterization,
$$\mathrm{\Delta }=\mathrm{cosh}(2\gamma )$$
4.1. Ferro-electric phase: $`\mathrm{\Delta }>1`$
This corresponds to the parameters $`\gamma `$ and $`t`$ real; we recall that the weights are given by
$$a=\mathrm{sinh}(t\gamma )b=\mathrm{sinh}(t+\gamma )c=\mathrm{sinh}(2\gamma )$$
with $`|\gamma |<t`$. This is the so-called ferroelectric phase. In the case of periodic boundary conditions, it is known that the system is frozen in its ground state configuration, in which all arrows are aligned: if $`a>b`$ all arrows point up and to the right or down and to the left, whereas if $`b>a`$ they point up and to the left or down and to the right. The domain wall boundary conditions do not allow all arrows to be aligned: the ground state will instead take the form of Fig. 4. However at leading order in the large $`N`$ limit, this does not affect the free energy, and we expect to find the same result as for periodic boundary conditions.
Fig. 4: Ground state configuration of the ferroelectric phase (for $`b>a`$; the case $`a>b`$ is obtained by taking the mirror image).
Indeed, it is easy to see that the relevant solution of Eq. (4.1) is
$$\mathrm{e}^{f(t)}=\frac{1}{\mathrm{sinh}(t|\gamma |)}$$
and therefore the bulk free energy takes the form
$$\mathrm{e}^{F(t)}=\mathrm{sinh}(t+|\gamma |)=\mathrm{max}(a,b)$$
in agreement with the case of periodic boundary conditions.
4.2. Disordered phase: $`1<\mathrm{\Delta }<1`$
In this regime, it is customary to make the following redefinitions:
$$\gamma ^{}=i(\gamma +i\pi /2)$$
$$t^{}=i(t+i\pi /2)$$
and divide all the weights by $`i`$, so that they take the form:
$$a=\mathrm{sin}(\gamma t)b=\mathrm{sin}(\gamma +t)c=\mathrm{sin}(2\gamma )$$
and $`\mathrm{\Delta }=\mathrm{cos}(2\gamma )`$. Using symmetry considerations, one can always assume $`0<\gamma <\pi /2`$. We only consider the region $`|t|<\gamma `$ (where the weights are positive).
Taking into account these redefinitions, the partition function becomes:
$$Z_N(t)=\frac{(\mathrm{sin}(t+\gamma )\mathrm{sin}(t\gamma ))^{N^2}}{\left(_{n=0}^{N1}n!\right)^2}\underset{1i,kN}{det}\left[\frac{\mathrm{d}^{i+k2}}{\mathrm{d}t^{i+k2}}\varphi (t)\right]$$
with a redefined $`\varphi (t)=\mathrm{sin}(2\gamma )/(\mathrm{sin}(t\gamma )\mathrm{sin}(t+\gamma ))`$; the determinant $`\tau _N`$ still satisfies Eq. (3.1) and $`f(t)`$ defined by (4.1) is still a solution of Eq. (4.1).<sup>2</sup> Note that the sign is unchanged in Eq. (4.1); this is the combined effect of the “Wick rotation” of $`t`$ ($`tit`$) and of dividing all the weights by $`i`$ ($`\mathrm{e}^fi\mathrm{e}^f`$).
Let us mention that the partition function has been computed exactly at three particular values of the parameters: $`t=0`$, $`\gamma =\pi /6`$, $`\pi /4`$ and $`\pi /3`$. In all three cases the expansion (4.1) and the assumption of smoothness of the sub-dominant corrections (which is necessary to derive the ordinary differential equation (4.1)) can be checked exactly. We have also checked it numerically for a variety of values of $`t`$ and $`\gamma `$.
We must now select the appropriate solution (of the form (4.1)) of the Eq. (4.1). Let us first assume that $`|t|<\gamma `$ (this is the only physical region, i.e. where all the weights are positive). It is easy to check that $`f(t)`$ must be an even function of $`t`$. The only even solution of Eq. (4.1) is
$$\mathrm{e}^{f(t)}=\frac{\alpha }{\mathrm{cos}(\alpha t)}$$
where $`\alpha `$ remains to be determined. Note that this implies for $`F(t)`$
$$\mathrm{e}^{F(t)}=\mathrm{sinh}(\gamma t)\mathrm{sinh}(\gamma +t)\frac{\alpha }{\mathrm{cos}(\alpha t)}$$
We must then use the boundary condition given by $`|t|=\pm \gamma `$. At these values one can compute directly $`Z_N(t)`$. Indeed the only non-zero configurations are of the form of Fig. 4, and we find
$$Z_N(t=\pm \gamma )=\mathrm{sin}(2\gamma )^{N^2}$$
and therefore $`\mathrm{e}^{F(t)}=\mathrm{sin}(2\gamma )`$. Since the prefactor in (4.1) vanishes when $`|t|=\gamma `$, we conclude that $`\alpha `$ must be chosen in such a way that $`\mathrm{cos}(\alpha t)`$ is non-zero for $`|t|<\gamma `$, but vanishes as $`|t|=\gamma `$. This uniquely determines $`\alpha `$ to be: $`\alpha =\frac{\pi }{2\gamma }`$. We obtain the final expression
$$\mathrm{e}^{F(t)}=\mathrm{sin}(\gamma t)\mathrm{sin}(\gamma +t)\frac{\pi /2\gamma }{\mathrm{cos}(\pi t/2\gamma )}$$
As a consistency check, one takes the limit $`t\pm \gamma `$ and finds $`\mathrm{e}^{F(t)}=\mathrm{sin}(2\gamma )`$, as it should be. Also, note that for $`\gamma =\pi /4`$, where the partition function is known and independent of $`t`$, one finds indeed that $`\mathrm{e}^{F(t)}=1`$.
For further checks, let us set $`t=0`$; a more standard normalization of the weights is then
$$a=b=1c=2\mathrm{cos}\gamma $$
and the bulk free energy becomes
$$\mathrm{e}^F=\frac{\pi }{2}\frac{\mathrm{sin}\gamma }{\gamma }$$
At $`\gamma =\pi /6`$, $`\pi /4`$, $`\pi /3`$, the values predicted by (4.1) coincide with the large $`N`$ limit of the expressions of . Also, this fits perfectly with some numerical computations of the determinant we have performed.
We can compute the bulk energy (energy per unit site), which turns out to be
$$E=(\mathrm{cot}\gamma 1/\gamma )\mathrm{cot}\gamma \mathrm{log}(2\mathrm{cos}\gamma )$$
Fig. 5 shows the comparison with Monte-Carlo simulations. The agreement is also very good.
Fig. 5: Energy $`E`$ as a function of the anisotropy $`\mathrm{\Delta }`$. The curve is given by Eq. (4.1), whereas the diamonds are the results of Monte-Carlo simulations on lattices of size $`N=64`$.
Finally let us mention that there seems to be no simple relation between the PBC and DWBC bulk free energies: from an analytic point of view, the DWBC free energy is an elementary function, whereas the PBC free energy is given by a non-trivial integral. Furthermore, the DWBC free energy is always greater then the PBC free energy, even at infinite temperature ($`\mathrm{\Delta }=1/2`$), see Fig. 6.
Fig. 6: Bulk free energies for PBC and DWBC as a function of $`\mathrm{\Delta }`$.
4.3. Anti-ferroelectric phase: $`\mathrm{\Delta }<1`$
In this phase, the smoothness assumption of the sub-dominant corrections to the bulk free energy is not satisfied, as can be clearly seen numerically. The ratio $`Z_{N+1}Z_{N1}/Z_N^2`$ does not converge in the large $`N`$ limit but instead has a pseudo-periodic behavior reminiscent of the one-matrix model with several cuts , and slightly more sophisticated methods are needed to analyze the large $`N`$ limit. We leave this to a future publication.
4.4. Phase transition at $`\mathrm{\Delta }=1`$
If the Boltzmann weights depend on a parameter (e.g. temperature), it is known that with periodic boundary conditions, the system undergoes phase transitions as $`\mathrm{\Delta }`$ crosses $`\pm 1`$. Let us use the expressions of the bulk free energy found above to clarify what happens in the case of domain wall boundary conditions.
Here we shall consider the transition from ferroelectric (low temperature) to disordered (high temperature) regime, that is from $`\mathrm{\Delta }>1`$ to $`\mathrm{\Delta }<1`$. The parameter that plays the role of deviation from criticality $`TT_c`$ can be defined as
$$TT_c1\mathrm{\Delta }$$
We assume that $`b>a`$ (the case $`a>b`$ can be treated similarly), and re-scale the weights so that $`b=1`$. With this convention, we simply have in the ferroelectric phase:
$$\mathrm{e}^F=1\mathrm{\Delta }>1$$
(cf Eq. (4.1)). Let us now consider $`\mathrm{\Delta }1^{}`$. The weights are
$$a=\frac{\mathrm{sin}(\gamma t)}{\mathrm{sin}(\gamma +t)}b=1c=\frac{\mathrm{sin}(2\gamma )}{\mathrm{sin}(\gamma +t)}$$
with $`\gamma =\pi /2+ϵ`$, $`t=\pi /2+ϵx`$; $`x`$ must be kept fixed as $`ϵ0`$. Note that $`\mathrm{\Delta }=\mathrm{cos}(2ϵ)`$, so that
$$TT_cϵ^2$$
Expanding the expression (4.1) for the free energy, we obtain:
$$\mathrm{e}^F=1\frac{2(x1)^2}{3\pi }ϵ^3+O(ϵ^4)$$
Comparing (4.1) and (4.1), we find a second order phase transition, with a singular part $`(TT_c)^{3/2}`$ corresponding to a critical exponent $`\alpha =1/2`$. This is to be contrasted with the first order phase transition that occurs in the case of PBC. Let us however emphasize that the difference of orders of the phase transitions is not that significant, since the phase transition is somewhat special (in the case of PBC, the correlation length jumps from zero for $`\mathrm{\Delta }>1`$ to infinity for $`\mathrm{\Delta }<1`$).
5. Some equivalences
We shall now review some alternative interpretations of the partition function of the six-vertex model with domain wall boundary conditions; these equivalent formulations will shed some light on the property of dependance on boundary conditions that was found.
5.1. Alternating sign matrices
Six-vertex model arrow configurations with domain wall boundary conditions on a $`N\times N`$ lattice are in one-to-one correspondence with alternating sign matrices (ASM) of size $`N`$, that is square matrices with entries $`0`$ or $`\pm 1`$ such that each row and column has an alternating sequence of $`+1`$ and $`1`$ (zeroes excluded) starting and ending with a $`+1`$. Recalling that there are 6 weights which we shall call $`a_1`$, $`a_2`$, $`b_1`$, $`b_2`$, $`c_1`$, $`c_2`$ in the order shown in Eq. (2.1), the correspondence goes as follows: given a six-vertex configuration, assign a $`0`$ to each vertex $`a`$ or $`b`$ and $`+1`$ (resp. $`1`$) to each vertex $`c_1`$ (resp. $`c_2`$). One can show that this map is bijective, and therefore, the number of ASM is exactly equal to the partition function considered before with $`a=b=c=1`$. For our purposes, let us define a refined counting of ASM ($`x`$-enumeration in the language of ) by assigning a weight $`x`$ to each entry $`1`$ of the ASM. The resulting quantity $`A(N,x)`$ is still related to the six-vertex model; indeed, one can easily show that
$$A(N,x)=x^{N/2}Z_N(a=b=1,c=\sqrt{x})$$
If $`0x4`$, one can set $`x=4\mathrm{cos}^2\gamma `$ and the weights are of the form (4.1). One can then prove (extending slightly the asymptotic expansion found in 4.2) that
$$\mathrm{log}A(N,x)=N^2\mathrm{log}\left[\frac{\pi }{2}\frac{\mathrm{sin}\gamma }{\gamma }\right]\frac{N}{2}\mathrm{log}x+O(\mathrm{log}N)x[0,4]$$
Though this equivalence does not directly provide any useful insight into the issue adressed in this paper, the result (5.1) itself might be of some mathematical interest.
5.2. Tilings of the Aztec diamond
A more illuminating equivalence is that of domino tilings (i.e. dimers on a square lattice in a dual description) and the six-vertex model at $`\mathrm{\Delta }=0`$ – both models are well-known to describe essentially one Dirac fermion. This is illustrated on Fig. 7. Since each vertex of type $`c_1`$ has $`2`$ possible corresponding domino tiling configurations, one must assign it a Boltzmann weight of $`2`$ in order to count each domino tiling exactly once; however with most boundary conditions there are as many vertices of type $`c_1`$ and $`c_2`$, and therefore one can give them both a weight of $`\sqrt{2}`$ instead, which leads to the values $`a=b=1`$, $`c=\sqrt{2}`$ of the parameters.
Fig. 7: Correspondence between vertices of the six-vertex model and small patches of a domino tiling.
The more precise statement is that the number of domino tilings of the Aztec diamond (see ) is equal (up to a small known prefactor) to the partition function of the six-vertex model with domain wall boundary conditions at $`a=b=1`$, $`c=\sqrt{2}`$, see Fig. 8.
Fig. 8: a) A configuration of the six-vertex model with DWBC, and b) one possible corresponding tiling of the Aztec diamond.
Since this a local correspondence of configurations it entends to all correlation functions. Also, introducing some weights for the local tiling patterns amounts to changing the weights $`a`$, $`b`$, $`c`$, but always in such a way that $`\mathrm{\Delta }`$ remains zero.
These tilings have been an object of interest for mathematicians, see in particular , . The “arctic circle theorem” shows that as the size of the system grows large, the domino configurations become frozen outside the circle inscribed inside the diamond, and remain disordered but still heterogeneous (i.e. non translationally invariant) inside the circle. These statements have a straightforward equivalent in the six-vertex language: we expect the one-point functions of the six-vertex model at $`\mathrm{\Delta }=0`$ with DWBC to be non-constant (following a similar pattern as the tilings), and presumably a similar behavior at $`\mathrm{\Delta }0`$. This gives a qualitative understanding of the dependence of the bulk free energy on the boundary conditions.
5.3. Height model
The general eight-vertex is well-known to be equivalent to a class of height models (SOS/RSOS model). In the case of the six-vertex model, there is a particularly simple equivalence which goes as follows: given a six-vertex configuration, integers are assigned to each face of the lattice in such a way that going from one face to a neighboring face, the number is increased by one if the arrow in between goes right (and so, is decreased by one if it goes left). Conservation of the arrows ensures consistency of this procedure. The Boltzman weights of the model are simply the weights of the original six-vertex model expressed in terms of the new height variables.
This equivalence is particularly interesting because it gives us a simple intuitive explanation of the lack of thermodynamic limit due to boundary conditions; this was studied in detail and proven rigorously in the case of tilings ($`\mathrm{\Delta }=0`$), see . Let us consider the domain wall boundary conditions and translate these into the language of our height model. They are fixed boundary conditions for the heights, of the form:
$$\begin{array}{ccccc}0& 1& \mathrm{}& N1& N\\ 1& & & & N1\\ \mathrm{}& & & & \mathrm{}\\ N1& & & & 1\\ N& N1& \mathrm{}& 1& 0\end{array}$$
where we have fixed arbitrarily the upper left height to be zero.
In the thermodynamic limit $`N\mathrm{}`$, let us define the rescaled coordinates on the square lattice to be $`x=ka`$, $`y=ia`$ where $`a1/N`$ is the lattice spacing. The heights $`h_{i,j}`$ are supposed to renormalize, according to standard lore, to a free massless bosonic field:
$$h_{i,j}\varphi (x,y)$$
However, it is reasonable to assume that in order to have a proper thermodynamic limit, the boundary conditions must be well-defined in terms of the limiting field $`\varphi `$. In the case of DWBC, one finds that the boundary conditions become $`\varphi (x,0)=x/a`$ etc, which do not have a limit as $`a0`$; in particular, the variations of $`\varphi `$ on the boundary diverge. More generally, we can conjecture that only the boundary conditions such that the variation of the function $`\varphi `$ on the boundary can remain bounded will lead to the usual thermodynamic limit. This is essentially what is proven in in the case $`\mathrm{\Delta }=0`$.
6. Conclusion
In this work, we have computed explicitly the large $`N`$ asymptotic behavior of a $`N\times N`$ determinant which plays the role of partition function of the six-vertex model with domain wall boundary conditions. This gives rise to particularly simple expressions for the bulk free energy of this model (Eqs. (4.1) and (4.1)). One important question is to physically interpret the discrepancy of the bulk free energy found when comparing domain wall and periodic boundary conditions of the six-vertex model, which is somewhat contrary to standard lore on the thermodynamic limit of statistical models. Some clues were given in the previous section, where various equivalences were discussed. In particular it was shown how “generic” fixed boundary conditions for the six-vertex model do not lead to a well-defined thermodynamic limit. It would be useful to make these arguments more rigorous. Also, it would be most interesting to find a more quantitative description of the non-translational invariance created by the boundary conditions, and in particular to prove a “generalized arctic circle theorem” for any value of the parameter $`\mathrm{\Delta }`$ of the six-vertex model.
Acknowledgements
It is a pleasure to acknowledge stimulating discussions with E. Lieb, F.Y. Wu, S. Ruijsenaars (who informed us that he had similar results) and R. Behrends (who participated in the early stages of this project). This work was supported by the National Science Foundation under grant number PHY-9605226 (V.K.).
References
relax E. Lieb, Phys. Rev. Lett. 18 (1967), 692. relax E. Lieb, Phys. Rev. Lett. 18 (1967), 1046. relax E. Lieb, Phys. Rev. Lett. 19 (1967), 108. relax E. Lieb, Phys. Rev. 162 (1967), 162. relax B. Sutherland, PRL 19 (1967), 103. relax R.J. Baxter, Exactly Solved Models in Statistical Mechanics (San Diego, CA: Academic). relax M.T. Batchelor, R.J. Baxter, M.J. O’Rourke and C.M. Yung, J. Phys. A28 (1995), 2759. relax P.W. Kasteleyn, Physica 27 (1960), 1209. relax M.E. Fisher, Phys. Rev. 124 (1961), 1664. relax W.T. Lu and F.Y. Wu, Phys. lett. A 259 (1999), 108. relax W. Jockush, J. Propp and P. Shor, preprint math.CO/9801068. relax H. Cohn, N. Elkies and J. Propp, Duke Math. Journal 85 (1996), 117. relax R. Kenyon, The planar dimer model with boundary: a survey, preprint (http://topo.math.u-psud.fr/$``$kenyon/papers.html). relax V.E. Korepin, Commun. Math. Phys 86 (1982), 391. relax A.G. Izergin, Sov. Phys. Dokl. 32 (1987), 878. relax A.G. Izergin, D.A. Coker and V.E. Korepin, J. Phys. A 25 (1992), 4315. relax K. Sogo, Journ. Phys. Soc. Japan 62, 6 (1993), 1887. relax N. Elkies, G. Kuperberg, M. Larsen and J. Propp, Journal of Algebraic Combinatorics 1 (1922), 111; 219. relax D.M. Bressoud, Proofs and Confirmations: The Story of the Alternating Sign Matrix Conjecture, Cambridge University Press, Cambridge, 1999 relax D. Bressoud and J. Propp, Notices of the AMS June/July (1999), 637. relax M. Toda, Journ. Phys. Soc. Japan 22 (1967), 431; Prog. Theor. Phys. Suppl. 45 (1970), 174. relax R. Hirota, Journ. Phys. Soc. Japan 56 (1987), 4285. relax H. Flaschka, Phys. Rev. B 9 (1974); Prog. Theor. Phys. 51 (1974), 703. relax K. Ueno and K. Takasaki, Adv. Studies in Pure Math. 4 (1984), 1. relax M. Adler and P. van Moerbeke, Duke Math. Journal 80 (1995), 863 (preprint solv-int/9706010); preprint math.CO/9912143. relax O. Lipan, P.B. Wiegmann, A. Zabrodin, preprint solv-int/9704015. relax P.B. Wiegmann, A. Zabrodin, preprint hep-th/9909147. relax I. Krichever, O. Lipan, P. Wiegmann and A. Zabrodin, preprint hep-th/9604080. relax A. Gerasimov, A. Marshakov, A. Mironov, A. Morozov and A. Orlov, Nucl. Phys. B 357 (1991), 565. relax G. Kuperberg, Internat. Math. Res. Notices 3 (1996), 139 (preprint math.CO/9712207). relax P. Deift, T. Kriecherbauer, K.T-R. McLaughlin, S. Venakides and X. Zhou, Commun. on Pure and Applied Math. 52 (1999), 1491.
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# QUANTUM FEEDBACK FOR PROTECTION OF SCHRÖDINGER CAT STATES
## 1 Introduction
One of the most fundamental issues in quantum theory is how the classical macroscopic world emerges from the quantum substrate. This question is also an important point in the interpretation of quantum mechanics and it is still the subject of an intense debate . The most striking example of this problem is given by the possibility, opened by quantum mechanics, of having linear superpositions of macroscopically distinguishable states, the so-called “Schrödinger-cat” states. Such paradoxical states are very sensitive to decoherence, i.e., the rapid transformation of these linear superpositions into the corresponding classical statistical mixture, caused by the unavoidable entanglement of the system with uncontrolled degrees of freedom of the environment . The decoherence time depends on the form of system-environment interaction but, in most cases, it is inversely proportional to the squared “distance” between the two states of the superposition . It is then clear, that for macroscopically distinguishable states, the decoherence process becomes thus practically instantaneous . Decoherence is therefore experimentally accessible only in the mesoscopic domain. In this case, one is able to monitor the progressive emergence of classical properties from the quantum ones. In this context, an important achievement has been obtained by Monroe et al. , who have prepared a trapped $`{}_{}{}^{9}\mathrm{Be}_{}^{+}`$ ion in a superposition of spatially separated coherent states and detected the quantum coherence between the two localized states. However, the decoherence of the superposition state has not been observed in this experiment. The progressive decoherence of a mesoscopic Schrödinger cat has been monitored for the first time in the experiment of Brune et al, where the linear superposition of two coherent states of the electromagnetic field in a cavity with classically distinct phases has been generated and detected.
Recently, the field of quantum information theory has undergone an impressive development, and the study of decoherence has then become important not only from a fundamental, but also from a more practical point of view. All the quantum information processing applications rely on the possibility of performing unitary transformations on a system of $`N`$ quantum bits (qubits), whose decoherence has to be made as small as possible. As a consequence, decoherence control is now a rapidly expanding field of investigation. In this respect, quantum error correction codes have been developed in which the entangled superposition state of $`N`$ qubits is “encoded” in a larger number of qubits. Assuming that only a fraction of qubits decoheres, it is then possible to reconstruct the original state with a suitable decoding procedure. These codes always require the entanglement of a large number of qubits, and will become practical only if quantum networks of tens of qubits become available. Up to now, the polarization states of three photons have been entangled at most . Entangled states of two Rydberg atoms or of two trapped ions at most can be generated. Therefore, in the present experimental situation, it is more realistic to study complementary and more “physical” ways to deal with decoherence, based on the explicit knowledge of the specific process causing decoherence, which could be applied with very few degrees of freedom. This is possible, in particular, in quantum optics, when information is encoded in the quantum states of an electromagnetic mode (see for example ). In this case decoherence is caused by photon leakage. It is then possible to develop experimental schemes able to face photon leakage and the associated decoherence.
We have already shown in some recent papers that a possible way to control decoherence in optical cavities is given by appropriately designed feedback schemes. Refs. show that a feedback scheme based on the continuous homodyne measurement of an optical cavity mode is able to increase the decoherence time of a superposition state. In Ref. a feedback scheme based on continuous photodetection and the injection of appropriately prepared atoms has been considered. This scheme, in the limit of very good detection efficiency, is able to obtain a significant “protection” of a generic quantum state in a cavity. In this photodetection-mediated scheme has been adapted to the microwave experiment of Ref. in which photodetectors cannot be used. The cavity state can only be indirectly inferred from measurements performed on probe atoms which have interacted with the cavity mode. Under ideal conditions, this adaptation to the microwave cavity case leads to a significant increase of the lifetime of the Schrödinger cat generated in . However, this scheme suffers from two important limitations, making it very inefficient when applied under the actual experimental situation. It first requires the preparation of samples containing exactly one Rydberg atom sent through the apparatus. Up to now, the experimental techniques allow only to prepare a sample containing a random atom number, with a Poisson statistics. Two-atom events are excluded only at the expense of a low average atom number, lengthening the feedback loop cycletime . The original scheme requires also a near unity atomic detection efficiency, which is extremely difficult to achieve even with the foreseeable improvements of the experimental apparatus.
Here we propose a significant improvement of the microwave feedback scheme described in . This new version, using a direct transmission of the quantum information from the probe to the feedback atom, does not require a large detection efficiency, removing one of the main difficulties of the previous design. It however also requires sub-poissonian atom statistics. We show briefly how such atomic packets could be in principle prepared with standard laser techniques. Finally, our scheme improves the efficiency of the feedback photon injection in the cavity by using an adiabatic rapid passage.
## 2 Detection-mediated feedback
In this section we briefly review the original “stroboscopic” feedback scheme for microwave cavities proposed in . This proposal is based on a very simple idea: whenever the cavity looses a photon, a feedback loop supplies the cavity mode with another photon, through the injection of an appropriately prepared atom. However, since there are no good enough photodetectors for microwaves, one has to find an indirect way to check if the high-Q microwave cavity has lost a photon or not. In the experiment of Brune et al. , information on the cavity field state is obtained by detecting the state of a circular Rydberg atom which has dispersively interacted with the superconducting microwave cavity. This provides an “instantaneous” measurement of the cavity field and suggests that continuous photodetection can be replaced by a series of repeated measurements, performed by non-resonant atoms regularly crossing the high-Q cavity.
The experimental scheme of the stroboscopic feedback loop is a simple modification of the scheme employed in Ref. . The relevant levels of the velocity-selected atoms are two adjacent circular Rydberg states with principal quantum numbers $`n=50`$ and $`n=51`$ (denoted by $`|g`$ and $`|e`$, respectively) and a very long lifetime ($`30`$ ms). The high-Q superconducting cavity is sandwiched between two low-Q cavities $`R_1`$ and $`R_2`$, in which classical microwave fields resonant with the transition between $`|e`$ and $`|g`$ can be applied.
The high-Q cavity $`C`$ is instead slightly off-resonance with respect to the $`eg`$ transition, with a detuning $`\delta =\omega \omega _{eg}`$, where $`\omega `$ is the cavity mode frequency and $`\omega _{eg}=(E_eE_g)/\mathrm{}`$. The Hamiltonian of the atom-microwave cavity mode system is the Jaynes-Cummings Hamiltonian,
$$H_{JC}=E_e|ee|+E_g|gg|+\mathrm{}\omega a^{}a+\mathrm{}\mathrm{\Omega }\left(|eg|a+|ge|a^{}\right),$$
(1)
where $`\mathrm{\Omega }`$ is the vacuum Rabi coupling between the atomic dipole on the $`eg`$ transition and the cavity mode. In the off-resonant case and perturbative limit $`\mathrm{\Omega }\delta `$, the Hamiltonian (1) assumes the dispersive form
$$H_{disp}=\mathrm{}\frac{\mathrm{\Omega }^2}{\delta }\left(|gg|a^{}a|ee|a^{}a\right).$$
(2)
A linear superposition state of two coherent states with opposite phases is generated when the cavity mode is initially in a coherent state $`|\alpha `$ and the Rydberg atom, which is initially prepared in the excited level $`|e`$, is subjected to a $`\pi /2`$ pulse both in $`R_1`$ and in $`R_2`$. In fact, when the atom has left the cavity $`R_2`$, the joint state of the atom-cavity system becomes the entangled state
$$|\psi _{atom+field}=\frac{1}{\sqrt{2}}(|e(|\alpha e^{i\varphi }|\alpha e^{i\varphi })+|g(|\alpha e^{i\varphi }+|\alpha e^{i\varphi })),$$
(3)
where $`\varphi =\mathrm{\Omega }^2t_{int}/\delta `$ and $`t_{int}`$ is the interaction time in $`C`$. A Schrödinger-cat state is then conditionally generated in the microwave cavity as soon as one of the two circular atomic states is detected.
As it was shown in Ref. , the stroboscopic feedback scheme works only for Schrödinger cat states with a definite parity, i.e. even or odd cat states, and therefore we shall restrict to $`\varphi =\pi /2`$ from now on. In fact, when the cavity field initial state is a generic density matrix $`\rho `$, the state of the probe atom-field system after the two $`\pi /2`$ pulses and the $`\varphi =\pi /2`$ conditional phase-shift can be written as
$$\rho _{atom+field}=|ee|\rho _e+|gg|\rho _g+|eg|\rho _++|ge|\rho _{},$$
(4)
where
$`\rho _e`$ $`=`$ $`P_{odd}\rho P_{odd}`$ (5)
$`\rho _g`$ $`=`$ $`P_{even}\rho P_{even},`$ (6)
are the projections of the cavity field state onto the subspace with an odd and even number of photons, respectively, and the operators $`\rho _\pm `$ (whose expression is not relevant here) are given in . Eq. (4) shows that there is a perfect correlation between the atomic state and the cavity field parity, which is the first step in an optimal quantum non demolition measurement of the photon number . It is possible to prove that this perfect correlation between the atomic state and a cavity mode property holds only in the case of an exact $`\varphi =\pi /2`$-phase shift sandwiched by two classical $`\pi /2`$ pulses in cavities $`R_1`$ and $`R_2`$ . Moreover, the entangled state of Eq. (4) allows to understand how it is possible to check if the microwave cavity $`C`$ has lost a photon or not and therefore to trigger the feedback loop, using atomic state detection only. The detection of $`e`$ or $`g`$ determines the parity of the field and, provided that the probe atomic pulses are frequent enough, indicates whether a microwave photon has left $`C`$ or not. In fact, let us consider for example the case in which an odd cat state is generated (first atom detected in $`e`$): a probe atom detected in state $`e`$ means that the cavity field has remained in the odd subspace. The cavity has therefore lost an even number of photons. If the time interval $`\tau _{pr}`$ between the two atomic pulses is much smaller than the cavity decay time $`\gamma ^1`$, $`\gamma \tau _{pr}1`$, the probability of loosing two or more photons is negligible and this detection of the probe atom in $`e`$ means that no photon has leaked out from the high-Q cavity $`C`$. On the contrary, when the probe atom is detected in $`g`$, the cavity mode state is projected into the even subspace. The cavity has then lost an odd number of photons. Again, in the limit of enough closely spaced sequence of probe atoms, $`\gamma \tau _{pr}1`$, the probability of loosing three or more photons is negligible. A detection in $`g`$ means that one photon has exited the cavity. Therefore, for achieving a good protection of the initial odd cat state, the feedback loop has to supply the superconducting cavity with a photon whenever the probe atom is detected in $`g`$, while feedback must not act when the atom is detected in the $`e`$ state.
In Ref. it has been proposed to realize this feedback loop with a switch connecting the $`g`$ state field-ionization detector with a second atomic injector, sending an atom in the excited state $`e`$ into the high-Q cavity. The feedback atom is put in resonance with the cavity mode by another switch turning on an electric field in the cavity $`C`$ when the atom enters it, so that the level $`e`$ is Stark-shifted into resonance with the cavity mode.
As it is shown in Ref. , if the probe atomic pulses are sufficiently frequent, this stroboscopic feedback scheme becomes extremely efficient and one gets a good preservation of an initial Schrödinger-cat state. However, if we consider the adaptation of this scheme to the present experimental apparatus of Ref. , we see that it suffers from two main limitations, which significantly decrease its efficiency. First of all the scheme is limited by the non-unit efficiency of the atomic state detectors ($`\eta _{det}0.4`$), since the feedback loop is triggered only when the $`g`$-detector clicks. Most importantly, the above scheme assumes one has perfect “atomic guns”, i.e. the possibility of having probe and feedback atomic pulses with exactly one atom. This is not experimentally achieved up to now. The actual experiment has been performed using atomic pulses with a probability of having exactly one atom $`p_10.2`$, close to the mean atom number in the sample. This low mean atom number has been chosen to minimize two–atom events. In this experimental situation, the proposed stroboscopic feedback scheme would have an effective efficiency $`\eta _{eff}=\eta _{det}p_1^20.016`$, too low to get an appreciable protection of the Schrödinger cat state. In the next section we show how this scheme may be improved and adapted to the experimental apparatus employed in Ref. .
## 3 The new automatic feedback scheme
From the above discussion, it is clear that the limitations due to the non-unit efficiency of the atomic detectors could be avoided if we eliminate the measurement step in the feedback loop and replace it with an “automatized” mechanism preparing the correct feedback atom whenever needed. This mechanism can be provided by an appropriate conditional quantum dynamics, and this, in turn, may be provided by a second high-Q microwave cavity $`C^{}`$, similar to $`C`$, replacing the atomic detectors, crossed by the probe atom first and by the feedback atom soon later, as described in Fig. 1.
The cavity $`C^{}`$ is resonant with the transition between an auxiliary circular state $`i`$ (an immediately lower circular Rydberg state), and level $`g`$. The interaction times have to be set so that both the probe and the feedback atom experience a $`\pi `$ pulse when they cross the empty cavity $`C^{}`$ in state $`g`$ (or when they enter in state $`i`$ with one photon in $`C^{}`$). This interaction copies the state of the probe atom onto the feedback atom, and thus removes any need for a unit detection efficiency.
This fine tuning of the interaction times to achieve the $`\pi `$-spontaneous emission pulse condition can be obtained applying through the superconducting mirrors of $`C^{}`$ appropriately shaped Stark-shift electric fields which put the atoms in resonance with the cavity mode in $`C^{}`$ only for the desired time. In this way, since $`C^{}`$ is initially in the vacuum state, when the probe atom crosses $`C^{}`$ one has
$`|e_p|0_C^{}`$ $``$ $`|e_p|0_C^{},`$ (7)
$`|g_p|0_C^{}`$ $``$ $`|i_p|1_C^{}.`$ (8)
Soon later a feedback atom enters $`C^{}`$ in the state $`|i_f`$ and one has
$`|i_f|0_C^{}`$ $``$ $`|i_f|0_C^{},`$ (9)
$`|i_f|1_C^{}`$ $``$ $`|g_f|0_C^{}.`$ (10)
In this way the cavity $`C^{}`$ is always left disentangled in the vacuum state. The feedback atom exiting $`C^{}`$ in $`|g`$ can be promoted to $`|e`$ before entering $`C`$, as required by the feedback scheme, by subjecting it to a $`\pi `$ pulse in the classical cavity $`R_2`$ (see Fig. 1). Therefore, the conditional dynamics provided by $`C^{}`$ eliminates any limitation associated to the measurement and leads to an “automatic feedback” scheme with unit efficiency in principle. Also the second important limitation of the stroboscopic feedback scheme of —i.e. that it requires exactly one probe and one feedback atom per loop—can be circumvented: A better control of the atom number, providing single atom events with a high probability, could be achieved by a modification of the Rydberg atoms preparation technique, in such a way that it is only triggered when the fluorescence detection signal (see Fig. 1) provides evidence of only one atom in the beam, implementing an atom counter.
Instead of preparing a random atom number at a given time, one thus prepares with a high probability a single Rydberg atom after a random delay. However, after a full quantum mechanical calculation and lengthy algebra, it is possible to determine the map of a generic feedback cycle, that is, the transformation connecting the states of the cavity field in $`C`$ soon after the passage of two successive feedback atoms in $`C`$, which also takes into account the non-unit efficiency of the Rydberg state preparation. This map, which is reported elsewhere , allows us to study the dynamics of the Schrödinger-cat state in the presence of feedback, and to compare it with the corresponding dynamics in absence of feedback.
In Fig. 2 we show the Wigner function of the initial odd cat state (top) and its dynamics in presence (left) and in absence (right) of feedback. The comparison between the two performances is striking: in absence of feedback the Wigner function becomes quickly positive definite, while in the presence of feedback the quantum aspects of the state remain well visible for many decoherence times.
## 4 Conclusions
In this paper we have proposed a method to significantly increase the “lifetime” of a Schrödinger cat state of a microwave cavity mode. However, as it can be easily expected, most of the techniques presented here could be applied to the case of a generic quantum state of a cavity mode (see also Ref.). After the first experimental evidences of decoherence mechanisms, decoherence control is an expanding field in quantum physics. An experimental realization of this realistic feedback scheme would be an important step in this direction.
## Acknowledgments
This work has been partially supported by INFM (through the 1997 Advanced Research Project “CAT”), by the European Union in the framework of the TMR Network “Microlasers and Cavity QED” and by MURST under the “Cofinanziamento 1997”.
## References
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# The asymptotic behavior of the colored Jones function of a knot and its volume
## 1. Introduction
The Jones polynomial $`V(L;t)[t^{1/2},t^{1/2}]`$ was first introduced by V. Jones as a link invariant which satisfies the following recursive formula:
$$\{\begin{array}{cc}V(O;t)=1,\hfill & \\ tV(L_+;t)t^1V(L_{};t)=(t^{1/2}t^{1/2})V(L_0;t),\hfill & \end{array}$$
where $`O`$ is the trivial knot and $`(L_+,L_{},L_0)`$ is a usual skein triple of links. (The original version uses a different normalization.) Let $`V_N(L;t)`$ denote the colored Jones polynomial, colored with the $`N`$-dimensional irreducible representation of $`sl_2()`$, which can be defined by using an enhanced Yang-Baxter operator . We use the normalization for $`V_N(L;t)`$ such that $`V_2(L;t)=V(L;t)`$ (Note that $`V_N(O;t)=1`$.) We also denote $`V_N(L;\mathrm{exp}(2\pi \sqrt{1}/N))`$ by $`J_N(L)`$ and call it the colored Jones function.
In we proved that $`J_N(L)`$ is the same as Kashaev’s invariant and generalized his conjecture to
###### Conjecture 1.1 (Volume Conjecture).
$$\frac{2\pi }{v_3}\underset{N\mathrm{}}{lim}\frac{\mathrm{log}|J_N(K)|}{N}=K$$
for any knot $`KS^3`$. Here $`v_3`$ is the volume of the ideal regular hyperbolic tetrahedron and $`K`$ is the Gromov norm of $`S^3K`$ .
###### Remark 1.2.
Kashaev’s conjecture is for hyperbolic knots;
###### Conjecture 1.3 (Kashaev ).
Let $`K`$ be a hyperbolic knot then
$$2\pi \underset{N\mathrm{}}{lim}\frac{\mathrm{log}|J_N(K)|}{N}=\mathrm{Vol}(S^3K).$$
Note that since for a hyperbolic knot $`K=\mathrm{Vol}(S^3K)/v_3`$, our conjecture is a generalization of Kashaev’s conjecture.
In this article I will describe how to calculate the colored Jones function and show various ways to confirm Kashaev’s conjecture for the figure-eight knot.
###### Remark 1.4.
1. Kashaev proved his conjecture for the knots $`4_1`$, $`5_2`$, and $`6_1`$ .
2. Kashaev and O. Tirkkonen proved the volume conjecture for torus knots .
3. Y. Yokota proved Kashaev’s conjecture for the knot $`6_2`$ in and suggests a proof for general hyperbolic knots. See also his forthcoming paper .
4. The volume conjecture is proved for the knots $`6_3`$, $`7_2`$ and $`8_9`$, and for the Whitehead link in .
###### Acknowledgments.
I would like to thank the participants in the workshops held at the International Institute for Advanced Study in October 1999 (some information is available on the WWW at http://www.iias.or.jp/research/suuken/1999.10.12.html) and at Kyushu University in December 1999. Especially I am grateful to J. Murakami, M. Okamoto, T. Takata, and Y. Yokota for useful conversations.
## 2. How to calculate the colored Jones function from a $`(1,1)`$-tangle
In this section I show how to calculate the colored Jones function of a knot from its $`(1,1)`$-tangle diagram.
Let $`D`$ be a $`(1,1)`$-tangle diagram of a given knot $`K`$. We assume that $`D`$ is generic with respect to the height function with end points at the top and at the bottom. In particular each crossing in $`D`$ is one of the following eight forms. The diagram is decomposed into arcs after we cut it into four arcs at each crossing. We label each arc and each crossing with parameter from $`0`$ to $`N1`$.
For each crossing we assign a complex number as follows, where $`\delta _{i,j}`$ is Kronecker’s delta.
$`:(_m^+)_{k,l}^{i,j}:=\delta _{l,i+m}\delta _{k,jm}{\displaystyle \frac{(q)_i^1(q)_j}{(q)_m(q)_k(q)_l^1}}(1)^{i+k+1}q^{ik+(i+k)/2+(N^2+1)/4}`$
$`:(_m^+)_{k,l}^{i,j}:=\delta _{i,l+m}\delta _{j,km}{\displaystyle \frac{(q)_i(q)_j^1}{(q)_m(q)_k^1(q)_l}}(1)^{j+l+1}q^{jl+(j+l)/2+(N^2+1)/4}`$
$`:(_m^{})_{k,l}^{i,j}:=\delta _{l,im}\delta _{k,j+m}{\displaystyle \frac{(\overline{q})_i(\overline{q})_j^1}{(\overline{q})_m(\overline{q})_k^1(\overline{q})_l}}(1)^{j+l+1}q^{jl(j+l)/2(N^2+1)/4}`$
$`:(_m^{})_{k,l}^{i,j}:=\delta _{i,lm}\delta _{j,k+m}{\displaystyle \frac{(\overline{q})_i^1(\overline{q})_j}{(\overline{q})_m(\overline{q})_k(\overline{q})_l^1}}(1)^{i+k+1}q^{ik(i+k)/2(N^2+1)/4}`$
$`:(_m^{})_{k,l}^{i,j}:=\delta _{j,km}\delta _{l,i+m}{\displaystyle \frac{(\overline{q})_i^1(\overline{q})_j^1}{(\overline{q})_m(\overline{q})_k^1(\overline{q})_l^1}}(1)^{i+j+1}q^{ij(k+l)/2(N^2+1)/4}`$
$`:(_m^{})_{k,l}^{i,j}:=\delta _{k,jm}\delta _{i,l+m}{\displaystyle \frac{(\overline{q})_i(\overline{q})_j}{(\overline{q})_m(\overline{q})_k(\overline{q})_l}}(1)^{k+l+1}q^{kl(i+j)/2(N^2+1)/4}`$
$`:(_m^+)_{k,l}^{i,j}:=\delta _{j,k+m}\delta _{l,im}{\displaystyle \frac{(q)_i(q)_j}{(q)_m(q)_k(q)_l}}(1)^{k+l+1}q^{kl+(i+j)/2+(N^2+1)/4}`$
$`:(_m^+)_{k,l}^{i,j}:=\delta _{k,j+m}\delta _{i,lm}{\displaystyle \frac{(q)_i^1(q)_j^1}{(q)_m(q)_k^1(q)_l^1}}(1)^{i+j+1}q^{ij+(k+l)/2+(N^2+1)/4}`$
For each local minimum and maximum where an arc labeled with $`i`$ goes from left to right, we assign the following quantities.
$`:_i:=q^{i(N1)/2}`$
$`:_i:=q^{i+(N1)/2}`$
Then we take the product of all the quantities above and take the summation with all the labels running non-negative integers less than $`N`$ keeping the labels of two end points of the $`(1,1)`$-tangle $`0`$. (We may choose the labels of the end points arbitrarily.)
Note that $`(_m^\pm )_{k,l}^{i,j}`$, $`_i`$, and $`_i`$ are obtained from the $`R`$-matrix $`R_J^\pm `$, $`\mu _J`$, and $`\mu _J^1`$ respectively which appear in the enhanced Yang-Baxter operator corresponding to the $`N`$-dimensional irreducible representation of $`sl(2,)`$ (see also ). We also note that $``$, $``$, and $``$ can be obtained from $``$, $``$, and $``$.
## 3. The colored Jones function of the figure-eight knot
Let us consider the figure-eight knot whose $`(1,1)`$-tangle description is shown in the following figure.
Now we attach a label to each arc and crossing noting that the difference of the labels should be non-negative if we go through an under-crossing and non-positive if we go through an over-crossing (note Kronecker’s delta in $`(_m^\pm )_{k,l}^{i,j}`$). The labels are indicated in the following figure, where integers of bold faces are attached to the crossings.
Then our invariant is
(3.1)
$$\begin{array}{cc}& J_N(4_1)\hfill \\ & =\underset{\begin{array}{c}0iN1\\ 0jN1\\ 0i+jN1\end{array}}{}(_0^+)_{0,i}^{i,0}\times (_j^+)_{j,i+j}^{i,\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}}\times (_0^{})_{0,j}^{j,0}\times _j\times (_i^{})_{\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0},j}^{i+j,i}\times _i\hfill \\ & =\underset{\begin{array}{c}0iN1\\ 0jN1\\ 0i+jN1\end{array}}{}\frac{(q)_i^1}{(q)_i^1}(1)^{i+1}q^{i/2+(N^2+1)/4}\frac{(q)_i^1}{(q)_j(q)_j^1(q)_{i+j}^1}(1)^{i+1}q^{i/2+(N^2+1)/4}\hfill \\ & \times \frac{(\overline{q})_j}{(\overline{q})_j}(1)^{j+1}q^{j/2(N^2+1)/4}\frac{(\overline{q})_{i+j}(\overline{q})_i}{(\overline{q})_i(\overline{q})_j}(1)^{j+1}q^{j/2(N^2+1)/4}\hfill \\ & \times q^{j(N1)/2}q^{i+(N1)/2}\hfill \\ & =\underset{\begin{array}{c}0iN1\\ 0jN1\\ 0i+jN1\end{array}}{}\frac{(q)_{i+j}(\overline{q})_{i+j}}{(q)_i(\overline{q})_j}.\hfill \end{array}$$
## 4. The limit of the colored Jones function of the figure-eight knot
In this section I describe the calculation of the limit of $`\mathrm{log}\left(J_N(4_1)\right)/N`$ due to T. Ekholm, and confirm Kashaev’s conjecture for the figure-eight knot. Putting $`k:=i+j`$ in (3.1) we have
$$\begin{array}{cc}& J_N(4_1)\hfill \\ & =\underset{k=0}{\overset{N1}{}}\underset{i=0}{\overset{k}{}}\frac{(q)_k(\overline{q})_k}{(q)_i(\overline{q})_{ki}}\hfill \\ & =\underset{k=0}{\overset{N1}{}}(1)^k(q^{1/2}q^{1/2})^{2k}[k]^2\hfill \\ & \times \underset{i=0}{\overset{k}{}}(1)^iq^{(k^2+k2ik2i)/4}\frac{1}{(q^{1/2}q^{1/2})^k[i]![ki]!}\hfill \\ & =\underset{k=0}{\overset{N1}{}}(1)^kq^{k(k+1)/4}(q^{1/2}q^{1/2})^k[k]\underset{i=0}{\overset{k}{}}(1)^iq^{i(k+1)/2}\left[\genfrac{}{}{0.0pt}{}{k}{i}\right]\hfill \end{array}$$
since $`(q)_l=(1)^lq^{l(l+1)/4}(q^{1/2}q^{1/2})^l[l]`$ and $`(\overline{q})_l=q^{l(l+1)/4}(q^{1/2}q^{1/2})^l[l]`$. Now using \[8, Lemma A.1\] we have
(4.1)
$$\begin{array}{cc}\hfill J_N(4_1)& =\underset{k=0}{\overset{N1}{}}(1)^kq^{k(k+1)/4}(q^{1/2}q^{1/2})^k[k](1q^1)(1q^2)\mathrm{}(1q^k)\hfill \\ & =\underset{k=0}{\overset{N1}{}}(1)^k\left\{(q^{1/2}q^{1/2})(q^{2/2}q^{2/2})\mathrm{}(q^{k/2}q^{k/2})\right\}^2\hfill \\ & =\underset{k=0}{\overset{N1}{}}\left\{2\mathrm{sin}\left(\frac{1}{N}\pi \right)\times 2\mathrm{sin}\left(\frac{2}{N}\pi \right)\times \mathrm{}\times 2\mathrm{sin}\left(\frac{k}{N}\pi \right)\right\}^2.\hfill \end{array}$$
###### Remark 4.1.
The formula (4.1) was obtained by Kashaev \[3, (2.2)\]. The colored Jones polynomial for generic $`t`$ is
$$\underset{k=0}{\overset{N1}{}}\underset{l=1}{\overset{k}{}}\left(t^{(N+l)/2}t^{(N+l)/2}\right)\left(t^{(Nl)/2}t^{(Nl)/2}\right),$$
which was first obtained by T. Le. One can obtain this formula by using technique described here using $`R`$-matrix for generic $`t`$.
Now we will calculate the $`N\mathrm{}`$ limit of $`\mathrm{log}(J_N(4_1))/N`$, which was first obtained by Kashaev . The following calculation is due to Ekholm.
###### Theorem 4.2 (Kashaev and Ekholm).
$$2\pi \underset{N\mathrm{}}{lim}\frac{\mathrm{log}(J_N(4_1))}{N}=6\mathrm{L}\left(\frac{\pi }{3}\right),$$
where $`\mathrm{L}(\alpha )`$ is the Lobachevsky function
$$\mathrm{L}(\alpha ):=_0^\alpha \mathrm{log}|2\mathrm{sin}\theta |d\theta .$$
Note that the volume of the ideal tetrahedron with face angles $`\alpha `$, $`\beta `$, and $`\gamma `$ $`(\alpha +\beta +\gamma =2\pi )`$ is $`\mathrm{L}(\alpha )+\mathrm{L}(\beta )+\mathrm{L}(\gamma )`$ and that the figure-eight knot complement can be decomposed into two regular ideal tetrahedra. Therefore the equality above shows that the left hand side equals the volume of the figure-eight knot, confirming Kashaev’s conjecture in this case.
###### Proof.
Put $`g_k:={\displaystyle \underset{j=1}{\overset{k}{}}}2\mathrm{sin}\left({\displaystyle \frac{j}{N}}\pi \right)`$ so that $`J_N(4_1)={\displaystyle \underset{k=0}{\overset{N1}{}}}g_k^2`$. Since
$$\{\begin{array}{cc}2\mathrm{sin}\left(\frac{j}{N}\pi \right)<1\text{when }\frac{j}{N}<\frac{1}{6}\text{ or }\frac{j}{N}>\frac{5}{6},\hfill & \\ 2\mathrm{sin}\left(\frac{j}{N}\pi \right)>1\text{when }\frac{1}{6}<\frac{j}{N}<\frac{5}{6},\hfill & \end{array}$$
$`g_k`$ is decreasing when $`k<{\displaystyle \frac{N}{6}}`$ or $`k>{\displaystyle \frac{5N}{6}}`$ and increasing when $`{\displaystyle \frac{N}{6}}<k<{\displaystyle \frac{5N}{6}}`$. Thus (roughly speaking) $`g_k`$ attains its maximum at $`k={\displaystyle \frac{5N}{6}}`$.
Since there are $`N`$ terms in the summation formula of $`J_N(4_1)`$, we have
$$g_{5N/6}^2J_N(4_1)Ng_{5N/6}^2.$$
Taking $`\mathrm{log}`$ and dividing by $`N`$ we have
$$\frac{2\mathrm{log}g_{5N/6}}{N}\frac{\mathrm{log}J_N(4_1)}{N}\frac{2\mathrm{log}g_{5N/6}}{N}+\frac{\mathrm{log}N}{N},$$
which turns out to be
$$2\underset{j=1}{\overset{5N/6}{}}\frac{\mathrm{log}2\mathrm{sin}\left(\frac{j}{N}\pi \right)}{N}\frac{\mathrm{log}J_N(4_1)}{N}2\underset{j=1}{\overset{5N/6}{}}\frac{\mathrm{log}2\mathrm{sin}\left(\frac{j}{N}\pi \right)}{N}+\frac{\mathrm{log}N}{N}.$$
Since $`\underset{N\mathrm{}}{lim}{\displaystyle \frac{\mathrm{log}N}{N}}=0`$, we have
$`\underset{N\mathrm{}}{lim}{\displaystyle \frac{\mathrm{log}J_N(4_1)}{N}}`$ $`=2\underset{N\mathrm{}}{lim}{\displaystyle \underset{j=1}{\overset{5N/6}{}}}{\displaystyle \frac{\mathrm{log}2\mathrm{sin}\left(\frac{j}{N}\pi \right)}{N}}`$
$`=2{\displaystyle _0^{5\pi /6}}{\displaystyle \frac{1}{\pi }}\mathrm{log}2\mathrm{sin}xdx`$
$`={\displaystyle \frac{2}{\pi }}\mathrm{L}\left({\displaystyle \frac{5\pi }{6}}\right).`$
On the other hand from \[7, Lemma 1\] we have
$`\mathrm{L}\left({\displaystyle \frac{\pi }{3}}\right)`$ $`=\mathrm{L}\left(2{\displaystyle \frac{\pi }{6}}\right)`$
$`=2\mathrm{L}\left({\displaystyle \frac{\pi }{6}}\right)+2\mathrm{L}\left({\displaystyle \frac{\pi }{6}}+{\displaystyle \frac{\pi }{2}}\right)`$
(since $`\mathrm{L}`$ has period $`\pi `$)
$`=2\mathrm{L}\left({\displaystyle \frac{\pi }{6}}\right)+2\mathrm{L}\left({\displaystyle \frac{\pi }{6}}{\displaystyle \frac{\pi }{2}}\right)`$
(since $`\mathrm{L}`$ is an odd function)
$`=2\mathrm{L}\left({\displaystyle \frac{\pi }{6}}\right)2\mathrm{L}\left({\displaystyle \frac{\pi }{3}}\right).`$
Therefore $`\mathrm{L}\left({\displaystyle \frac{\pi }{6}}\right)={\displaystyle \frac{3}{2}}\mathrm{L}\left({\displaystyle \frac{\pi }{3}}\right)`$ and
$$\mathrm{L}\left(\frac{5\pi }{6}\right)=\mathrm{L}\left(\frac{\pi }{6}\right)=\frac{3}{2}\mathrm{L}\left(\frac{\pi }{3}\right).$$
We finally have
$$\underset{N\mathrm{}}{lim}\frac{\mathrm{log}J_N(4_1)}{N}=\frac{1}{2\pi }6\mathrm{L}\left(\frac{\pi }{3}\right),$$
completing the proof. ∎
## 5. Saddle point method
In this section I follow Kashaev and calculate the limit of $`\mathrm{log}\left(J_N(4_1)\right)/N`$ directly from (3.1) using the saddle point method. Note that Kashaev calculated the limit from (4.1) using the same method.
From , $`(q)_i=S_\gamma (\gamma \pi )/S_\gamma (\pi +\gamma +2i\gamma )`$ and $`(\overline{q})_i=S_\gamma (\pi \gamma 2i\gamma )/S_\gamma (\pi \gamma )`$ with $`\gamma =\pi /N`$ and $`S_\gamma (p)`$ an analytic function of $`p`$ which behaves like
$$S_\gamma (p)\mathrm{exp}\left[\frac{\mathrm{Li}_2(e^{p\sqrt{1}})}{2\gamma \sqrt{1}}\right]$$
for small $`\gamma `$, where $`\mathrm{Li}_2(\zeta )`$ is Euler’s dilogarithm
$$\mathrm{Li}_2(\zeta )=_0^\zeta \frac{\mathrm{log}(1\xi )}{\xi }𝑑\xi .$$
Putting $`z:=q^i`$ we have
$`(q)_i`$ $`={\displaystyle \frac{S_\gamma (\gamma \pi )}{S_\gamma (\pi +\gamma +2i\gamma )}}`$
$`\mathrm{exp}\left[{\displaystyle \frac{\mathrm{Li}_2\left(e^{\gamma \sqrt{1}}\right)\mathrm{Li}_2\left(ze^{\gamma \sqrt{1}}\right)}{2\gamma \sqrt{1}}}\right]\mathrm{exp}\left[{\displaystyle \frac{\mathrm{Li}_2(z)}{2\gamma \sqrt{1}}}\right]`$
for small $`\gamma `$. Similarly putting $`w:=q^j`$ we have
$`(\overline{q})_j`$ $`\mathrm{exp}\left[{\displaystyle \frac{\mathrm{Li}_2(z^1)}{2\gamma \sqrt{1}}}\right],`$
$`(q)_{i+j}`$ $`\mathrm{exp}\left[{\displaystyle \frac{\mathrm{Li}_2(zw)}{2\gamma \sqrt{1}}}\right],`$
$`(\overline{q})_{i+j}`$ $`\mathrm{exp}\left[{\displaystyle \frac{\mathrm{Li}_2(z^1w^1)}{2\gamma \sqrt{1}}}\right].`$
Therefore for large $`N`$ we can replace the sum (3.1) with the following double contour integral
(5.1)
$$\mathrm{exp}\left[\frac{\mathrm{Li}_2(zw)+\mathrm{Li}_2(z^1w^1)+\mathrm{Li}_2(z)\mathrm{Li}_2(w^1)}{2\gamma \sqrt{1}}\right]𝑑z𝑑w$$
with suitably chosen contour. Then the saddle point method (or the method of steepest descent) tells us that (5.1) behaves as a function of $`\gamma `$ like
(5.2)
$$\mathrm{exp}\left[\frac{F(z_0,w_0)}{2\gamma \sqrt{1}}\right],$$
where $`F(z,w):=\mathrm{Li}_2(zw)+\mathrm{Li}_2(z^1w^1)+\mathrm{Li}_2(z)\mathrm{Li}_2(w^1)`$ and $`(z_0,w_0)`$ is a solution to the following system of partial differential equations:
(5.3)
$$\{\begin{array}{cc}\frac{F}{z}=0,\hfill & \\ \frac{F}{w}=0.\hfill & \end{array}$$
Since $`d\mathrm{Li}_2(z)/dz=\mathrm{log}(1z)/z`$, (5.3) turns out to be
$$\{\begin{array}{cc}\mathrm{log}(1zw)+\mathrm{log}(1z^1w^1)\mathrm{log}(1z)=0,\hfill & \\ \mathrm{log}(1zw)+\mathrm{log}(1z^1w^1)\mathrm{log}(1w^1)=0\hfill & \end{array}$$
and so we have
(5.4)
$$\{\begin{array}{cc}(1zw)(1z^1w^1)=1z,\hfill & \\ (1zw)(1z^1w^1)=1w^1.\hfill & \end{array}$$
Therefore we want to find a solution to the following system of equations:
(5.5)
$$\{\begin{array}{cc}z^2w^2zwz^2w+1=0,\hfill & \\ z^2w^2zwz+1=0,\hfill & \end{array}$$
giving the trivial solution $`z=w=1`$.
To get a non-trivial solution we will regard $`z`$ and $`w`$ as elements in $`\widehat{}=\{\mathrm{}\}`$ and put $`u:=zw`$. Then (5.5) becomes
$$\{\begin{array}{cc}u^2uzu+1=0,\hfill & \\ u^2uz+1=0.\hfill & \end{array}$$
Throwing away the trivial solution $`u=z=1`$, we have
(5.6)
$$\{\begin{array}{cc}z=0,\hfill & \\ u^2u+1=0,\hfill & \end{array}$$
giving non-trivial solutions $`(z=0,u=\mathrm{exp}(\pi \sqrt{1}/3))`$ and $`(z=0,\mathrm{exp}(5\pi \sqrt{1}/3))`$ ($`w=\mathrm{}`$). (I learned this ‘blow-up’ technique from J. Murakami.)
We denote by $`F_0`$ the $`F(z_0,w_0)`$ corresponding to the solution $`(z=0,w=\mathrm{},u=\mathrm{exp}(5\pi \sqrt{1}/3))`$. Then
$`\mathrm{}\left(F_0\right)`$ $`=\mathrm{}\left(\mathrm{Li}_2\left(\mathrm{exp}(5\pi \sqrt{1}/3)\right)\right)+\mathrm{}\left(\mathrm{Li}_2\left(\mathrm{exp}(5\pi \sqrt{1}/3)\right)\right)`$
$`+\mathrm{}\left(\mathrm{Li}_2(0)\right)\mathrm{}\left(\mathrm{Li}_2(0)\right)`$
$`=2\mathrm{}\left(\mathrm{Li}_2\left(\mathrm{exp}(\pi \sqrt{1}/3)\right)\right),`$
where $`\mathrm{}`$ denotes the imaginary part. Since $`\mathrm{}\left(\mathrm{Li}_2\left(\mathrm{exp}(\theta \sqrt{1})\right)\right)=\mathrm{L}(\theta )+2\mathrm{L}(\pi \theta /2)`$ from \[3, (3.7)\] (see also ), $`\mathrm{}\left(F_0\right)`$ equals the volume of the figure-eight knot complement.
From (5.2)
$$\underset{N\mathrm{}}{lim}\frac{\mathrm{log}|J_N(4_1)|}{N}=\mathrm{}\left(\frac{F(z_0,w_0)}{2\pi \sqrt{1}}\right)=\frac{\mathrm{}\left(F(z_0,w_0)\right)}{2\pi },$$
giving $`\mathrm{Vol}(S^34_1)/2\pi `$ again, where $`\mathrm{}`$ is the real part.
## 6. A cheating calculation
In this section I follow D. Thurston \[10, page 5\] to get the limit by a ‘formal’ calculation.
We put
$$f(i,j):=\frac{(q)_{i+j}(\overline{q})_{i+j}}{(q)_i(\overline{q})_j}$$
so that $`J_N(4_1)=_{i,j}f(i,j)`$. Now consider the ratios $`f(i,j)/f(i1,j)`$ and $`f(i,j)/f(i,j1)`$ to find ‘(local) maxima/minima’ of the function $`f`$ (I am cheating here!). To do that we will find a solution to the ‘partial difference equations’ $`f(i,j)/f(i1,j)=f(i,j)/f(i,j1)=1`$, which might give ‘(local) maxima/minima’.
Since
$$\{\begin{array}{cc}\frac{f(i,j)}{f(i1,j)}\hfill & =\frac{(1q^{i+j})(1q^{(i+j)})}{(1q^i)},\hfill \\ \frac{f(i,j)}{f(i,j1)}\hfill & =\frac{(1q^{i+j})(1q^{(i+j)})}{(1q^j)},\hfill \end{array}$$
we have
(6.1)
$$\{\begin{array}{cc}(1zw)(1z^1w^1)=1z,\hfill & \\ (1zw)(1z^1w^1)=1w^1\hfill & \end{array}$$
putting $`z:=q^i`$ and $`w:=q^j`$.
Now we have the same system of equations as in the previous section (5.4). We denote by $`f_{\text{MAX}}`$ <sup>1</sup><sup>1</sup>1MAX are Nana, Reina, Mina and Lina. Aquarius! the $`f_{i,j}`$ corresponding to the solution $`(z=0,w=\mathrm{},u=\mathrm{exp}(5\pi \sqrt{1}/3))`$ ($`u=zw`$). (I am also cheating since $`|q|=1`$ and so $`z=q^i`$ cannot be $`0`$!) Since $`(q)_i=(\overline{q})_j1`$ for large $`N`$ when $`q^i=z=0=w^1=q^j`$, we have
$`f_{\text{MAX}}`$ $`=\left|(q^{1/2}q^{1/2})(q^{2/2}q^{2/2})\times \mathrm{}\times (q^{(i+j)/2}q^{(i+j)/2})\right|^2`$
$`=\left\{2\mathrm{sin}\left({\displaystyle \frac{1}{N}}\right)\times 2\mathrm{sin}\left({\displaystyle \frac{2}{N}}\right)\times \mathrm{}\times 2\mathrm{sin}\left({\displaystyle \frac{5\pi }{6}}\right)\right\}^2`$
$`=g_{5N/6}^2,`$
which is the same value as before and gives the same limit.
Note that a similar calculation using (4.1) was indicated by Thurston and gives the same result as in § 4.
## 7. Geometry
As seen in §§ 5 and 6, if we use the triangulation described in (see the picture below), the only two tetrahedra corresponding to $`u^1=q^{(i+j)}=\mathrm{exp}(\pi \sqrt{1}/3)`$ survive after taking the limit.
$``$
Note that the equation $`u^2u+1=0`$ in (5.6) is the hyperbolicity equation for the figure-eight knot complement which determines its hyperbolic structure.
The calculations here suggest that in the limit each $`R`$-matrix corresponds to five ideal tetrahedra, some of which may collapse, and the partial differential equations (appeared in § 5) and the partial difference equations (appeared in § 6) give the same algebraic equations (5.6), which coincide with the hyperbolicity equations. Due to Yokota this holds in much more general situation . Therefore it is now very natural to expect that Kashaev’s conjecture is true for any hyperbolic knot and that there should be a rich theory behind it.
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# Symmetric products, permutation orbifolds and discrete torsion
## Appendix
This appendix is devoted to the proof of Eqs.(19,20) for the discrete torsion coefficients. To this end, let’s introduce for $`x,yS_n`$ the quantity
$$\widehat{\epsilon }(x,y)=\{\begin{array}{cc}\frac{\vartheta (x,y)}{\vartheta (y,x)}& \mathrm{if}xy=yx\\ \frac{\vartheta (x,y)\vartheta (x,x^1)}{\vartheta (y,x^1)}(1)^{\frac{\left|x\right|(\left|x\right|+1)}{2}}& \mathrm{if}xy=yx^1\end{array}$$
where $`\vartheta `$ denotes the non-trivial 2-cocycle of $`S_n`$. The above quantity is well defined - this needs to be checked for $`x^2=1`$, because in that case $`xy=yx`$ is equivalent to $`xy=yx^1`$ , so the two expressions for $`\widehat{\epsilon }(x,y)`$ should coincide - because $`\vartheta (x,x)=(1)^{\frac{\left|x\right|(\left|x\right|+1)}{2}}`$ for any $`xS_n`$ of order $`2`$, and satisfies the following conditions as a consequence of the cocycle relation satisfied by $`\vartheta `$ :
1. $`\widehat{\epsilon }(x,x)=1`$
2. $`\widehat{\epsilon }(x,y)\widehat{\epsilon }(y,x)=1`$ for $`xy=yx`$
3. $`\widehat{\epsilon }(x^z,y^z)=\widehat{\epsilon }(x,y)`$
4. $`\widehat{\epsilon }(x,yz)=\widehat{\epsilon }(x,y)\widehat{\epsilon }(x^y,z)`$ for $`y,z\widehat{C}(x)=\left\{y\right|xy=yx^{\pm 1}\}`$
Note that $`\widehat{C}(x)`$ is a subgroup of $`S_n`$. Our claim is that $`\widehat{\epsilon }(x,y)`$ equals the quantity
$$\epsilon (x,y):=(1)^{\left(\right|x|1)\left(\right|y|1)+|x,y|1}$$
for any $`y\widehat{C}(x)`$. This follows if we can show that $`\epsilon (x,y)`$ satisfies the above conditions for $`\widehat{\epsilon }`$ and is non-trivial, because $`H^2(S_n)=_2`$ implies that there could be at most one such quantity.
First, $`|x,x|=\left|x\right|`$implies that
$$\epsilon (x,x)=(1)^{(\left|x\right|1)\left|x\right|}=1$$
because the exponent is always even. On the other hand, the obvious relation $`|y,x|=|x,y|`$ implies that
$$\epsilon (y,x)=\epsilon (x,y)^1$$
(29)
because the values of $`\epsilon `$ are $`\pm 1`$. That $`\epsilon (x^z,y^z)=\epsilon (x,y)`$ holds follows from its definition. It remains to show that
$$\epsilon (x,yz)=\epsilon (x,y)\epsilon (x^y,z)$$
(30)
for $`y,z\widehat{C}(x)`$. First, let’s note that $`\sigma (x)=(1)^{\left|x\right|}`$ is the sign of the permutation $`x`$, and it is well known that
$$\sigma (xy)=\sigma (x)\sigma (y)$$
(31)
for any $`x,yS_n`$. But we may rewrite $`\epsilon `$ in the form
$$\epsilon (x,y)=\sigma (y)^{\left|x\right|1}(1)^{\left|x\right||x,y|}$$
It is straightforward to show that the cyclic subgroup $`x`$ generated by $`x`$ is normal in $`\widehat{C}(x)`$, consequently for any orbit $`\zeta `$ of $`\widehat{C}(x)`$, the orbits of $`x`$ on $`\zeta `$ form a block system. In other words, for any orbit $`\zeta 𝒪\left(\widehat{C}(x)\right)`$, all $`x`$ orbits $`\xi 𝒪_\zeta (x)`$ contained in $`\zeta `$ have the same length, and there is a homomorphism
$$\pi _\zeta :\widehat{C}(x)Sym\left(𝒪_\zeta (x)\right)$$
But for $`y\widehat{C}(x)`$ we have
$$\left|x\right||x,y|=\underset{\zeta 𝒪(\widehat{C}(x))}{}\left|\pi _\zeta (y)\right|$$
(32)
which implies
$$\epsilon (x,y)=\sigma (y)^{\left|x\right|1}\underset{\zeta 𝒪(\widehat{C}(x))}{}\sigma \left(\pi _\zeta (y)\right)$$
The homomorphism property of $`\pi _\zeta `$ and Eq.(31) then imply Eq.(30).
To complete the proof of $`\widehat{\epsilon }(x,y)=\epsilon (x,y)`$ we have to show that $`\epsilon `$ is non-trivial for $`n4`$, but this is obvious since $`\epsilon (x,y)=1`$ for any two commuting transpositions $`x,y`$.
*Work supported by grant OTKA T32453.*
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# The decay 𝜌⁰→𝜋⁺+𝜋⁻+𝛾 and the coupling constant gρσγ
## Abstract
The experimental branching ratio for the radiative decay $`\rho ^0\pi ^++\pi ^{}+\gamma `$ is used to estimate the coupling constant $`g_{\rho \sigma \gamma }`$ for a set of values of $`\sigma `$-meson parameters M<sub>σ</sub> and $`\mathrm{\Gamma }_\sigma `$. Our results are quite different than the values of this constant used in the literature.
The radiative decay process $`\rho ^0\pi ^++\pi ^{}+\gamma `$ has been studied employing different approaches . There are two mechanisms that can contribute to this radiative decay: the first one is the internal bremsstrahlung where one of the charged pions from the decay $`\rho ^0\pi ^++\pi ^{}`$ emits a photon, and the second one is the structural radiation which is caused by the internal transformation of the $`\rho `$-meson quark structure. Since the bremsstrahlung is well described by quantum electrodynamics, different methods have been used to estimate the contribution of the structural radiation.
Singer calculated the amplitude for this decay by considering only the bremsstrahlung mechanism since the decay $`\rho ^0\pi ^++\pi ^{}`$ is the main decay mode of $`\rho ^0`$-meson. He also used the universality of the coupling of the $`\rho `$-meson to pions and nucleons to determine the coupling constant $`g_{\rho \pi \pi }`$ from the knowledge of the coupling constant $`g_{\rho NN}`$. Later, Renard studied this decay among other vector meson decays into $`2\pi +\gamma `$ final states in a gauge invariant way with current algebra, hard-pion and Ward-identities techniques. He, moreover, established the correspondence between these current algebra results and the structure of the amplitude calculated in the single particle approximation for the intermediate states. In corresponding Feynman diagrams the structural radiation proceeds through the intermediate states as $`\rho ^0S+\gamma `$ where the meson S subsequently decays into a $`\pi ^+\pi ^{}`$ pair. He concluded that the leading term is the pion bremsstrahlung and that the largest contribution to the structural radiation amplitude results from the scalar $`\sigma `$-meson intermediate state. He used the rough estimate $`g_{\rho \sigma \gamma }1`$ for the coupling constant $`g_{\rho \sigma \gamma }`$ which was obtained with the spin independence assumption in the quark model. The coupling constant $`g_{\rho \pi \pi }`$ was determined using the then available experimental decay rate of $`\rho `$-meson and also current algebra results as $`3.2g_{\rho \pi \pi }4.9`$. On the other hand, the coupling constant $`g_{\sigma \pi \pi }`$ was deduced from the assumed decay rate $`\mathrm{\Gamma }100`$ MeV for the $`\sigma `$-meson as $`g_{\sigma \pi \pi }`$=3.4 with $`M_\sigma =400`$ MeV. Furthermore, he observed that the $`\sigma `$\- contribution modifies the shape of the photon spectrum for high momenta differently depending on the mass of the $`\sigma `$-meson. We like to note, however, that the nature of the $`\sigma `$-meson as a $`\overline{q}q`$ state in the naive quark model and therefore the estimation of the coupling constant $`g_{\rho \sigma \gamma }`$ in the quark model have been a subject of controversy. Indeed, Jaffe lately argued within the framework of lattice QCD calculation of pseudoscalar meson scattering amplitudes that the light scalar mesons are $`\overline{q}^2q^2`$ states rather than $`\overline{q}q`$ states.
Recently, on the other hand, the coupling constant $`g_{\rho \sigma \gamma }`$ has become an important input for the studies of $`\rho ^0`$-meson photoproduction on nucleons. The presently available data on the photoproduction of $`\rho ^0`$-meson on proton targets near threshold can be described at low momentum transfers by a simple one-meson exchange model . Friman and Soyeur showed that in this picture the $`\rho ^0`$-meson photoproduction cross section on protons is given mainly by $`\sigma `$-exchange. They calculated the $`\gamma \sigma \rho `$-vertex assuming Vector Dominance of the electromagnetic current, and their result when derived using an effective Lagrangian for the $`\gamma \sigma \rho `$-vertex gives the value $`g_{\rho \sigma \gamma }`$2.71 for this coupling constant. Later, Titov et al. in their study of the structure of the $`\varphi `$-meson photoproduction amplitude based on one-meson exchange and Pomeron-exchange mechanisms used the coupling constant $`g_{\varphi \sigma \gamma }`$ which they calculated from the above value of $`g_{\rho \sigma \gamma }`$ invoking unitary symmetry arguments as $`g_{\varphi \sigma \gamma }`$0.047. They concluded that the data at low energies near threshold can accommodate either the second Pomeron or the scalar mesons exchange, and the differences between these competing mechanisms have profound effects on the cross sections and the polarization observables.
It, therefore, appears of much interest to study the coupling constant $`g_{\rho \sigma \gamma }`$ that plays an important role in scalar meson exchange mechanism from a different perspective other than Vector Meson Dominance as well. For this purpose we calculate the branching ratio for the radiative decay $`\rho ^0\pi ^++\pi ^{}+\gamma `$, and using the experimental value $`0.0099\pm 0.0016`$ for this branching ratio , we estimate the coupling constant $`g_{\rho \sigma \gamma }`$. Our calculation is based on the Feynman diagrams shown in Fig. 1. The first two terms in this figure are not gauge invariant and they are supplemented by the direct term shown in Fig. 1(c) to establish gauge invariance. Guided by Renard’s current algebra results, we assume that the structural radiation amplitude is dominated by $`\sigma `$-meson intermediate state which is depicted in Fig. 1(d).
We describe the $`\rho \sigma \gamma `$-vertex by the effective Lagrangian
$`_{\rho \sigma \gamma }^{int.}={\displaystyle \frac{e}{M_\rho }}g_{\rho \sigma \gamma }[^\alpha \rho ^\beta _\alpha A_\beta ^\alpha \rho ^\beta _\beta A_\alpha ]\sigma `$ (1)
which also defines the coupling constant $`g_{\rho \sigma \gamma }`$. The $`\rho \pi \pi `$-vertex is described by the effective Lagrangian
$`_{\rho \pi \pi }^{int}=g_{\rho \pi \pi }\stackrel{}{\rho }_\mu (^\mu \stackrel{}{\pi }\times \stackrel{}{\pi })`$ (2)
using which we obtain the decay width $`\mathrm{\Gamma }(\rho \pi \pi )`$ as
$`\mathrm{\Gamma }(\rho \pi \pi )={\displaystyle \frac{g_{\rho \pi \pi }^2}{4\pi }}{\displaystyle \frac{M_\rho }{12}}\left[1({\displaystyle \frac{2M_\pi }{M_\rho }})^2\right]^{3/2}.`$ (3)
The experimental value of the width $`\mathrm{\Gamma }`$=151 MeV then yields the value $`\frac{g_{\rho \pi \pi }^2}{4\pi }`$=2.91 for the coupling constant $`g_{\rho \pi \pi }`$. For the $`\sigma \pi \pi `$-vertex we use the effective Lagrangian
$`_{\sigma \pi \pi }^{int}={\displaystyle \frac{1}{2}}g_{\sigma \pi \pi }M_\sigma \stackrel{}{\pi }\stackrel{}{\pi }\sigma .`$ (4)
The decay width of the $`\sigma `$-meson that follows from this effective Lagrangian is given as
$`\mathrm{\Gamma }_\sigma \mathrm{\Gamma }(\sigma \pi \pi )={\displaystyle \frac{g_{\sigma \pi \pi }^2}{4\pi }}{\displaystyle \frac{3M_\sigma }{8}}\left[1({\displaystyle \frac{2M_\pi }{M_\sigma }})^2\right]^{1/2}.`$ (5)
In our calculation of the invariant amplitude for the process $`\rho ^0\pi ^++\pi ^{}+\gamma `$, in the $`\sigma `$-meson propagator in Fig. 1(d) we make the replacement $`M_\sigma M_\sigma \frac{1}{2}i\mathrm{\Gamma }_\sigma `$, where $`\mathrm{\Gamma }_\sigma `$ is given by Eq. (5). Since the experimental candidate for $`\sigma `$-meson $`f_0`$(400-1200) has a width (600-1000) MeV , we obtain a set of values for the coupling constant $`g_{\rho \sigma \gamma }`$ by considering the ranges $`M_\sigma `$=400-1200 MeV, $`\mathrm{\Gamma }_\sigma `$=600-1000 MeV for the parameters of the $`\sigma `$-meson.
In terms of the invariant amplitude $``$(E<sub>γ</sub>, E<sub>1</sub>), the differential decay probability for an unpolarized $`\rho ^0`$-meson at rest is given by
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dE_\gamma dE_1}}={\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle \frac{1}{8M_\rho }}^2,`$ (6)
where E<sub>γ</sub> and E<sub>1</sub> are the photon and pion energies respectively, and we perform an average over the spin states of $`\rho ^0`$-meson and a sum over the polarization states of the photon. The decay width $`\mathrm{\Gamma }(\rho \pi ^++\pi ^{}+\gamma )`$ is the obtained by integration
$`\mathrm{\Gamma }={\displaystyle _{E_{\gamma ,min.}}^{E_{\gamma ,max.}}}𝑑E_\gamma {\displaystyle _{E_{1,min.}}^{E_{1,max.}}}𝑑E_1{\displaystyle \frac{d\mathrm{\Gamma }}{dE_\gamma dE_1}}.`$ (7)
The maximum photon energy E<sub>γ,max.</sub> is given as $`E_{\gamma ,max.}=(M_\rho ^24M_\pi ^2)/2M_\rho `$=334 MeV, and the minimum photon energy is taken as E<sub>γ,min.</sub>=50 MeV, since the experimental value of the branching ratio is determined for this range of photon energies . The maximum and minimum values for pion energy E<sub>1</sub> are given by
$`{\displaystyle \frac{1}{2(2E_\gamma M_\rho M_\rho ^2)}}[2E_\gamma ^2M_\rho +3E_\gamma M_\rho ^2M_\rho ^3`$ (8)
$`\pm E_\gamma \sqrt{(2E_\gamma M_\rho +M_\rho ^2)(2E_\gamma M_\rho +M_\rho ^24M_\pi ^2)}].`$ (9)
The photon spectra for the branching ration of the decay $`\rho ^0\pi ^++\pi ^{}+\gamma `$ are plotted in Fig. 2 as a function of photon energy E<sub>γ</sub>. The contributions of bremsstrahlung and structural radiation amplitude calculated with $`\sigma `$-meson intermediate state as well as the contribution of the interference term are shown together with the phase space for this decay as a function of photon energy. The parameters of the $`\sigma `$-meson adapted for the numerical calculations are M<sub>σ</sub>=500 MeV, $`\mathrm{\Gamma }_\sigma `$=800 MeV resulting in the coupling constants $`g_{\sigma \pi \pi }=8.04`$ and $`g_{\rho \sigma \gamma }=8.45`$. The contribution of the $`\sigma `$-term becomes increasingly important in the region of high photon energies dominating the contribution of the bremsstrahlung amplitude. Although its contribution is somewhat reduced by the interference term, the structural radiation makes the main contribution to the branching ratio in the region of high photon energies. The pion energy spectra for the branching ratio as a function of pion energy are shown in Fig. 3 for the same set of $`\sigma `$-meson parameters. The contribution of the $`\sigma `$-meson term is again significant in the region of low pion energies becoming quite insignificant for high pion energies compared to bremsstrahlung term and the interference term in general makes a small but negative contribution.
In Fig. 4 we show the dependence of the branching ratio calculated with the same parameters for the $`\sigma `$-meson on the minimum detected photon energy. This dependence is quite strong and therefore in our calculations we use minimum detected photon energy as 50 MeV as quoted in the experimental results . We furthermore plot the ratio R<sub>β</sub> as a function of $`\beta `$ in Fig. 5. This ratio is defined by
$`R_\beta ={\displaystyle \frac{\mathrm{\Gamma }_\beta }{\mathrm{\Gamma }_{tot}(\rho ^0\pi ^++\pi ^{}+\gamma )}}`$ (10)
where the numerator and denominator are given as
$`\mathrm{\Gamma }_\beta ={\displaystyle _{50}^{50+\beta }}𝑑E_\gamma {\displaystyle \frac{d\mathrm{\Gamma }}{dE_\gamma }},\mathrm{\Gamma }_{tot.}={\displaystyle _{50}^{E_{\gamma ,max.}}}𝑑E_\gamma {\displaystyle \frac{d\mathrm{\Gamma }}{dE_\gamma }}.`$ (11)
The shape of the curve for this ratio further reflects that the dependence of the branching ratio on the $`\sigma `$-term is pronounced in the region of the high photon energies.
We present the results of our calculation in Table 1. We note that by using the experimental value of the decay rate for the radiative decay $`\rho ^0\pi ^++\pi ^{}+\gamma `$ in the expression we obtain in our calculation for this decay rate, we arrive at a quadric equation for the coupling constant $`g_{\rho \sigma \gamma }`$, the coefficient of the quadric term resulting from the $`\sigma `$-meson contribution of Fig. 1(d) and the coefficient of the linear term from the interference of the $`\sigma `$-meson and the pion bremsstrahlung terms of Fig. 1(a,b,c). Therefore, our analysis produces a set of values for the coupling constant $`g_{\rho \sigma \gamma }`$ depending on the $`\sigma `$-meson parameters, and for a given set of $`\sigma `$-meson parameters results in two values for the coupling constant $`g_{\rho \sigma \gamma }`$, one being positive and the other one negative. Note that from the specific structure of the matrix element for the $`\sigma `$-meson contribution we can obtain in a straightforward manner the approximate relation
$`{\displaystyle \frac{g_{\rho \sigma \gamma }^2}{\mathrm{\Gamma }_\sigma M_\sigma ^3}}constant.`$ (12)
Furthermore, the values of the coupling constant $`g_{\rho \sigma \gamma }`$ resulting from our estimation are in general quite different than the values of this constant usually adopted for the one-meson exchange mechanism calculations existing in the literature. For example, Titov et al. uses the value $`g_{\rho \sigma \gamma }`$=2.71 which they obtain from Friman and Soyeur’s analysis of $`\rho `$-meson photoproduction using Vector Meson Dominance. It is interesting to note that in their study of pion dynamics in Quantum Hadrodynamics II, which is a renormalizable model constructed using local gauge invariance based on SU(2) group, that has the same Lagrangian densities for the vertices we use, Serot and Walecka come to the conclusion that in order to be consistent with the experimental result that s-wave $`\pi N`$-scattering length is anomalously small, in their tree-level calculation they have to choose $`g_{\sigma \pi \pi }`$=12. Since they use $`M_\sigma `$=520 MeV this implies $`\mathrm{\Gamma }_\sigma `$ 1700 MeV. If we use these values in our analysis, we then obtain $`g_{\rho \sigma \gamma }`$=11.91. Soyeur , on the other hand, uses quite arbitrarly the values $`M_\sigma `$=500 MeV, $`\mathrm{\Gamma }_\sigma `$=250 MeV, which in our calculation results in the coupling constant $`g_{\rho \sigma \gamma }`$=6.08. We like to note, however, that these values for $`\sigma `$-meson parameters are not consistent with the experimental data on $`\sigma `$-meson .
Our analysis and estimation of the coupling constant $`g_{\rho \sigma \gamma }`$ using the experimental value of the branching ratio of the radiative decay $`\rho ^0\pi ^++\pi ^{}+\gamma `$ give quite different values for this coupling constant than used in the literature. Furthermore, since we obtain this coupling constant as a function of $`\sigma `$-meson parameters, it will be of interest to study the dependence of the observables of the reactions, such as for example the photoproduction of vector mesons on nucleons $`\gamma +NN+V`$ where V is the neutral vector meson, analyzed using one-meson exchange mechanism on these parameters.
Acknowledgments
We thank Prof. Dr. M. P. Rekalo for suggesting this problem to us and for his guidance during the course of our work. We also wish to thank Prof. Dr. T. M. Aliev for helpful discussions.
Figure Captions:
Diagrams for the decay $`\rho ^0\pi ^++\pi ^{}+\gamma `$
The photon spectra for the decay width of $`\rho ^0\pi ^++\pi ^{}+\gamma `$. The contributions of different terms are indicated.
The pion energy spectra for the decay width of $`\rho ^0\pi ^++\pi ^{}+\gamma `$. The contributions of different terms are indicated.
The decay width of $`\rho ^0\pi ^++\pi ^{}+\gamma `$ as a function of minimum detected photon energy.
The ratio R$`{}_{\beta }{}^{}=\frac{\mathrm{\Gamma }_\beta }{\mathrm{\Gamma }_{tot}}`$ as a function of $`\beta `$.
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# One-loop effective potential
## One-loop effective potential
Here we will consider a simple non-supersymmetric tachyon-free $`Z_2`$ orientifold of type IIB superstring compactified to four dimensions on $`T^4/Z_2\times T^2`$ . Cancellation of Ramond-Ramond charges requires the presence of 32 D9 and 32 anti-D5 (D$`\overline{5}`$) branes <sup>4</sup><sup>4</sup>4In general arbitrary numbers of pairs D9+D$`\overline{9}`$ and D5+D$`\overline{5}`$ can also be added .. The bulk (closed strings) as well as the D9 branes are $`N=2`$ supersymmetric while supersymmetry is broken on the world-volume of the D$`\overline{5}`$’s. The massless closed string spectrum contains the graviton-, 19 vector- and 4 hyper-multiplets, while the massless open string spectrum on the D9 branes contains an $`N=2`$ vector multiplet in the adjoint of the $`SO(16)\times SO(16)`$ gauge group and a hypermultiplet in the (16,16) representation. When all D$`\overline{5}`$ branes are put at the origin of $`T^4`$, the non-supersymmetric D$`\overline{5}`$ sector contains gauge fields and complex scalars in the adjoint representation of $`USp(16)\times USp(16)`$ gauge group, a pair of complex scalars in the (16,16) representation, and Dirac fermions in the (120,1) \+ (1,120) \+ (16,16) representations. Finally there are 9$`\overline{5}`$ strings giving rise to complex scalars in the (16,1;1,16) \+ (1,16;16,1) together with Weyl fermions in the (16,1;16,1) \+ (1,16;1,16) representations, with respect to $`SO(16)\times SO(16)\times USp(16)\times USp(16)`$. Note that the $`9\overline{5}`$ spectrum is supersymmetric when D$`\overline{5}`$ gauge interactions are turned off.
We will restrict ourselves to the effective potential involving the scalars of the D$`\overline{5}`$ branes, namely the adjoints and bifundamentals of the $`USp(16)\times USp(16)`$ gauge group. The relevant part of the one-loop partition function corresponding to $`\overline{5}\overline{5}`$ open strings is
$`𝒜_{\overline{5}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(d_1+d_2)^2{\displaystyle \frac{V_8S_8}{\eta ^8}}W_4P_2+{\displaystyle \frac{1}{4}}(d_1d_2)^2{\displaystyle \frac{V_4O_4O_4V_4C_4C_4+S_4S_4}{\eta ^8}}\left({\displaystyle \frac{2\eta }{\theta _2}}\right)^2P_2`$
$`_{\overline{5}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(d_1+d_2)\left\{{\displaystyle \frac{\widehat{V}_8+\widehat{S}_8}{\widehat{\eta }^8}}W_4+{\displaystyle \frac{\widehat{V}_4\widehat{O}_4\widehat{O}_4\widehat{V}_4+\widehat{C}_4\widehat{C}_4\widehat{S}_4\widehat{S}_4}{\widehat{\eta }^8}}\left({\displaystyle \frac{2\widehat{\eta }}{\widehat{\theta }_2}}\right)^2\right\}P_2`$ (1)
where $`𝒜`$ and $``$ denote the contributions from the annulus and Möbius strip, respectively. In the above equation $`d_1=d_2=16`$, while $`V_{2n}`$, $`O_{2n}`$, $`C_{2n}`$ and $`S_{2n}`$ are the $`SO(2n)`$ characters,
$$V_{2n}=\frac{\theta _3^n\theta _4^n}{2\eta ^n},O_{2n}=\frac{\theta _3^n+\theta _4^n}{2\eta ^n},C_{2n}=\frac{\theta _2^ni^n\theta _1^n}{2\eta ^n},S_{2n}=\frac{\theta _2^n+i^n\theta _1^n}{2\eta ^n},$$
where $`\theta _i`$ are the Jacobi theta functions and $`\eta `$ the Dedekind eta function, depending on the usual complex variable $`\tau =it/2`$, with $`t`$ being the (real) annulus parameter. In the product of characters, the first factor stands for the contribution of space-time and $`T^2`$ world-sheet fermions, while the second factor represents the corresponding contribution from the internal $`T^4`$. The hatted functions are defined by $`\widehat{f}f(\tau +1/2)`$. Finally, $`P_2`$ ($`W_4`$) denotes the momentum (winding) lattice sum along the $`T_2`$ ($`T_4`$) torus ; for one dimension, they read:
$$P_1(\tau )=\underset{m}{}e^{2i\pi \tau m^2\alpha ^{}/R_{}^2};W_1(\tau )=\underset{n}{}e^{2i\pi \tau n^2R_{}^2/\alpha ^{}},$$
(2)
where $`\alpha ^{}M_s^2`$ is the Regge slope, and $`R_{}`$ ($`R_{}`$) denotes the radius of the corresponding dimension parallel (transverse) to the D-brane.
In both the annulus and Möbius amplitudes the first term stands for the untwisted contribution while the second term accounts for the $`Z_2`$ orbifold projection which differentiates $`T^4`$ and $`T^2`$ contributions. Its presence is due to the non-freely action of $`Z_2`$ at the origin of $`T^4`$ and thus it depends only on the lattice of $`T^2`$. It is obvious from Eq. (One-loop effective potential) that the $`Z_2`$ projection acts in a supersymmetric way, and therefore the second terms containing the twisted contribution vanish identically and will not play any role in our calculation.
In the first terms containing the untwisted contribution, $`V_8`$ and $`S_8`$ arise from bosons and fermions, respectively. Here, supersymmetry is explicitly broken via the orientifold projection realized by the Möbius amplitude. Indeed, from the change of sign of $`S_8`$ between $`𝒜`$ and $``$, it is manifest that the orientifold projection acts in opposite ways for bosons and fermions and breaks supersymmetry. More precisely, it symmetrizes the bosons and antisymmetrizes the fermions in each $`USp(16)`$ factor.
The tree-level scalar potential can be obtained by a truncation of an $`N=2`$ supersymmetric theory and has flat directions corresponding to the Wilson lines $`a`$ along the $`T^2`$ or $`T^4`$ directions. For longitudinal directions they amount to shifting the momenta $`mm+a`$ in Eq. (2), while for transverse directions they shift the windings $`nn+a`$ and describe brane separation. It follows that at one-loop level the flat directions are lifted since the Wilson lines acquire a potential from the Möbius amplitude which breaks supersymmetry. Without loss of generality we will consider a Wilson line $`a`$ along one direction of $`T^2`$ of radius $`R`$, and treat the other, upon T-duality, on the same footing as the dimensions of $`T^4`$ with a common radius $`r`$. After transforming the amplitudes (One-loop effective potential) in the transverse (closed string) channel and using the standard $`\theta `$-function Riemann identity, the one loop effective potential for the Wilson line is given by :
$`V_{\mathrm{eff}}(a)`$ $`=`$ $`{\displaystyle \frac{1}{32\pi ^4\alpha ^2}}{\displaystyle _0^{\mathrm{}}}𝑑l{\displaystyle \frac{\theta _2^4}{4\eta ^{12}}}\left(il+{\displaystyle \frac{1}{2}}\right){\displaystyle \frac{R}{r^5}}{\displaystyle \underset{\stackrel{}{m}}{}}e^{2\pi \frac{\stackrel{}{m}^2}{r^2}l}{\displaystyle \underset{n}{}}e^{4i\pi na}e^{2\pi n^2R^2l}`$
$`=`$ $`{\displaystyle \frac{1}{32\pi ^4\alpha ^2}}{\displaystyle _0^{\mathrm{}}}𝑑l{\displaystyle \frac{\theta _2^4}{4\eta ^{12}}}\left(il+{\displaystyle \frac{1}{2}}\right){\displaystyle \frac{R}{r^5}}{\displaystyle \underset{\stackrel{}{m}}{}}e^{2\pi \frac{\stackrel{}{m}^2}{r^2}l}\left(1+2{\displaystyle \underset{n>0}{}}\mathrm{cos}(4\pi na)e^{2\pi n^2R^2l}\right),`$
where the radii $`R`$ and $`r`$ are defined in units of $`\alpha ^{}`$.
In this setup, the canonically normalized scalar field $`h`$ associated to the Wilson line $`a`$ is $`h=a/gR`$, where $`g`$ is the gauge coupling, as can be easily seen by dimensional reduction. Let us first expand the effective potential in powers of $`h`$ and extract its quadratic (squared mass) term $`\mu ^2h^2/2`$. The result is:
$$\mu ^2=\frac{g^2}{2\pi ^2\alpha ^{}}_0^{\mathrm{}}𝑑l\frac{\theta _2^4}{4\eta ^{12}}\left(il+\frac{1}{2}\right)\frac{R^3}{r^5}\underset{\stackrel{}{m}}{}e^{2\pi \frac{\stackrel{}{m}^2}{r^2}l}\underset{n}{}n^2e^{2\pi n^2R^2l}.$$
(4)
It is easy to see that the integral converges. In fact, in the limit $`l\mathrm{}`$ the integrand falls off exponentially, while for $`l0`$ one can use the Poisson resummations
$`{\displaystyle \underset{m}{}}e^{2\pi \frac{m^2}{r^2}l}`$ $`=`$ $`{\displaystyle \frac{r}{\sqrt{2l}}}{\displaystyle \underset{p}{}}e^{\pi \frac{r^2}{2l}p^2},`$ (5)
$`{\displaystyle \underset{n}{}}n^2e^{2\pi n^2R^2l}`$ $`=`$ $`{\displaystyle \frac{1}{R\sqrt{2l}}}{\displaystyle \underset{n}{}}\left({\displaystyle \frac{1}{4\pi R^2l}}{\displaystyle \frac{n^2}{4R^4l^2}}\right)e^{\frac{\pi }{2R^2l}n^2},`$ (6)
and the identity
$$\frac{\theta _2^4}{\eta ^{12}}\left(il+\frac{1}{2}\right)=(2l)^4\frac{\theta _2^4}{\eta ^{12}}\left(\frac{i}{4l}+\frac{1}{2}\right),$$
to show that the integrand goes to a constant. Moreover, $`\mu ^2`$ is negative which implies that the origin is unstable and $`h`$ must acquire a non trivial VEV breaking the gauge symmetry. Note that the negative sign comes from the expansion of $`\mathrm{cos}4\pi na`$ in Eq. (One-loop effective potential) and is correlated with the positive sign of the contribution from the same states to the cosmological constant. Although this seems to be a general property in these models, we do not have a deeper understanding of the correlation between the sign of the mass term and the (massive) spectrum of the theory.
Even if $`a`$ is a periodic variable of period 1, $`V_{\mathrm{eff}}`$ is periodic under the shift $`aa+1/2`$, since its contribution originates from the Möbius amplitude. Moreover, in this particular example, the one-loop effective potential has a global minimum at $`a=1/4`$. This follows trivially from its expression (One-loop effective potential), whose derivative with respect to $`a`$ is a sum of terms proportional to $`\mathrm{sin}4\pi na`$, while its second derivative gives
$$V_{\mathrm{eff}}^{\prime \prime }_{a=1/4}=\frac{1}{2\pi ^2\alpha ^2}_0^{\mathrm{}}𝑑l\frac{\theta _2^4}{4\eta ^{12}}\left(il+\frac{1}{2}\right)\frac{R}{r^5}\underset{\stackrel{}{m}}{}e^{2\pi \frac{\stackrel{}{m}^2}{r^2}l}\underset{n}{}()^{n+1}n^2e^{2\pi n^2R^2l}.$$
(7)
Positivity of the integrand is manifest for all factors with the exception of the last sum for which a careful analysis is required. This sum can be written as $`_\tau \theta _4(\tau )/2i\pi `$ with $`\tau =2iR^2l`$, which can be easily shown to be a positive function.
In the $`T`$-dual picture, the VEV $`a=1/4`$ corresponds to separating a brane at a distance from the origin equal to half the compactification interval $`\pi R`$. By turning on all Wilson lines $`a_I`$, the effective potential becomes a sum $`_IV_{\mathrm{eff}}(a_I)`$, with $`V_{\mathrm{eff}}(a_I)`$ given in (One-loop effective potential), which upon minimization fixes all $`a_I`$ at the same value $`1/4`$. Thus, the global minimum of all Wilson lines corresponds to put all branes at the same point in the middle of the compactification interval. The $`USp(16)\times USp(16)`$ gauge group is then broken down to a $`U(8)\times U(8)`$ or $`USp(16)`$ subgroup, corresponding to turning on Wilson lines along the $`T^2`$ or $`T^4`$ directions, transforming in the adjoint or in the bifundamental representation, respectively.
In order to make a numerical estimate of the results, we will consider the case of a 4-brane with five large transverse dimensions by taking the limit $`r\mathrm{}`$ and keeping the radius $`R`$ (along the 4-brane) as a parameter. To take the limit $`r\mathrm{}`$, we use Eq. (5) for each of the five transverse dimensions, and note that only $`p=0`$ contributes in the sum. In fact, non vanishing values of $`p`$ may contribute only in the region $`l\mathrm{}`$, in which case the corresponding integrand in Eq. (One-loop effective potential) vanishes as $`l^{5/2}`$. It follows that in the limit $`r\mathrm{}`$ the potential becomes:
$$V_{\mathrm{eff}}(a,R)=\frac{R}{32\pi ^4\alpha ^2}_0^{\mathrm{}}\frac{dl}{\left(2l\right)^{5/2}}\frac{\theta _2^4}{4\eta ^{12}}\left(il+\frac{1}{2}\right)\underset{n}{}e^{4i\pi na}e^{2\pi n^2R^2l}.$$
(8)
The effective potential (8) is plotted in Fig. 1 for the range of values of the radius $`2R3`$ as a function of $`a`$ inside its period $`(1/2,1/2)`$. Following our previous analysis, it has a maximum at the origin and a minimum at $`a=\pm 1/4`$ for any value of $`R`$.
The mass term at the origin, in the limit $`r\mathrm{}`$ and for arbitrary $`R`$, can be equally computed from Eq. (4) using the Poisson resummation (5). The result is:
$$\mu ^2(R)=\epsilon ^2(R)g^2M_s^2$$
(9)
with
$$\epsilon ^2(R)=\frac{1}{2\pi ^2}_0^{\mathrm{}}\frac{dl}{\left(2l\right)^{5/2}}\frac{\theta _2^4}{4\eta ^{12}}\left(il+\frac{1}{2}\right)R^3\underset{n}{}n^2e^{2\pi n^2R^2l}.$$
(10)
The parameter $`\epsilon `$ is plotted in Fig. 2 as a function of $`R`$ in a typical range $`1/4<R<5`$.
At the lower end, it has almost reached its asymptotic value for $`R0`$ <sup>5</sup><sup>5</sup>5This limit corresponds, upon T-duality, to a large transverse dimension of radius $`1/R`$., $`\epsilon (0)0.14`$, and the effective cutoff for the mass term at the origin is $`M_s`$, as can be seen from Eq. (9). At large $`R`$, $`\mu ^2(R)`$ falls off as $`1/R^2`$, which is the effective cutoff in the limit $`R\mathrm{}`$, in agreement with field theory results in the presence of a compactified extra dimension <sup>6</sup><sup>6</sup>6Actually this effect is at the origin of thermal squared masses, $`T^2`$, in four-dimensional field theory at finite temperature, $`T`$, where the time coordinate is compactified on a circle of inverse radius $`1/RT`$ and the Boltzmann suppression factor generates an effective cutoff at momenta $`pT`$.. In fact, in the limit $`R\mathrm{}`$ an analytic approximation to $`\epsilon (R)`$ can be computed as,
$$\epsilon (R)\frac{\epsilon _{\mathrm{}}}{M_sR},\epsilon _{\mathrm{}}^2=\frac{3\zeta (5)}{4\pi ^4}0.008,$$
(11)
which approximately describes Fig. 2 for large values of $`R`$.
Notice that the mass term (9) we found for the Wilson line $`a`$ also applies, by gauge invariance, to the charged massless fields which belong to the same representation.
## Electroweak symmetry breaking
In the previous example we obtained a VEV of the order of the string scale, because we only considered Wilson lines, which correspond to tree-level flat directions in the Cartan subalgebra of the gauge group, and have put to zero the VEV’s of all other fields. Thus, the total potential to be minimized appeared at the one-loop level. Had we minimized the effective potential with respect to fields charged under the Cartan subalgebra, we would have found the same solution (which corresponds to a true minimum in the multidimensional field space) since the charged fields acquire, from the Wilson lines, positive tree-level squared masses and have vanishing VEV’s. In more realistic models, the Wilson lines are at least partially projected away by an orbifold projection which also breaks the gauge group. If the orbifold projection acts in a supersymmetric way, as was the case of the $`Z_2`$ in the previous example, the calculation of the squared mass term remains valid for the left-over charged scalars in the spectrum, up to an overall numerical factor given by the order of the orbifold group ($`1/N`$ for a $`Z_N`$ orbifold). Moreover, the charged scalars have a tree-level potential which can be obtained by an appropriate truncation, dictated by the orbifold, of a supersymmetric theory. These two facts allow the existence of a (local) perturbative minimum, around which higher order terms in the expansion of the one loop potential can be neglected since the charged scalars would acquire a VEV controlled by the quadratic terms.
We will illustrate these points within the context of the toy model described in the previous section. The crucial property is that the bosonic sector of the non-supersymmetric D$`\overline{5}`$ branes is identical to the one of an $`N=2`$ supersymmetric theory obtained by a $`Z_2`$ orbifold projection from an $`N=4`$ theory based on a “fictitious” $`USp(32)`$ gauge group. The latter contains six adjoint scalars that can be organized in three $`N=1`$ chiral multiplets $`\mathrm{\Phi }_i`$ with $`i=1,2,3`$. Notice that in this model supersymmetry is explicitly broken because the fermions belong to the antisymmetric instead of the adjoint (symmetric) representation of $`USp(32)`$. The $`Z_2`$ projection breaks $`USp(32)`$ into $`USp(16)\times USp(16)`$ and keeps the adjoint of $`USp(16)\times USp(16)`$ from $`\mathrm{\Phi }_1`$ and the $`(\mathrm{𝟏𝟔},\mathrm{𝟏𝟔})`$ components from $`\mathrm{\Phi }_{2,3}`$. The tree-level scalar potential can be obtained straightforwardly by a corresponding truncation of the potential of the $`N=4`$ theory:
$$V_{N=4}=\frac{g^2}{2}\mathrm{Tr}\left(\underset{i,j}{}\left|[\mathrm{\Phi }_i,\mathrm{\Phi }_j]\right|^2+\left(\underset{i}{}[\mathrm{\Phi }_i,\mathrm{\Phi }_i^{}]\right)^2\right).$$
(12)
The result is identical to the potential of an $`N=2`$ theory with $`USp(16)\times USp(16)`$ gauge group and one hypermultiplet in the $`(\mathrm{𝟏𝟔},\mathrm{𝟏𝟔})`$ representation. In $`N=1`$ notation, it corresponds to the superpotential $`W=g/\sqrt{2}\phi _2\phi _1\phi _3`$ where $`\phi _1`$ is the adjoint from $`\mathrm{\Phi }_1`$ and $`\phi _{2,3}`$ are the two bifundamental chiral multiplets from $`\mathrm{\Phi }_{2,3}`$. The $`F`$\- and $`D`$-term contributions to the potential come from the first and second term of Eq. (12), respectively.
As we discussed in detail after Eq. (2), the $`Z_2`$ orbifold projection does not by itself break all supersymmetries and does not play any role in the computation of the potential. As a result, the scalar mass terms generated at one loop receive contributions only from the untwisted sector which treats the adjoint and the $`(\mathrm{𝟏𝟔},\mathrm{𝟏𝟔})`$ scalars in the same way, as an adjoint of $`USp(32)`$. Thus, the generated masses of the different scalars can be obtained from the same functional of the radii through permutations. In particular, this means that scalars describing displacement of branes in dimensions of the same size acquire equal masses. For instance, in the isotropic 3-brane limit of six large transverse dimensions, $`r\mathrm{}`$ and $`R0`$, the result (9) applies for all scalar components.
We would like now to discuss possible phenomenological applications of these results. Let us assume that there is a sequence of “supersymmetric” orbifold projections that lead to the Standard Model living on some non-supersymmetric brane configuration along the line of the toy model presented above. In the minimal case, where there is only one Higgs doublet $`h`$ originating from the untwisted sector, the scalar potential would be:
$$V=\lambda (h^{}h)^2+\mu ^2(h^{}h),$$
(13)
where $`\lambda `$ arises at tree-level and is given by an appropriate truncation of a supersymmetric theory. Within the minimal spectrum of the Standard Model, $`\lambda =(g_2^2+g^2)/8`$, with $`g_2`$ and $`g^{}`$ the $`SU(2)`$ and $`U(1)_Y`$ gauge couplings, as in the MSSM. On the other hand, $`\mu ^2`$ is generated at one loop and can be estimated by Eqs. (9) and (10).
The potential (13) has a minimum at $`h=(0,v/\sqrt{2})`$, where $`v`$ is the VEV of the neutral component of the $`h`$ doublet, fixed by $`v^2=\mu ^2/\lambda `$. Using the relation of $`v`$ with the $`Z`$ gauge boson mass, $`M_Z^2=(g_2^2+g^2)v^2/4`$, and the fact that the quartic Higgs interaction is provided by the gauge couplings as in supersymmetric theories, one obtains for the Higgs mass a prediction which is the MSSM value for $`\mathrm{tan}\beta \mathrm{}`$ and $`m_A\mathrm{}`$:
$$M_h=M_Z.$$
(14)
Furthermore, one can compute $`M_h`$ in terms of the string scale $`M_s`$, as $`M_h^2=2\mu ^2=2\epsilon ^2g^2M_s^2`$, or equivalently
$$M_s=\frac{M_h}{\sqrt{2}g\epsilon }$$
(15)
The lowest order relations (14) and (15) receive in general two kinds of higher order corrections. On the one hand, there might be important string corrections that we will discuss in the next section. On the other hand, from the point of view of the effective field theory, they are valid at the string scale $`M_s`$, and Standard Model radiative corrections should be taken into account for scales between $`M_s`$ and $`M_Z`$. In particular, the tree level Higgs mass has been shown to receive important radiative corrections from the top-quark sector. For present experimental values of the top-quark mass, the Higgs mass in Eqs. (14) and (15) is raised to values around 120 GeV . Moreover from Eq. (15), we can compute the string scale $`M_s`$. There is a first ambiguity in the value of the gauge coupling $`g`$ at $`M_s`$, which depends on the details of the model. Here, we use a typical unification value $`g1/\sqrt{2}`$. A second ambiguity concerns the numerical coefficient $`\epsilon `$ which is in general model dependent. In our calculation, this is partly reflected in its $`R`$-dependence, as seen in Fig. 2. Varying $`R`$ from 0 to 5, that covers the whole range of values for a transverse dimension $`1<1/R<\mathrm{}`$, as well as a reasonable range for a longitudinal dimension 1<R
<
51𝑅
<
51<R\mathrel{\lower 2.5pt\vbox{\hbox{$<$}\hbox{$\sim$}}}5, one obtains $`M_s15`$ TeV. Note that in the $`R1`$ (large longitudinal dimension) region our theory is effectively cutoff by $`1/R`$ and the Higgs mass is then related to it by,
$$\frac{1}{R}=\frac{M_h}{\sqrt{2}g\epsilon _{\mathrm{}}}.$$
(16)
Using now the value for $`\epsilon _{\mathrm{}}`$ in the present model, Eq. (11), we find 1/R
>
1
>
1𝑅11/R\mathrel{\lower 2.5pt\vbox{\hbox{$>$}\hbox{$\sim$}}}1 TeV.
A further model dependence of $`\epsilon `$ comes from the order of the orbifold group. As mentioned above, had we considered a higher order orbifold, e.g. $`Z_{2N}`$ instead of $`Z_2`$ as required by more realistic models, $`\epsilon `$ would decrease by a factor $`\sqrt{N}`$. As a result, the radiative electroweak symmetry breaking can be consistent with a string scale as heavy as $`𝒪(10)`$ TeV and a compactification scale 1/R
>
2
>
1𝑅21/R\mathrel{\lower 2.5pt\vbox{\hbox{$>$}\hbox{$\sim$}}}2 TeV.
In a more general context, the Higgs sector may be more complicated and the scalar potential could have classically undetermined flat directions as discussed in the introduction. For concreteness we will consider the case of two Higgs doublets $`h_1`$ and $`h_2`$ with a tree-level potential, obtained by an appropriate truncation of a supersymmetric theory, and equal to that of the MSSM. We are also assuming two different one-loop generated squared mass terms $`\mu _1^2`$ and $`\mu _2^2`$ for the Higgs fields:
$$V=\lambda \left(\left|h_1\right|^2\left|h_2\right|^2\right)^2+\rho \left|h_1^{}h_2\right|^2+\mu _1^2\left|h_1\right|^2+\mu _2^2\left|h_2\right|^2$$
(17)
where $`\lambda =(g_2^2+g^2)/8`$ and $`\rho =g_2^2/2`$. The conditions for having a stable minimum are $`\mu _1^2+\mu _2^2>0`$ and $`\mu _1^2\mu _2^2<0`$. These conditions are fulfilled provided that one of the masses, say $`\mu _2^2`$, is negative and the other, say $`\mu _1^2`$, is positive. In this case we get the VEV’s $`h_1=0`$ and $`h_2=(0,v/\sqrt{2})`$, where $`v^2=\mu _2^2/\lambda `$. Using again the relation of $`v`$ with $`M_Z`$, we obtain the tree-level Higgs mass spectrum :
$$M_{h_2}=M_Z,M_{h_1^0}^2=\mu _1^2+\mu _2^2,M_{h_1^{}}^2=M_{h_1^0}^2+M_W^2,$$
(18)
where $`h_2`$ corresponds to the Standard Model Higgs, and $`h_1^0`$, $`h_1^{}`$ to the neutral and charged components of the $`h_1`$ doublet. Moreover, the string scale is given by
$$M_s=\frac{M_{h_2}}{\sqrt{2}g\epsilon _2}$$
(19)
with $`\mu _2^2=\epsilon _2^2g^2M_s^2`$.
Again, these are tree-level relations which are subject to both string and Standard Model radiative corrections. In particular, the latter provide important contributions to the mass of the Standard Model Higgs $`h_2`$, which is increased roughly to $``$ 120 GeV, and accordingly to the string scale given in Eq. (19). It is interesting that we obtained the same relations as in the previous example with a single Higgs field. The difference is that there is also a left-over scalar doublet whose neutral and charged components acquire masses given in Eq. (18). As we have pointed out, in this case one needs the one-loop generated squared masses for the two scalar doublets, $`\mu _1^2`$, $`\mu _2^2`$, to be different and opposite in sign. Although our toy string example allows for different values by introducing different radii, the change in sign requires more general models, such as those obtained for instance by introducing additional pairs of branes - anti-branes .
## Discussion on string threshold corrections
We discuss now string threshold corrections to the relations (14) and (15). These are moduli dependent and may become very important only when some radii become large compared to the string length. Otherwise, if all radii are of order one in string units, higher loop corrections are order one numbers multiplied by loop factors which are suppressed when string theory is weakly coupled. Of course, these (model dependent) corrections are needed for a detailed phenomenological analysis and could be as important as those of the MSSM that increase the Higgs mass by roughly 10%. An estimate of these corrections can be done by an explicit computation of the $`a^4`$ terms in the expansion of the potential (8). Notice though that these terms do not determine uniquely the one-loop corrections to the quartic couplings of the charged fields, partly because there are more than one gauge invariant combinations. An additional subtlety is the existence of an infrared divergence as $`l0`$, which is due to the low energy running of the couplings and must be appropriately subtracted to obtain the string threshold corrections in a definite renormalization scheme .
For dimensions longitudinal to our world brane, the large radius limit leads in general the theory very rapidly to a non perturbative regime, since the (ten-dimensional) string coupling becomes strong when four-dimensional gauge couplings are of order unity. On the other hand, for large transverse dimensions, the tree-level string coupling remains perturbative (of order of the gauge couplings), and therefore their size can in principle become as large as desired. If this is the case, the decompactification limit exists, and threshold corrections are again controlled by the string coupling and are suppressed by loop factors. However, this limit does not exist in general when there are massless bulk fields that propagate in one or two transverse dimensions, and threshold corrections become very important .
A way to see how these large corrections to the parameters of the effective lagrangian on the brane arise, is to look at the ultraviolet open string loop diagrams as emission of massless closed strings in the bulk at the location of distant sources created by other branes or orientifold planes. This emission leads to corrections that diverge linearly or logarithmically with the size of transverse space, if there are massless closed string states propagating in one or two dimensions, respectively. The case of one large transverse dimension is similar to that of a large longitudinal one, since threshold corrections grow linearly with the radius and bring rapidly the theory to a non-perturbative regime . In this case, one can fine-tune the radius to a narrow region near the string scale and the low energy parameters will be very sensitive to the initial conditions.
In the case of two large transverse dimensions, the logarithmic contributions to the parameters of the effective action on the brane are similar to those in a renormalizable theory and can be resummed as in the renormalization group improved MSSM . In this analogy, the string scale $`M_s`$ plays the role of the supersymmetry breaking scale, while the size of the transverse space replaces the ultraviolet cutoff at the Planck mass, $`M_P`$. For instance, if the bulk contains $`n`$ large transverse dimensions of common radius $`R_{}`$, while the remaining $`6n`$ have string size, one obtains the familiar relation $`M_P^2=M_s^{2+n}R_{}^n`$. When there are massless bulk fields propagating in two of them, like e.g. twisted moduli localized at an $`n2`$ dimensional subspace, the logarithmic corrections are $`\mathrm{log}(R_{}M_s)=(2/n)\mathrm{log}(M_P/M_s)`$.
Concerning the Higgs mass considered here, such large radius dependent contributions would arise if there are bulk massless fields emitted by the Higgs at zero external momentum. The vanishing of such tree-level couplings, as for instance with bulk gravitons, implies the absence of large threshold corrections for the Higgs mass at the one-loop level. This is in agreement with our result (4) which remains finite in the decompactification limit for any number of large transverse dimensions. However, large corrections can arise at higher orders, e.g. through gravitons emitted from open string loops. While computation of such effects is out of the scope of this work, we would like to discuss the general structure of such corrections and comment on their phenomenological implications.
In the simplest case, the relevant part of the world brane action in the string frame is:
$$_{\mathrm{brane}}=e^\varphi \left\{\omega ^2|DH|^2+\frac{1+\mathrm{tan}^2\theta _W}{8}\omega ^4(H^{}H)^2+\frac{1}{4}(F_{SU(2)}^2+\mathrm{cot}^2\theta _WF_Y^2)\right\}\epsilon ^2M_s^2\omega ^4|H|^2,$$
(20)
where $`\varphi `$ is the string dilaton, $`\omega `$ the scale factor of the four-dimensional (world brane) metric, $`H`$ the Higgs scalar (in the string frame) and $`D`$ the gauge covariant derivative. The weak angle at the string scale $`\theta _W`$ must be correctly determined in the string model. Notice that the last term has no $`e^\varphi `$ dependence since it corresponds to a one loop correction. The bulk fields $`\varphi `$ and $`\omega `$ are evaluated in the transverse coordinates at the position of the brane. The physical couplings $`g_2`$, $`\lambda `$ and the mass $`\mu ^2`$ are given by
$$g_2=e^{\varphi /2},\lambda =\frac{1+\mathrm{tan}^2\theta _W}{8}e^\varphi ,\mu ^2=\epsilon ^2e^\varphi \omega ^2M_s^2,$$
(21)
while Eq. (14) remains unchanged and the relation (15) becomes
$$M_s=\frac{M_h}{\sqrt{2}\epsilon e^{\varphi /2}\omega }.$$
(22)
The lowest order result (15) corresponds to the (bare) value $`\omega =1`$.
As we discussed above, when the bulk fields $`\varphi `$ and $`\omega `$ propagate in two large transverse dimensions, they acquire a logarithmic dependence on these coordinates due to distant sources. Since the value of $`\varphi `$ at the position of the world brane is fixed by the value of the gauge coupling in Eq. (21), the relation (14) for the Higgs mass is not affected, while Eq. (22) for the string scale is corrected by a renormalization of $`\omega `$ which takes the generic form:
$$\omega =1+b_\omega g_2^2\mathrm{ln}(R_{}M_s),$$
(23)
where $`b_\omega `$ is a numerical coefficient. This correction is similar to a usual renormalization factor in field theory, which here is due to an infrared running in the transverse space. Depending on the sign of $`b_\omega `$, it can enhance ($`b_\omega <0`$) or decrease ($`b_\omega >0`$) the value of the string scale by the factor $`1/\omega `$. This effect can be important since the involved logarithm is large, varying between 7 and 35, for $`R_{}`$ between 1 fm and 1 mm.
In more general models, there are additional bulk fields entering in the expression of low energy couplings on the brane, such as the twisted moduli localized at the orbifold fixed points. As a result, every term in the lagrangian (20) may be multiplied by a different combination of the bulk fields that acquires an independent correction, similarly to Eq. (23). Thus, in the generic case, both relations (14) and (15) may be modified by corresponding renormalization factors that are computable in every specific model. In particular, the prediction of $`120`$ GeV for the Higgs mass, which coincides with that of the lightest Higgs in the MSSM for large values of $`\mathrm{tan}\beta `$ and $`m_A`$, can change by this effect.
A final important question that we have not addressed in this letter is the possible signatures of Higgs production in brane world models. Previous works done in the context of the effective field theory suggest that there may be new effects, leading in general to signatures that are different from those in the Standard Model or the MSSM . It will be interesting to study this issue in the framework of the non-supersymmetric type I string models we discussed here.
## Acknowledgments
This work was partly supported by the EU under TMR contracts ERBFMRX-CT96-0090 and ERBFMRX-CT96-0045, by CICYT (Spain) under contract AEN98-0816 and by IN2P3-CICYT contract Pth 96-3. K.B. thanks the CPHT of Ecole Polytechnique for hospitality.
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# On the efficiency of internal shocks in gamma-ray bursts
## 1. Introduction
Cosmological gamma-ray bursts (GRBs) are huge explosions of energy $`10^{53}`$ ergs, which may be triggered, e.g., by coalescence of neutron stars and/or black holes. The created fireball expands relativistically, with a Lorentz factor $`\mathrm{\Gamma }10^2`$. In the most likely scenario, the observed gamma-rays are produced at times $`t10^210^5`$ s after the explosion, when the relativistic outflow gets optically thin and before it is decelerated by the surrounding medium (see, e.g., Piran 1999 for a review). The outflow has a fluctuating velocity profile, and the $`\gamma `$-rays are generated by internal shocks that develop when faster shells try to overtake slower ones (Rees & Mészáros 1994).
The internal dissipation was discussed and simulated numerically in a number of works (e.g., Kobayashi, Sari, & Piran 1997; Daigne & Mochkovitch 1998; Lazzati, Ghisellini, & Celotti 1999; Panaitescu, Spada, & Mészáros 1999; Spada, Panaitescu, Mészáros 1999; Kumar 1999). It was concluded that only a few percent of the total kinetic energy is emitted to infinity (e.g., Panaitescu et al. 1999; Kumar 1999). Then the energy of the prompt GRB should be dominated by the early afterglow associated with an external shock. Observations show the opposite. Besides, the low efficiency requires huge explosion energies.
The conclusion that the internal dissipation has a low efficiency was based on certain assumptions:
(1) A specific probability distribution was assumed for the outflow Lorentz factor, such that the root-mean-square (rms) of the fluctuations was less than $`1/\sqrt{3}`$ (see eq. ). In fact, the rms can be much higher.
(2) The fluctuations were taken with a white (flat) spectrum.
(3) It was often assumed that the dissipated energy is equally distributed between the electrons, protons, and magnetic field, and only the electron part is radiated. In fact, the radiative capability of shocked matter is highly uncertain and dependent on poorly understood plasma processes. The high radiative capability is favored by the observed high luminosities of the prompt GRBs. If most of the dissipated energy is transferred to the electrons and radiated, the outflow stays at low pressure and the coasting matter gets concentrated into thin shells (caustics).
In § 2, we give a simple analytical description for dissipation of white fluctuations with modest rms (the linear regime). Then the efficiency is proportional to the mean square of the initial fluctuations. In §3, we study high-amplitude fluctuations. The efficiency then approaches 100 %. The analytical results are illustrated with numerical simulations.
## 2. Low-amplitude fluctuations
The size of the central engine, $`r_{\mathrm{in}}10^710^8`$ cm corresponds to a typical time scale $`t_{\mathrm{in}}=r_{\mathrm{in}}/c1`$ ms. The outflow generated by the engine can have highly correlated parameters on this time scale, so that portions $`\mathrm{\Delta }t<t_{\mathrm{in}}`$ can be considered as individual shells in the continuous outflow. The outflow is then discretized as a sequence of such shells, $`i=1,\mathrm{},N`$, with mass $`m_i`$ and energy input $`e_i`$. One after another, the shells are accelerated to a Lorentz factor $`\mathrm{\Gamma }_i=e_i/m_i`$ ($`e_i`$ gets converted into the bulk kinetic energy) on the time-scale $`\mathrm{\Gamma }_it_{\mathrm{in}}<1`$ s (see Piran 1999). Subsequent evolution of the outflow proceeds in the coasting regime until the shells begin to collide.
Suppose $`\mathrm{\Gamma }_i`$ is a white noise with the average value $`\mathrm{\Gamma }_{\mathrm{av}}`$ and initial rms $`\mathrm{\Gamma }_{\mathrm{rms}}^{}{}_{0}{}^{}<\mathrm{\Gamma }_{\mathrm{av}}`$. Masses of the shells can be taken equal, $`m_i=m_0`$, or also fluctuating – the results will be the same. The initial scale of the fluctuations is $`\lambda _0r_{\mathrm{in}}`$. The collision process starts at
$$t_0\frac{\mathrm{\Gamma }_{\mathrm{av}}^3}{\mathrm{\Gamma }_{\mathrm{rms}}^{}{}_{0}{}^{}}\frac{\lambda _0}{c}.$$
(1)
After time $`t2t_0`$, the first generation of shell coalescence has been done (the typical mass of a shell increases by a factor of two), then the second generation occurs, and so on. In the process of hierarchical coalescence, the scale of the fluctuations increases, $`\lambda =\lambda (t)`$, as well as the average mass of an individual shell, $`\lambda (t)/\lambda _0=m(t)/m_0K(t)`$.
The fluctuations look especially simple when viewed from the frame moving with Lorentz factor $`\mathrm{\Gamma }_{\mathrm{av}}`$. We will denote quantities measured in this frame by symbols with tilde. The fluctuation velocity is $`\stackrel{~}{v}_i/c=(\mathrm{\Gamma }_i^2\mathrm{\Gamma }_{\mathrm{av}}^2)/(\mathrm{\Gamma }_i^2+\mathrm{\Gamma }_{\mathrm{av}}^2)`$. Given $`\mathrm{\Gamma }_{\mathrm{rms}}<\mathrm{\Gamma }_{\mathrm{av}}`$, we have $`\stackrel{~}{v}_i/c(\mathrm{\Gamma }_i\mathrm{\Gamma }_{\mathrm{av}})/\mathrm{\Gamma }_{\mathrm{av}}`$ with the average $`\stackrel{~}{v}_{\mathrm{av}}=0`$ and the rms,
$$\frac{\stackrel{~}{v}_{\mathrm{rms}}}{c}\frac{\mathrm{\Gamma }_{\mathrm{rms}}}{\mathrm{\Gamma }_{\mathrm{av}}}.$$
(2)
The outflow motion is thus decomposed into two parts: the relativistic motion of the center-of-momentum (CM) with a Lorentz factor $`\mathrm{\Gamma }_{\mathrm{CM}}\mathrm{\Gamma }_{\mathrm{av}}`$ and superimposed non-relativistic fluctuations.
The hierarchical coalescence of shells can be expected to occur in a self-similar regime, so that the collisional time-scale for one generation is about the time passed since the explosion,
$$\frac{\stackrel{~}{\lambda }}{\stackrel{~}{v}_{\mathrm{rms}}}\stackrel{~}{t}=\frac{t}{\mathrm{\Gamma }_{\mathrm{CM}}}.$$
(3)
Here, $`\stackrel{~}{\lambda }=\mathrm{\Gamma }_{\mathrm{CM}}\lambda `$ is the scale of the fluctuations in the CM-frame. After coalescence of $`K`$ shells with initial momentum $`\stackrel{~}{p}_0m_0\stackrel{~}{v}_{\mathrm{rms0}}`$, the new big shell has a momentum $`\stackrel{~}{p}\sqrt{K}\stackrel{~}{p}_0`$ (this is the random walk formula: $`K`$ momenta $`p_0`$ are summed with random sign). We thus have,
$$m\stackrel{~}{v}_{\mathrm{rms}}\sqrt{K}m_0\stackrel{~}{v}_{\mathrm{rms0}}=\sqrt{\frac{m}{m_0}}m_0\stackrel{~}{v}_{\mathrm{rms0}}.$$
(4)
Combining (2), (3) and (4), we get the self-similar solution describing the hierarchical coalescence at $`t>t_0`$,
$$\frac{m}{m_0}=\frac{\lambda }{\lambda _0}=\left(\frac{t}{t_0}\right)^{2/3},\frac{\mathrm{\Gamma }_{\mathrm{rms}}}{\mathrm{\Gamma }_{\mathrm{rms}}^{}{}_{0}{}^{}}=\left(\frac{t}{t_0}\right)^{1/3}.$$
(5)
Note that the collisions establish a Gaussian distribution of $`\stackrel{~}{v}`$. One can therefore prescribe a temperature to the system of shells, $`T`$, which is related to the average kinetic energy by $`m\stackrel{~}{v}_{\mathrm{rms}}^2/2=kT/2`$ (shells have one degree of freedom). From equation (5) it follows that $`T=const`$. The evolution of shells can thus be described as isothermal sticking together.
The free energy of the outflow is given by
$$U_{\mathrm{free}}=\mathrm{\Gamma }_{\mathrm{CM}}M\frac{\stackrel{~}{v}_{\mathrm{rms}}^2}{2}Mc^2\frac{\mathrm{\Gamma }_{\mathrm{rms}}^2}{2\mathrm{\Gamma }_{\mathrm{CM}}},$$
(6)
where $`M`$ is the total mass of the outflow. This energy is available for dissipation. From equation (5) one gets
$$U_{\mathrm{free}}=U_{\mathrm{free}}^0\left(\frac{t}{t_0}\right)^{2/3}.$$
(7)
In Figure 1, we illustrate solution (5) by direct numerical simulations.
Most of the free energy is dissipated at $`tt_0`$. It should be compared with the time at which the outflow gets optically thin,
$$t_{}\frac{\dot{M}\sigma _\mathrm{T}}{8\pi m_pc^2\mathrm{\Gamma }^2}=\frac{L\sigma _\mathrm{T}}{8\pi m_pc^4\mathrm{\Gamma }^3}2\times 10^2L_{52}\mathrm{\Gamma }_2^3,$$
(8)
where $`\dot{M}`$ is the mass outflow rate and $`L=\dot{M}c^2\mathrm{\Gamma }_{\mathrm{av}}`$ is the kinetic luminosity (assuming a spherical outflow). In the case of $`\mathrm{\Gamma }_{\mathrm{CM}}<\mathrm{\Gamma }_{}180(\mathrm{\Gamma }_{\mathrm{rms}}/\mathrm{\Gamma }_{\mathrm{CM}})^{1/5}L_{52}^{1/5}(\lambda _0/3\times 10^7)^{1/5}`$, dissipation starts at $`t_0<t_{}`$. The free energy remaining in the outflow by the transparency moment is $`U_{\mathrm{free}}^{}=U_{\mathrm{free}}^0(t_{}/t_0)^{2/3}`$, and most of the energy is dissipated at the optically thick stage. The radiation produced at $`t<t_{}`$ is trapped in the plasma whose volume increases proportionally to $`t^2`$. As a result of adiabatic cooling, the volume-integrated radiation energy is reduced as $`t^{2/3}`$. The energy conservation law implies that the adiabatic cooling is accompanied by a regular radial acceleration of shells in the outward direction. The energy of trapped radiation is thus spent to accelerate the outflow CM, and it is lost as a free energy. A fraction $`(t_{}/t)^{2/3}`$ of the radiation energy survives till $`t_{}`$ and contributes to the observed luminosity,
$$L_{}_{t_0}^t_{}\left(\frac{\mathrm{d}U_{\mathrm{free}}}{\mathrm{d}t}\right)\left(\frac{t_{}}{t}\right)^{2/3}dt+U_{\mathrm{free}}^{}.$$
The radiative efficiency (the ratio of $`L_{}`$ to the explosion energy) is
$$ϵ\frac{L_{}}{\mathrm{\Gamma }_{\mathrm{CM}}Mc^2}\frac{A^2}{2}\left(\frac{t_{}}{t_0}\right)^{2/3}\left(\frac{2}{3}\mathrm{ln}\frac{t_{}}{t_0}+1\right).$$
(9)
Here, $`A=\mathrm{\Gamma }_{\mathrm{rms}}^{}{}_{0}{}^{}/\mathrm{\Gamma }_{\mathrm{CM}}`$ is the initial amplitude of the fluctuations.
In the case of $`\mathrm{\Gamma }_{\mathrm{CM}}>\mathrm{\Gamma }_{}`$, we have $`t_0>t_{}`$ and all the radiated free energy will escape. Then,
$$ϵ\frac{U_{\mathrm{free}}^0}{\mathrm{\Gamma }_{\mathrm{CM}}Mc^2}=\frac{A^2}{2}.$$
(10)
If the initial Lorentz factor takes random values between $`\mathrm{\Gamma }_{\mathrm{min}}`$ and $`\mathrm{\Gamma }_{\mathrm{max}}`$ (as assumed in most of the previous works), then
$$A^2=\frac{1}{3}\frac{4\psi }{3(1+\psi )^2},\psi \frac{\mathrm{\Gamma }_{\mathrm{min}}}{\mathrm{\Gamma }_{\mathrm{max}}}.$$
(11)
Since $`A^2<1/3`$ in equation (11), $`ϵ`$ does not exceed $`15`$ %. If one further assumes that only 1/3 of the energy is actually radiated and the other 2/3 are stored as internal energy subject to adiabatic cooling, then the efficiency is reduced by a factor of 3 (2/3 of the dissipated energy is then spent to accelerate the CM). One thus arrives at $`5`$ % limit on $`ϵ`$.
Numerical simulations of Kobayashi et al. (1997; see their Table 1) are in good agreement with equations (10,11), except for the case where $`\psi 0`$ and $`m_i\mathrm{\Gamma }_i^1`$ (their parameter $`\eta =1`$). One can show that in this special case $`\mathrm{\Gamma }_{\mathrm{CM}}/\mathrm{\Gamma }_{\mathrm{rms}}\sqrt{\psi }0`$ and $`ϵ1+\sqrt{\psi }\mathrm{ln}\psi 1`$. The outflow dynamics is then non-linear (even though $`A=1/\sqrt{3}<1`$), similar to the high amplitude case we study in § 3.
## 3. High-amplitude fluctuations
Suppose that $`\mathrm{\Gamma }1`$ fluctuates with a log-normal distribution,
$$\mathrm{ln}\frac{\mathrm{\Gamma }1}{\mathrm{\Gamma }_01}=A\xi ,P(\xi )=\frac{e^{\xi ^2/2}}{\sqrt{2\pi }}.$$
(12)
Here, $`A`$ measures the amplitude of the fluctuations. At $`A<1`$, we have $`(\mathrm{\Gamma }\mathrm{\Gamma }_0)/\mathrm{\Gamma }_0A\xi `$, so that $`\mathrm{\Gamma }_{\mathrm{av}}=\mathrm{\Gamma }_0`$ and $`\mathrm{\Gamma }_{\mathrm{rms}}/\mathrm{\Gamma }_{\mathrm{av}}=A`$. At $`A>1`$ we are in the high-amplitude regime.
The possibility of high efficiencies at $`A>1`$ can be understood when comparing two global characteristics of the outflow: the CM Lorentz factor and the specific energy, $`\eta `$ (in units of $`c^2`$). For definiteness, take a uniform initial distribution of mass. Then,
$$\beta _{\mathrm{CM}}=\frac{P(\xi )\mathrm{\Gamma }\beta d\xi }{P(\xi )\mathrm{\Gamma }d\xi },\mathrm{\Gamma }_{\mathrm{CM}}=\frac{1}{\sqrt{1\beta _{\mathrm{CM}}^2}},$$
(13)
$$\eta =P(\xi )(\mathrm{\Gamma }1)d\xi .$$
(14)
It easy to show that at high amplitudes $`\eta \mathrm{\Gamma }_{\mathrm{CM}}`$, i.e., the chaotic (free) component of the kinetic energy is much larger than the regular component. If the dissipation occurs in the optically thin regime and the emitted radiation is approximately isotropic in the CM-frame, then the momentum conservation law implies $`\mathrm{\Gamma }_{\mathrm{CM}}const`$. When the dissipation is done, the final kinetic energy is about $`(\mathrm{\Gamma }_{\mathrm{CM}}1)Mc^2\eta Mc^2`$, i.e., almost all the explosion energy has been radiated to infinity.
In fact, in the high-amplitude case, the fluctuations start to dissipate very early, much before the transparency moment. The optically thick dissipation accelerates the CM, as discussed in § 2: the free energy gets transformed into the bulk motion of the outflow as a whole. The process is non-linear: the CM is substantially accelerated, in contrast to the low-amplitude case. It tends to reduce the initially big difference between $`\mathrm{\Gamma }_{\mathrm{CM}}`$ and $`\eta `$ by increasing $`\mathrm{\Gamma }_{\mathrm{CM}}`$, while $`\eta `$ stays at the initial value, $`\eta =\eta _0`$.
If at the transparency moment the ratio $`\eta /\mathrm{\Gamma }_{\mathrm{CM}}`$ is still high, then a high efficiency can be expected: the dissipation at the optically thin stage converts the difference between $`\eta Mc^2`$ and $`(\mathrm{\Gamma }_{\mathrm{CM}}1)Mc^2`$ into the observed radiation. Note also, that the radiation produced before the transparency moment is not completely lost. Even being cooled adiabatically, it contributes substantially to the outgoing luminosity.
We now illustrate with numerical simulations. We generate a sequence of $`3\times 10^3`$ thin shells with initial separation 1 ms. Their Lorentz factor fluctuates according to equation (12). The shells have equal mass $`m_0`$. Note that the results depend on the initial mass distribution (in contrast to the linear case), and $`m_i=m_0`$ is taken as a simple example. The duration of the central engine activity is 3 s. The results do not change substantially if one takes longer activity: each 1 s portion of the outflow is causally disconnected from the other portions during the main emission time, $`t<10^4`$ s.
The transparency moment, $`t_{}`$, is roughly estimated by equation (8) with $`\mathrm{\Gamma }\mathrm{\Gamma }_{\mathrm{CM}}`$. The transition from the optically thick to optically thin regime is treated in the simplest way: (1) If two shells merge at $`t<t_{}`$, we assume that no radiation is emitted. Instead, the radiation is trapped and contributes to the kinetic energy of the new big shell. In other words, shell coalescence at $`t<t_{}`$ proceeds with conservation of energy, $`\eta =const`$. A fraction $`(t_{}/t)^{2/3}`$ of radiation trapped at moment $`t`$ survives till $`t_{}`$ and contributes to the luminosity, $`L_{}`$. (2) At $`t>t_{}`$, only specific momentum conserves in the coalescence events. The energy released in the inelastic collision is radiated away isotropically in the rest frame of the newly formed shell.
We consider models of two types. Model I: $`t_{}=200`$ s and $`\mathrm{\Gamma }_0`$ in equation (12) is adjusted in such a way that $`\mathrm{\Gamma }_{\mathrm{CM}}100`$ after the dissipation is finished. Model II: $`t_{}=7`$ s and $`\mathrm{\Gamma }_0`$ is adjusted to get the final $`\mathrm{\Gamma }_{\mathrm{CM}}300`$. In both models, the isotropic kinetic luminosity, $`L10^{52}`$ erg s<sup>-1</sup>.
The outflow forms by $`t=3`$ s and after that it has a well defined CM with a velocity $`\beta _{\mathrm{CM}}=Pc/E`$, where $`P`$ is the total momentum and $`E=(\eta +1)Mc^2`$ is the total energy of the outflow. (Note that the initial $`\eta =\eta _0`$ depends on the specific statistical realization $`\mathrm{\Gamma }_i`$ and it varies around the average expected value given by eq. ). In Figure 2, we show examples of the non-linear evolution of $`\eta `$ and $`\mathrm{\Gamma }_{\mathrm{CM}}`$ in Models I and II. At high $`A`$, the radiation is mostly produced by shells moving faster than the CM, and this results in the CM deceleration at $`t>t_{}`$.
Then, we compute a sequence of models and find the efficiency, $`ϵL_{}/\eta _0Mc^2`$, as a function of $`A`$ (Fig. 3). At $`A<1`$, we are in the linear regime of § 2. Shells begin to collide at $`t_0A^1\mathrm{\Gamma }_0^2(\lambda _0/c)`$ and then evolve according to solution (5). Model II ($`\mathrm{\Gamma }_{\mathrm{CM}}300>\mathrm{\Gamma }_{}`$) is in perfect agreement with equation (10). The efficiency of Model I ($`\mathrm{\Gamma }_{\mathrm{CM}}100<\mathrm{\Gamma }_{}`$) is reduced as a result of adiabatic losses during the optically thick stage.
At $`A>1`$, the efficiency increases up to $`83`$ % in Model I, and up to 96 % in Model II. The global parameters of the outflow show substantial variations depending on the particular realization of the initial $`\mathrm{\Gamma }_i`$. Figure 3 shows $`ϵ`$ averaged for many realizations and its standard deviation.
Figure 4 shows examples of the generated light curves. Each coalescence event produces a pulse of a standard shape corresponding to a thin instantaneously radiating shell. The observed pulse has a width $`\mathrm{\Delta }t_{\mathrm{obs}}t/\mathrm{\Gamma }^2`$, where $`t`$ is the time at which the radiation leaves the shell. Although the initial fluctuations are Gaussian, the non-linear dissipation generates a highly correlated signal which is a mixture of quiescent periods and periods of strong activity. Such a behavior is observed in GRBs and this special feature is naturally explained by the non-linear model.
There is another special feature of the non-linear regime. The effective temperature of fluctuations in the CM-frame is relativistic, and an additional emission mechanism appears: inverse Compton scattering (IC) on the bulk motions (see also Lazzati et al. 1999). This mechanism does not work at the optically thick stage since the radiation is trapped and follows the motions of shells. However, when the outflow gets optically thin, the radiation can propagate and promote momentum exchange between the shells without direct collisions (like Silk damping in the early Universe). This effect should be accounted for in future. We expect that it will not change crucially the efficiency because: (1) Anyway, the free energy is radiated away, whatever dissipation mechanism works. (2) Photon exchange does not crucially enhance the rate of dissipation because direct collisions also occur with relativistic velocities in the CM-frame.
The IC by bulk motions can provide an observational test for the model. In particular, it can explain the high energy tail sometimes observed in GRB spectra. The photons emitted by internal shocks (with energy $`101000`$ keV) are boosted in energy by a factor of $`(\mathrm{\Gamma }_{\mathrm{rms}}/\mathrm{\Gamma }_{\mathrm{CM}})^2`$, resulting in GeV emission.
The study of fluctuating fireballs in this paper was limited to the case of white fluctuations. Fluctuations with an arbitrary spectrum will be studied in a future paper (Beloborodov 2000).
This work was supported by the Swedish Natural Science Research Council and RFBR grant 00-02-16135.
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# Rotational Evolution during Type I X-Ray Bursts
## 1. Introduction
Type I X-ray bursts have been long understood as thermonuclear flashes on the surfaces of neutron stars accreting at rates of $`10^{11}M_{}\mathrm{yr}^1<\dot{M}<10^8M_{}\mathrm{yr}^1`$ in low mass X-ray binaries (LMXBs). The accreted hydrogen and helium accumulates on the surface of the neutron star and periodically ignites and burns. Thermonuclear flash models successfully explain the burst recurrence times (hours to days), energetics ($`10^{39}\mathrm{ergs}`$), and durations ($`10`$$`100\mathrm{s}`$) (Lewin, van Paradijs, & Taam 1993; Bildsten 1998), though many quantitative comparisons to observations are less successful (for example, see Fujimoto et al. 1987; Bildsten 2000).
Evolutionary scenarios connecting the neutron stars in LMXBs to the millisecond radio pulsars (see Bhattacharya 1995 for a review) predict that neutron stars in LMXBs should be spinning rapidly. This has been confirmed for one system, the $`\nu _s=401`$ Hz accreting pulsar SAX J1808.4-3658 (Wijnands & van der Klis 1998; Chakrabarty & Morgan 1998), in which the neutron star magnetic field ($`B10^8`$$`10^9\mathrm{G}`$; Psaltis & Chakrabarty 1999) channels the accretion flow onto the magnetic polar caps, creating an asymmetry which is modulated by rotation. However, most neutron stars in LMXB’s show no evidence for coherent periodicity in the persistent emission, implying that the neutron stars do not possess magnetic fields strong enough to make a permanent asymmetry ($`B10^{10}\mathrm{G}`$).
Type I X-ray bursts have provided a new way to determine the spin of these neutron stars. Observations with the Proportional Counter Array (PCA) on the Rossi X-Ray Timing Explorer (RXTE) of neutron stars in six LMXBs have shown coherent oscillations during Type I X-ray bursts, with frequencies $`\nu _0`$ that range from 300 to 600 Hz (see Table 1). The simplest interpretation is that the burning is not spherically symmetric, providing a temporary asymmetry on the neutron star that allows for a direct measurement of rotation. The coherent nature of the periodicities ($`Q300`$), large modulation amplitudes and stability of the frequency over at least a year support its interpretation as the neutron star spin (Strohmayer 1999b, and references therein).
It is expected theoretically that the burning should not be spherically symmetric. Joss (1978) and Shara (1982) suggested that ignition of a burst occurs at a local spot on the star, and not simultaneously over the whole surface. This is because it takes hours to days to accumulate the fuel, but only a few seconds for the thermal instability to grow. Simultaneous ignition thus requires synchronization of the thermal state of the accreted envelope to one part in $`10^3`$$`10^4`$ over the surface. More likely is that ignition is local and a burning front then spreads laterally (Fryxell & Woosley 1982; Nozakura, Ikeuchi, & Fujimoto 1984; Bildsten 1995), burning the rest of the accreted fuel and creating a temporary brightness asymmetry on the neutron star surface (Schoelkopf & Kelley 1991; Bildsten 1995). Consistent with the picture of a spreading burning front, the pulsation amplitude is observed to decrease during the burst rise while the emitting area increases, reaching a constant value during the decay (Strohmayer, Zhang & Swank 1997b; Strohmayer et al. 1998c).
The oscillations during the burst rise are well-explained by the picture of a spreading hotspot. However, several mysteries remain. First, oscillations are often seen in the burst tail, when the whole surface of the star has presumably ignited. The cause of the azimuthal asymmetry at late times is not understood. Secondly, for those objects with $`\nu _0550\mathrm{Hz}`$ in Table 1, there is evidence that the burst oscillation frequency is twice the spin frequency of the star, including observations of a $`290\mathrm{Hz}`$ subharmonic (Miller 1999) and the fact that the kHz QPO separation is approximately half the burst oscillation frequency in these sources (see van der Klis 2000 for a review). What might cause an $`m=2`$ azimuthal asymmetry is unknown. Third, oscillations are not seen in all sources, or all Type I bursts from the same source.
An initially puzzling feature of the observations was that the oscillation frequency often changes during the burst, increasing by $`\mathrm{\Delta }\nu 1`$$`2`$ Hz in the burst tail. Strohmayer et al. (1997a) proposed a simple explanation — that this frequency shift results from angular momentum conservation. The slight hydrostatic radial expansion (contraction) of the burning layers as the temperature increases (decreases) results in spin down (spin up) if angular momentum is conserved. The time for radial heat transport from the burning layers to the photosphere is about one second, which means that the layers are puffed up and spun down by the time the observer sees the burst. As the burning layers cool during the tail of the burst, they contract and spin up. Strohmayer & Markwardt (1999) and Miller (1999, 2000) have modelled the observed frequency evolution. They find that, in the evolving frame of the burning shell, the oscillations are coherent, as expected if they are due to rotation. Calculations show that the change in thickness of the burning layers during a burst is $`\mathrm{\Delta }z20\mathrm{m}`$ (Hanawa & Fujimoto 1984; Ayasli & Joss 1982; Bildsten 1998). A simple estimate of the spin down of the burning shell due to this change in thickness is then $`\mathrm{\Delta }\nu \nu _s(2\mathrm{\Delta }z/R)1\mathrm{Hz}(\nu _s/300\mathrm{Hz})(\mathrm{\Delta }z/20\mathrm{m})(10\mathrm{km}/R)`$ where $`R`$ is the neutron star radius and $`\nu _s`$ the spin frequency. This roughly agrees with the observed values (Table 1).
In this paper, we make a first attempt at understanding the evolution of the neutron star atmosphere during a Type I X-ray burst on a rotating neutron star. We use hydrostatic models of the neutron star atmosphere to calculate the expansion and resulting spin evolution during a burst. Our aim is to investigate whether the observed frequency changes during bursts are consistent with spin down of the burning shell, thus lending support to the interpretation that the burst oscillation frequency is intimately related to the neutron star spin. This simple picture demands that the hot burning material decouples from the cooler underlying ashes and conserves its angular momentum, completing a few phase wraps with the underlying star during the burst. We thus examine mechanisms that might couple the burning layers to the neutron star, and ask whether it is plausible that the burning layers can remain decoupled for the $`10\mathrm{s}`$ duration of the burst.
We stress that we consider only one-dimensional models for these hydrostatic and coupling calculations. We do not consider the complex question of how the burning front spreads around the star during the burst rise, what determines the number of hotspots on the surface, or what causes the asymmetry at late times during the cooling tail of the burst. We leave these questions for future investigations.
We start in §2 by describing the hydrostatic structure of the atmosphere during fuel accumulation and the X-ray burst. In §3, we calculate the expansion of the atmosphere and show that the expected spin down of a decoupled layer is consistent with observed values. We consider heat transport through the atmosphere, and ask how a single coherent frequency is transmitted to the observer. In §4, we discuss hydrodynamic mechanisms that could transport angular momentum within the burning layers or couple the hot burning material to the underlying colder and denser ashes. In §5, we summarize our results and discuss the many remaining puzzles. Finally, we present our conclusions in §6.
## 2. Thermal Structure and Expansion of the Burning Layers
Most neutron stars in LMXBs accrete hydrogen and helium rich material from their companions at rates $`\dot{M}10^{11}10^8M_{}\mathrm{yr}^1`$. For accretion rates $`\dot{M}2\times 10^{10}M_{}\mathrm{yr}^1`$ (see Bildsten 1998 and references therein), the accumulating hydrogen is thermally stable and burns via the hot CNO cycle of Hoyle & Fowler (1965). The temperature of most of the atmosphere is $`8\times 10^7\mathrm{K}`$ so that the time for a proton capture onto a <sup>14</sup>N nucleus is less than the time for the subsequent beta-decays. This fixes the energy production rate at the value
$$ϵ_H=5.8\times 10^{13}\left(\frac{Z_{\mathrm{CNO}}}{0.01}\right)\mathrm{ergs}\mathrm{g}^1\mathrm{s}^1,$$
(1)
where $`Z_{\mathrm{CNO}}`$ is the mass fraction of CNO nuclei. This energy production rate is independent of temperature or density, and simply set by the beta-decay timescales of <sup>14</sup>O (half-life 71 s) and <sup>15</sup>O (half-life 122 s). Because the hydrogen burns at a constant rate, the time to burn all of it depends only on the metallicity and initial hydrogen abundance. The hydrogen mass fraction $`X`$ in a given fluid element changes at a rate $`dX/dt=ϵ_H/E_H`$, where $`E_H6.7\mathrm{MeV}/\mathrm{proton}6.4\times 10^{18}\mathrm{erg}\mathrm{g}^1`$ is the energy release per gram from burning hydrogen to helium. The time to burn all the hydrogen is then
$$t_H22\mathrm{hours}\left(\frac{0.01}{Z_{\mathrm{CNO}}}\right)\left(\frac{X_0}{0.71}\right),$$
(2)
where $`X_0`$ is the initial hydrogen mass fraction.
The X-ray burst is triggered when helium burning becomes unstable at the base of the accumulated layer, at a density $`10^5`$$`10^6\mathrm{g}\mathrm{cm}^3`$ and temperature $`2\times 10^8\mathrm{K}`$. The composition at the base of the layer depends on how much hydrogen has burned during the accumulation (Fujimoto, Hanawa & Miyaji 1981, hereafter FHM; Bildsten 1998), which is determined by the local accretion rate, $`\dot{m}`$. The local Eddington accretion rate is $`\dot{m}_{\mathrm{Edd}}=2m_pc/(1+X)R\sigma _{\mathrm{Th}}`$, where $`\sigma _{\mathrm{Th}}`$ is the Thomson scattering cross-section, $`m_p`$ is the proton mass, $`c`$ is the speed of light, and $`R`$ is the stellar radius. In this paper, we use the Eddington accretion rate for solar composition ($`X_0=0.71`$) and $`R=10`$ km, $`\dot{m}_{\mathrm{Edd}}8.8\times 10^4\mathrm{g}\mathrm{cm}^2\mathrm{s}^1`$, as our basic unit for the local accretion rate (for a 10 km neutron star, this corresponds to a global rate $`\dot{M}=1.7\times 10^8M_{}\mathrm{yr}^1`$). For $`\dot{m}0.03\dot{m}_{\mathrm{Edd}}`$, the helium burning becomes unstable before all the hydrogen is burned and the helium ignites and burns in a hydrogen rich environment. At lower accretion rates $`\dot{m}0.03\dot{m}_{\mathrm{Edd}}`$, there is enough time to burn all the hydrogen, and a pure helium layer accumulates which eventually ignites. We now describe simple models of the atmosphere immediately prior to (§2.1) and during (§2.2) the X-ray burst for both of these cases. At lower accretion rates still, $`\dot{m}0.01\dot{m}_{\mathrm{Edd}}`$, the hydrogen burning becomes unstable and triggers a mixed hydrogen/helium burning flash (FHM).
### 2.1. The Accumulating Atmosphere
The neutron star atmosphere is in hydrostatic balance as the accreted hydrogen and helium accumulates. The pressure obeys $`dP/dz=\rho g`$, where $`\rho `$ is the density and the gravitational acceleration $`gGM/R^2`$ is constant in the thin envelope. A useful variable is the column depth $`y`$ (units g cm<sup>-2</sup>), defined by $`dy\rho dz`$, giving $`P=gy`$. As the matter accumulates, a given fluid element moves to greater and greater column depth. In this paper, we take $`g_{14}g/10^{14}\mathrm{cm}\mathrm{s}^2=1.9`$, appropriate for a $`M=1.4M_{}`$ and $`R=10\mathrm{km}`$ neutron star. As described above, the hydrogen burning rate is a constant so that the change of $`X`$ with column depth $`y`$ is $`dX/dy=ϵ_H/\dot{m}E_H`$, where we take $`E_H6.4\times 10^{18}\mathrm{erg}\mathrm{g}^1`$ (we have neglected the $`0.5\mathrm{MeV}/\mathrm{proton}`$ lost as neutrinos in the hot CNO cycle, see Wallace & Woosley 1981). Integrating this equation, we find the hydrogen abundance as a function of depth is
$$X(y)=\{\begin{array}{cc}X_0\left[1(y/y_d)\right]\hfill & y>y_d\text{,}\hfill \\ 0\hfill & y<y_d\text{,}\hfill \end{array}$$
(3)
where $`X_0`$ is the initial hydrogen abundance, and the column depth at which the hydrogen runs out is $`y_dX_0\dot{m}E_H/ϵ_H`$, or
$$y_d=6.8\times 10^8\mathrm{g}\mathrm{cm}^2\left(\frac{\dot{m}}{0.1\dot{m}_{\mathrm{Edd}}}\right)\left(\frac{0.01}{Z_{\mathrm{CNO}}}\right)\left(\frac{X_0}{0.71}\right).$$
(4)
If helium ignites at a column depth $`y<y_d`$, a mixed hydrogen/helium burning flash occurs, otherwise a pure helium layer accumulates which eventually ignites at a column $`y>y_d`$.
The thermal profile of the accumulating layer is described by the heat equation,
$$\frac{dT}{dy}=\frac{3\kappa F}{4acT^3},$$
(5)
where $`\kappa `$ is the opacity and $`F`$ is the outward heat flux. The entropy equation is
$$ϵ+\frac{dF}{dy}=c_P\frac{T}{t}+\frac{c_PT\dot{m}}{y}\left(\frac{d\mathrm{ln}T}{d\mathrm{ln}y}n\right)$$
(6)
(Bildsten & Brown 1997), where $`n(d\mathrm{ln}T/d\mathrm{ln}P)_S`$, $`c_P`$ is the heat capacity at constant pressure, and the terms on the right hand side describe the compression of the accumulating matter. The compressional terms contribute $`c_PT5k_BT/2\mu m_p20T_8\mathrm{keV}`$ per accreted nucleon, where $`\mu `$ is the mean molecular weight, and $`T_8T/10^8\mathrm{K}`$. In comparison, the hot CNO cycle hydrogen burning gives $`7X_0\mathrm{MeV}`$ per accreted nucleon for pure helium ignition, or $`7X_0(y_b/y_d)\mathrm{MeV}`$ for mixed H/He ignition, where $`y_b`$ is the column depth at the base. There is additional flux from heat released by electron captures and pycnonuclear reactions deep in the crust, giving $`100\mathrm{keV}`$ per nucleon (Brown & Bildsten 1998; Brown 2000). For the purposes of this paper, we neglect the compressional terms in the entropy equation, and take
$$ϵ=\frac{dF}{dy},$$
(7)
but include a flux at the base ($`y=y_b`$) of $`F_b=150\mathrm{keV}`$ per nucleon as an approximation to the heat from the crust and compressional heating.
We find the temperature profile by integrating equations (5) and (7) from the top of the atmosphere to the base at column depth $`y_b`$. At the top, which we arbitrarily place at $`y=5\times 10^4\mathrm{g}\mathrm{cm}^2`$, we set the temperature using the analytic radiative zero solution for a constant flux atmosphere with Thomson scattering opacity (the solutions are not sensitive to this upper boundary condition). For $`y<y_d`$ we take $`ϵ=ϵ_H`$, and for $`y>y_d`$ we set $`ϵ=0`$. The flux at the top $`F_t`$ is set by the energy release from hot CNO burning in the atmosphere and the flux at the base, giving $`F_t=F_b+ϵ_Hy_b`$ or $`F_t=F_b+ϵ_Hy_d`$, whichever is smaller<sup>1</sup><sup>1</sup>1As pointed out by Taam & Picklum (1978), some helium burning occurs during accumulation, which increases the number of CNO nuclei at ignition (for example, see Fujimoto et al. 1987, Figure 4). FHM and Hanawa & Fujimoto (1982) estimated the increase in CNO abundance, concluding that the effect on the ignition temperature and density would be small because of the large temperature sensitivity of the triple alpha reaction. In calculations which include helium burning during accumulation, we find that helium burning reactions decrease the ignition column by 10–20%, increase the ignition temperature by less than a few percent and increase the metallicity at ignition by factors of 2–3. For the purposes of this paper, we adopt simple models with hot CNO burning only.. The opacity $`\kappa `$ has contributions from electron scattering, free-free absorption and electron conduction, which we calculate as described in Schatz et al. (1999). Since in the hot CNO cycle the seed nuclei are mostly <sup>14</sup>O or <sup>15</sup>O waiting to $`\beta `$-decay, we take the gas to be a mixture of hydrogen (mass fraction $`X`$ given by eq.), <sup>14</sup>O and <sup>15</sup>O (mass fraction $`Z_{\mathrm{CNO}}`$) and helium (mass fraction $`Y=1XZ_{\mathrm{CNO}}`$). The ratio by number of <sup>14</sup>O to <sup>15</sup>O is given by the ratio of the beta-decay timescales, giving $`Z_{14}=0.352Z_{\mathrm{CNO}}`$ and $`Z_{15}=0.648Z_{\mathrm{CNO}}`$. To obtain a simple analytic estimate of the base temperature, we first integrate equation (7) with $`ϵ_H`$ a constant to find $`F(y)=ϵ_H(y_by)`$, where we ignore compressional heating and the flux from the base. Substituting this into equation (5) and integrating assuming constant opacity, we find $`T^4(y)=(3\kappa /ac)ϵ_Hyy_b(1y/2y_b)`$, giving
$`T_b2.6\times 10^8\mathrm{K}\left({\displaystyle \frac{\kappa }{0.08\mathrm{cm}^2\mathrm{g}^1}}\right)^{1/4}\left({\displaystyle \frac{Z_{\mathrm{CNO}}}{0.01}}\right)^{1/4}`$
$`\left({\displaystyle \frac{y_b}{2\times 10^8\mathrm{g}\mathrm{cm}^2}}\right)^{1/2},`$ (8)
where we have inserted a typical value for $`\kappa `$.
We find the column depth at the base of the accumulated layer just before the thermally unstable helium ignition by comparing the temperature sensitivity of the heating and cooling rates (FHM; Fushiki & Lamb 1987b; Bildsten 1998). The heating rate due to the triple-alpha reaction is
$$ϵ_{3\alpha }=5.3\times 10^{21}\mathrm{erg}\mathrm{g}^1\mathrm{s}^1f\frac{\rho _5^2Y^3}{T_8^3}\mathrm{exp}\left(\frac{44}{T_8}\right),$$
(9)
where $`f`$ is the screening factor (Fushiki & Lamb 1987a). In addition, <sup>12</sup>C rapidly captures protons once it is made, increasing the energy release from the triple alpha reaction. To account for this, we multiply $`ϵ_{3\alpha }`$ by a factor $`1+Q_{12}/Q_{3\alpha }=1.9`$, where $`Q_{3\alpha }=7.274`$ MeV, $`Q_{12}=Q(^{12}\mathrm{C}+2p^{14}\mathrm{O})=6.57`$ MeV. An effective local cooling rate is obtained from equations (5) and (7), giving
$$ϵ_{\mathrm{cool}}\frac{acT^4}{3\kappa y^2}.$$
(10)
As the matter accumulates, the column depth at the base $`y_b`$ increases until $`dϵ_{3\alpha }/dT>dϵ_{\mathrm{cool}}/dT`$, at which point a thermal instability occurs (FHM). We find the temperature profile at ignition by choosing $`y_b`$ such that $`dϵ_{3\alpha }/dT=dϵ_{\mathrm{cool}}/dT`$ at the base. In Figure 1, the hatched region in the temperature-column depth plane indicates where helium ignition occurs for abundances ranging from the initial value ($`Y=0.3`$) to pure helium ($`Y=1.0`$).
The conditions at the base of the accumulated column at the time of ignition for a variety of metallicities and accretion rates are shown in Table 2. We have separated our solutions into two groups depending on whether the helium unstably ignites before (mixed H/He ignition) or after (pure He ignition) the hydrogen is completely burned. We give the temperature, density and composition at the base, the time to accumulate the unstable column $`y_b/\dot{m}`$, and the physical distance from the base to the place where $`y=y_b/10`$, $`\mathrm{\Delta }z(90`$%$`)`$. For mixed H/He (pure He) ignition, $`\mathrm{\Delta }z(90`$%$`)5\mathrm{m}`$ ($`3.5`$$`7\mathrm{m}`$).
Figure 1(a) shows the temperature profile of models with mixed H/He ignition. The solid lines are for $`Z_{\mathrm{CNO}}=0.01`$ and (bottom to top) $`\dot{m}/\dot{m}_{\mathrm{Edd}}=0.03,0.1,0.3`$. The dashed lines are for $`\dot{m}=0.1\dot{m}_{\mathrm{Edd}}`$ and $`Z_{\mathrm{CNO}}=0.005`$ (bottom curve) or $`0.02`$ (top curve). Figure 1(b) shows temperature profiles for pure He ignition. The solid lines show $`Z_{\mathrm{CNO}}=0.01`$ and (bottom to top) $`\dot{m}/\dot{m}_{\mathrm{Edd}}=0.01,0.015,`$ and $`0.02`$. The dashed lines are for $`\dot{m}=0.01\dot{m}_{\mathrm{Edd}}`$ and $`Z_{\mathrm{CNO}}=0.005`$ (top curve) or $`0.02`$ (bottom curve). The black dots show the depth where hydrogen runs out (at a column $`y_d`$, eq. ).
For the mixed H/He ignition models, the ignition temperature and density do not depend sensitively on $`\dot{m}`$. Slight differences arise because at higher accretion rates less helium is made from hydrogen burning, requiring a higher density and temperature for ignition. For pure He ignition, the ignition column and temperature are much more sensitive to accretion rate. After the hydrogen is burned, there is a slight temperature gradient to carry the flux from the base $`F_b`$, but the temperature at ignition is mainly set by the temperature at the base of the hydrogen burning shell (FHM; Wallace, Woosley and Weaver 1982). This temperature is greater at higher accretion rates, leading to a smaller column depth at ignition.
The effect of increasing metallicity for mixed H/He ignition models is to increase the energy generation rate, and thus the temperature at the base. Together with the increased rate of helium production, this allows He ignition at a lower column depth, as seen in Figure 1(a). For pure He ignition, the effect of increasing metallicity is exactly opposite. The ignition temperature is lowered with increasing metallicity because the hydrogen runs out at a smaller column, giving a smaller temperature at the base of the hydrogen burning shell. Ignition then requires a greater column depth at the base of the almost isothermal pure helium layer, as seen in Figure 1(b).
The $`\dot{m}`$ at which the transition between pure He and mixed H/He ignitions occurs depends on metallicity. The $`\dot{m}=0.02\dot{m}_{\mathrm{Edd}}`$ case in Figure 1 (right panel) just burns the hydrogen before ignition; this is the transition $`\dot{m}`$ for a metallicity $`Z_{\mathrm{CNO}}=0.01`$. For $`Z_{\mathrm{CNO}}=0.005`$, the transition occurs for $`\dot{m}0.01\dot{m}_{\mathrm{Edd}}`$, and $`\dot{m}0.03\dot{m}_{\mathrm{Edd}}`$ for $`Z_{\mathrm{CNO}}=0.02`$. This is in rough agreement with the $`Z_{\mathrm{CNO}}^{13/18}`$ scaling estimated by Bildsten (1998). The transition accretion rate increases with increasing metallicity because the time to burn all the hydrogen decreases, allowing a pure helium layer to build up even at large accretion rates.
Our results agree well with those of previous authors. The calculations we have presented here are similar to those of Taam (1980) and Hanawa & Fujimoto (1982). Hanawa & Sugimoto (1982) simulated pure He ignition bursts on neutron stars with $`g_{14}=0.93`$ and $`7.1`$. Conditions at the base of the hydrogen burning shell that they find agree with our calculations to 10%, and the helium ignition column agrees to within a factor of two. We find similar agreement with the pure He ignition calculations of Wallace, Woosley & Weaver (1982, hereafter W82). However, the density at the base of the hydrogen burning shell that we compute is twice as great as that given in Table 1 of W82, whereas the pressure and temperature agree. We do not know the reason for this discrepancy, but it may be that W82 used a fixed solar composition for the hydrogen burning shell, giving a different value for $`\mu _e`$ at the base. Hanawa & Fujimoto (1984) simulated mixed H/He ignition bursts. Taking into account their gravity, $`g_{14}=3.3`$, our ignition column and temperature agree with their calculations to 10–30%.
### 2.2. The Atmosphere During the Burst
Many authors have simulated X-ray bursts for both mixed H/He and pure He ignitions, for reviews see Lewin, van Paradijs, & Taam (1993), and Bildsten (1998). In this section, we make simplified models which allow us to calculate the hydrostatic expansion of the atmosphere during a burst. We do not discuss the lateral spreading of the burning front in this paper. In addition, our models are not appropriate for the radius expansion phase of bursts with super-Eddington luminosities (radius expansion bursts; Lewin, van Paradijs, & Taam 1993). We discuss the effect of radius expansion on the observability of oscillations in §3.5. We start in §2.2.1 by considering the convective stage of a burst. As we discuss later, convection is important because it likely enforces rigid rotation, and may affect the observability of a coherent signal (§3). In §2.2.2 we compute fully-radiative atmospheric models, appropriate for those bursts that do not become convective, or for the later stages of a convective burst, when the convection zone retreats.
#### 2.2.1 Models with Convection
In one-dimensional models, the energy release from the temperature sensitive helium burning reaction makes the early stages of many bursts convective. The convection zone expands outwards, extending over a few scale heights to the radiative outermost layers (Joss 1977). For example, Hanawa & Fujimoto (1984) computed a mixed H/He ignition model, in which the atmosphere was convective for $`1\mathrm{s}`$ after ignition, reaching a maximum extent of $`99\%`$ of the accumulated mass. Hanawa & Sugimoto (1982) found for pure He ignition that a convection zone rapidly grew to encompass most of the atmosphere, and then shrunk back, disappearing $`1\mathrm{s}`$ later when the nuclear fuel was almost exhausted.
The convection in the neutron star atmosphere is very efficient, since the time for sound to cross a scale height ($`c_s/g10^6\mathrm{s}`$) is much shorter than the local thermal time ($`1`$$`10\mathrm{s}`$). Thus the atmosphere has a nearly adiabatic profile when convective. For a given temperature $`T_b`$ and column depth $`y_b`$ at the base, the thermal profile of the convective zone just follows the adiabat $`T=T_b(y/y_b)^n`$, where $`n(\mathrm{ln}T/\mathrm{ln}P)_S`$. We take $`y_b`$ to be the value given by the settling solution at ignition; the convection zone profile is then determined by the single parameter $`T_b`$ when the mean molecular weight, $`\mu `$, is fixed.
During the burst, the temperature at the base reaches $`10^9\mathrm{K}`$, high enough that radiation pressure is important (Joss 1977). At these temperatures, the degeneracy is lifted (we find the electron pressure differs from the ideal gas value by less than a few percent) so the total pressure at the base is
$$P_b=\frac{\rho _bk_BT_b}{\mu m_p}+\frac{1}{3}aT_b^4,$$
(11)
where $`\rho _b`$ is the density at the base. However, the total pressure is also fixed by the weight of the overlying atmosphere, which gives $`P_b=gy_b`$. Thus $`T_b`$ cannot exceed the critical value $`T_{\mathrm{max}}=(3gy_b/a)^{1/4}`$, or
$$T_{\mathrm{max}}=2.2\times 10^9\mathrm{K}\left(\frac{y_b}{3\times 10^8\mathrm{g}\mathrm{cm}^2}\right)^{1/4}\left(\frac{g_{14}}{1.9}\right)^{1/4}.$$
(12)
As $`T_b`$ approaches $`T_{\mathrm{max}}`$, radiation pressure becomes increasingly important, forcing the density to decrease. We show below that this greatly enhances the expansion of the atmosphere.
As clearly argued by Joss (1977), the convection zone cannot extend all the way to the photosphere. A radiative layer is always needed to transport the heat to the photosphere. We assume that most of the burning occurs in the convective layer, in which case the radiative atmosphere carries a constant flux and is described by the heat equation (5) with constant $`F`$. For the convective zone, we integrate $`dT/dy=n(T/y)`$ for a given base temperature $`T_b`$ (the super-adiabaticity needed to carry the flux in the convective zone is negligible when the convection is so efficient; see Cox & Giuli 1968). For these simple models, we treat the flux $`F`$ and base temperature $`T_b`$ as free parameters (in reality, the temperature throughout the convective region determines the nuclear energy generation rate and thus the flux). To compute the thermal profile of the atmosphere, we integrate the convection zone outwards from the base, and the radiative zone inwards from the surface, varying the column depth of the top of the convection zone $`y_c`$ until the temperature matches at the interface.
We compute convective models for a range of fluxes and base temperatures for two of the settling solutions of Figure 1. The first is for the mixed H/He ignition model with $`\dot{m}=0.1\dot{m}_{\mathrm{Edd}}`$ and $`Z_{\mathrm{CNO}}=0.01`$. Figure 2(a) shows temperature profiles of convective models with flux $`F/F_{\mathrm{Edd}}=0.31,0.74,0.89`$ and $`1.15`$, the fraction of the accumulated matter that becomes part of the convective zone, $`(y_by_c)/y_b=0.5,0.8,0.99`$ and $`0.95`$, and base temperature $`T_b/10^9\mathrm{K}=1.2,1.5,1.7`$ and $`1.7`$ respectively. Here we measure flux in units of the Eddington flux (for the accreted solar composition) $`F_{\mathrm{Edd}}=GM\dot{m}_{\mathrm{Edd}}/R`$, giving $`F_{\mathrm{Edd}}1.67\times 10^{25}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. Figure 2(b) shows convective models corresponding to the pure He ignition settling solution with $`\dot{m}=0.015\dot{m}_{\mathrm{Edd}}`$ and $`Z_{\mathrm{CNO}}=0.01`$. We make convective models with flux $`F/F_{\mathrm{Edd}}=0.35,0.59,1.06`$ and $`1.21`$, the fraction of the accumulated matter that becomes part of the convective zone, $`(y_by_c)/y_b=0.5,0.8,0.99`$ and $`0.95`$, and base temperature $`T_b/10^9\mathrm{K}=1.3,1.6,2.0`$ and $`1.95`$ respectively. Properties of all these models are given in Table 3. In each case, we take the composition profile of the atmosphere to be that immediately prior to ignition, but assume full mixing of matter in the convective zone (we average over the convection zone according to $`\overline{X_i}𝑑yX_i(y)/(y_by_c)`$ for each species $`i`$).
We give the thickness of the convection zone $`\mathrm{\Delta }z_c`$ in Table 3, as well as the height above the base which contains 90% of the mass, $`\mathrm{\Delta }z(90\%)`$. A simple analytic estimate of the thickness, $`\mathrm{\Delta }z_c`$, is obtained as follows. First, the adiabatic index of the gas is
$$n\left(\frac{d\mathrm{ln}T}{d\mathrm{ln}P}\right)_s=\frac{86\beta }{3224\beta 3\beta ^2},$$
(13)
(Clayton 1983), where the ratio of gas pressure to total pressure is
$$\beta =10.044\left(\frac{T}{10^9\mathrm{K}}\right)^4\left(\frac{y_b}{3\times 10^8\mathrm{g}\mathrm{cm}^2}\right)^1\left(\frac{g_{14}}{1.9}\right)^1.$$
(14)
For $`\beta =1`$ ($`P=P_g`$), we obtain the ideal gas value $`n=2/5`$; for $`\beta =0`$ ($`P=P_r`$), $`n=1/4`$. We take $`PT^n`$, and presume $`\beta `$ is a constant throughout the convective zone, so that $`n`$ is a constant, given by equation (13). Integrating $`dz=dP/\rho g`$, we find the thickness of the convection zone is
$$\mathrm{\Delta }z_c\frac{H_b}{n}\left[1\left(\frac{P_t}{P_b}\right)^n\right],$$
(15)
when it extends to a pressure $`P_t`$. We denote the scale height at the base as $`H_bP_b/\rho _bg=k_BT_b/\mu m_pg\beta `$ and assume $`\mu `$ is constant, giving
$$\mathrm{\Delta }z_c11\mathrm{m}\frac{T_{9,b}}{\mu \beta }\left(\frac{0.4}{n}\right)\left(\frac{1.9}{g_{14}}\right)\left[1\left(\frac{P_t}{P_b}\right)^n\right].$$
Using the value of $`\beta `$ at the base (denoted $`\beta _b`$ in Table 3) to evaluate equation (2.2.1), we find this analytic estimate agrees to better than 5% with the numerical integrations. Equation (2.2.1) shows that radiation pressure strongly affects the radial extent of the convective zone, $`\mathrm{\Delta }z_c1/\beta `$.
We show the density as a function of height above the base for the mixed H/He and pure He ignition models in Figure 3(a) and (b) respectively. The jump in density at the interface between the convective and radiative zones is due to the jump in composition because the convective zone is fully mixed. The black dots mark the height which encloses 90% of the mass as found numerically. This ranges from $`15`$$`50\mathrm{m}`$ for the mixed H/He models, and $`10`$$`30\mathrm{m}`$ for the pure He models (Table 3), compared to $`5\mathrm{m}`$ before the burst (Table 2). The difference between the mixed H/He and pure He models arises because $`\mathrm{\Delta }z_c`$ depends strongly on the mean molecular weight (eq. \[2.2.1\] gives $`\mathrm{\Delta }z_c1/\mu `$). For a mixture of H and He, $`\mu =4/(8X+3Y)`$, giving $`\mu =4/3`$ for pure He, but $`\mu =0.6`$ for solar composition, almost a factor of two different. The mean molecular weight for each model is given in Table 3.
#### 2.2.2 Fully-Radiative Models
We consider two different kinds of radiative models: a radiative atmosphere carrying constant flux, and a radiative atmosphere with a constant energy production rate $`ϵ`$. The first case is relevant when the energy production is concentrated near the base, for example in the initial He burning stages of a burst which does not convect. The second case is relevant when the burning region is more extended. For example, in bursts with mixed H/He ignition, the helium burns rapidly (by purely strong reactions) during the initial stages of the burst, whereas the hydrogen burns more slowly (by the rp-process involving weak reactions; Wallace & Woosley 1981), giving rise to the long tails seen in some X-ray bursts (for a recent example, see Figure 2 of Kong et al. 2000). Because the timescale for H burning by the rp process ($`10`$$`100\mathrm{s}`$) is long compared to the time for radiative heat transport ($`1`$$`10\mathrm{s}`$), the atmosphere is radiative during the rp process tail of the burst (Hanawa & Sugimoto 1982).
We give details of our fully-radiative models in Table 4. For the constant $`ϵ`$ models, we choose $`ϵ`$ to give the required flux at the top, $`F=ϵy_b+F_b`$. We assume the composition of the atmosphere is the same as immediately prior to ignition. Figure 4 shows the temperature profiles. We compute models with $`F/F_{\mathrm{Edd}}=0.31,0.74`$ and $`1.15`$ (mixed ignition) and $`F/F_{\mathrm{Edd}}=0.35,0.59`$ and $`1.21`$ (pure He ignition). Figure 5 shows the density profiles for these models. For the mixed ignition models, $`\mathrm{\Delta }z(`$90%$`)15`$$`33\mathrm{m}`$; for the pure He models, $`\mathrm{\Delta }z(`$90%$`)9`$$`17\mathrm{m}`$. Again, the pure He models are less extended because of the larger $`\mu `$.
## 3. Spin Evolution of the Burning Layers
We have shown that the atmosphere expands outwards by $`10`$$`50\mathrm{m}`$ ($`5`$$`25\mathrm{m}`$) during a mixed H/He (pure He) burst. In this section, we first calculate how the spin of the burning layers evolves during the burst, assuming they conserve angular momentum as they hydrostatically expand and contract, and that they remain decoupled from the bulk of the neutron star (§3.1–3.3). In §3.4, we point out (in agreement with Miller 2000) that for the observer to see a coherent oscillation requires either a mechanism to enforce rigid rotation of the burning layers (such as convection), or that the burning layers be geometrically thin. In addition, we consider the heat propagation through the radiative layers to the surface, and discuss how this can “wash-out” any coherent pulse emanating from deeper layers. We conclude in §3.5 by discussing what happens during radius expansion bursts.
### 3.1. Expansion and Spin Down
We assume that the action of hydrodynamic instabilites (Fujimoto 1988, 1993) or a weak poloidal magnetic field will force the accumulating pile of fuel to be rigidly rotating with the spin frequency of the star, $`\mathrm{\Omega }_{}`$. The time for angular momentum to diffuse across a scale height $`H`$ due to molecular viscosity $`\nu `$ alone is $`t_{visc}H^2/\nu `$. In §5, we show that this time is $`t_{visc}38\mathrm{h}(H/5\mathrm{m})^2`$ for $`T_82`$, comparable to the time between bursts. Thus, even microphysical mechanisms might bring the accumulated fuel into corotation.
We further assume that the burning is spherically symmetric, ignoring the complicated, and important, question of how the burning spreads over the stellar surface, and the fact that some asymmetry is needed to give an observable oscillation. Thus we assume the angular velocity is constant on spherical shells, and depends on radius only. This remains true as the atmosphere expands, since a rigidly rotating spherical shell stays rigidly rotating as it expands outwards (because the fractional change in the distance from the rotation axis is the same for all latitudes). Conservation of angular momentum demands that $`r^2\mathrm{\Omega }(r)`$ remains constant as a fluid element expands outwards, giving the spin frequency as a function of column depth during the burst
$$\mathrm{\Omega }(y)=\mathrm{\Omega }_{}\left(\frac{R+z_1(y)}{R+z_2(y)}\right)^2,$$
(16)
where $`z_1(y)`$ ($`z_2(y)`$) is the height of column $`y`$ above the base before (during) the burst. We take the base at $`y_b`$ to be at fixed radius $`R`$, even though the old ashes below heat up and expand a little. We ignore the general relativistic corrections to the frequency shift, which Strohmayer (1999a) estimates are a $``$10%–20% effect.
Equation (16) shows that the spin frequency depends on depth in the radiative layers. The picture is different in a convective zone. We expect the convective motions will rapidly bring the convection zone into rigid rotation, since the convective turnover time is very short ($`10^3\mathrm{s}`$). Thus, if there is no angular momentum transport between the convective zone and neighboring layers, the spin frequency of the convective zone is
$$\mathrm{\Omega }_c=\mathrm{\Omega }_{}\left(\frac{I_1}{I_c}\right),$$
(17)
where $`I_1`$ is the total moment of inertia of the material between $`y_c`$ and $`y_b`$ before the ignition, and $`I_c`$ is the moment of inertia of the convection zone. We calculate the moment of inertia as follows. A spherical shell at radius $`r`$, mass $`dm=4\pi r^2\rho dr`$ has a moment of inertia $`dI=(8\pi /3)\rho (r)r^4dr`$, where the $`8\pi /3`$ comes from integrating over angles. Since $`\rho drdy`$, this is $`dI=(8\pi /3)r^4dy`$. Thus the moment of inertia is
$$I(y)=\frac{8\pi R^4}{3}_{y_b}^y𝑑y\left(1+\frac{z(y)}{R}\right)^4,$$
(18)
where $`z(y)`$ is the height above the base, given by integrating $`dz/dy=1/\rho `$. We take $`P=gy`$ during this integration, neglecting the small variation of $`g`$ with depth, which has an $`𝒪(H/R)^2`$ effect on the moment of inertia.
We derive a simple analytic estimate of the moment of inertia of the convective zone, by assuming $`\beta =P_g/P`$ is constant, so that $`z(P)`$ is given by an equation similar to equation (15). Inserting this into equation (18), writing $`(1+z/R)^41+4z/R`$ and integrating, we find
$$I_c=\frac{2}{3}\mathrm{\Delta }M_cR^2\left[1+\frac{4H_b}{nR}\left(1\frac{1}{n+1}\frac{1x^{n+1}}{1x}\right)\right],$$
(19)
where $`xy_c/y_b`$, and $`\mathrm{\Delta }M_c4\pi R^2(y_by_c)`$ is the mass of the convection zone. The moment of inertia has two pieces. The zeroth order piece, $`2\mathrm{\Delta }M_cR^2/3`$ is the moment of inertia of mass $`\mathrm{\Delta }M_c`$ concentrated at radius $`R`$. The second term is the $`𝒪(H/R)`$ correction to this because the envelope is extended. Equation (19) reproduces the numerical results to a few percent.
The lower panels of Figure 2 show $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }(y)(\mathrm{\Omega }_{}\mathrm{\Omega }(y))/\mathrm{\Omega }_{}`$ as a function of column depth for the different convective models in Table 3. We define $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ so that a positive value means that the layers have spun down. Figure 4 shows $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }(y)`$ for the fully radiative models of Table 4. The values of $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ that we find are similar to the observed frequency shifts during bursts (Table 1). As we described above, in radiative zones the spin down depends on height, whereas convective zones are rigidly rotating. In the convective models, this gives rise to jumps in $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ at the boundaries between the convection zone and both the underlying ashes and the overlying radiative layer. For both convective and fully-radiative models, the $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }(y)`$ is smaller by roughly a factor of two for pure He ignition as opposed to mixed H/He ignition. This is because the pure He models expand outwards less, as we discussed in §2.2.
### 3.2. Shrinking of the Convection Zone
As we described in §2.2.1, the atmosphere does not remain convective for the whole duration of the burst. Eventually, the convection zone shrinks back. What is the profile of $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ left by the retreating convection zone? Imagine that the convection zone retreats by an amount $`dy_c`$ (so the top of the convection zone moves from column $`y_c`$ to $`y_c+dy_c`$). The material which decouples from the convection zone takes away an amount of angular momentum given by $`dJ_c=8\pi (R+\mathrm{\Delta }z_c)^4\mathrm{\Omega }_cdy_c/3`$, where $`\mathrm{\Omega }_c`$ is the spin frequency of the convection zone, and $`\mathrm{\Delta }z_c`$ is the height of the top of the convection zone above the base. The loss of angular momentum from the convection zone must be compensated by a change in its spin or moment of inertia $`dJ_c=d(I_c\mathrm{\Omega }_c)=I_cd\mathrm{\Omega }_c+\mathrm{\Omega }_cdI_c`$. Dividing by $`dy_c`$, and changing variables to $`xy_c/y_b`$, we obtain the differential equation<sup>2</sup><sup>2</sup>2Since $`I_c`$ depends on both $`H_b`$ and $`x`$, one can write $`dI_c/dx=I_c/x+(dH_b/dx)I_c/H_b`$. In fact, the first term in equation (20) exactly cancels the $`I_c/x`$ piece of the second term, giving $`d\mathrm{\Omega }_c/dx=(\mathrm{\Omega }_c/I_c)(dH_b/dx)I_c/H_b`$. Thus if the convection zone shrank at fixed $`H_b`$, it would not change its spin frequency because the angular momentum taken away by the decoupled fluid exactly cancels the change in the moment of inertia due to the change in $`x`$.
$$\frac{d\mathrm{\Omega }_c}{dx}=\frac{\mathrm{\Omega }_c}{I_c}\left[\frac{8\pi R^4y_b}{3}\left(1+\frac{\mathrm{\Delta }z_c}{R}\right)^4\frac{dI_c}{dx}\right].$$
(20)
We integrate this equation from the starting values of $`\mathrm{\Omega }_c`$ and $`x`$, towards $`x1`$ as the convection zone vanishes. There is one difficulty, however, which is that $`\mathrm{\Delta }z_c`$ and $`I_c`$ are functions of not only $`x`$ but also the scale height at the base $`H_b`$. If we assume the composition $`\mu `$ does not change, then $`H_b`$ only depends on $`T_b`$. We model $`T_b(x)`$ as follows,
$$T_b(x)=T_0+\mathrm{\Delta }T\left(\frac{1x}{1x_i}\right),$$
(21)
where $`T_0`$ is the final temperature after the convection vanishes ($`x1`$), $`T_0+\mathrm{\Delta }T`$ is the initial temperature at the base of the convection zone, and $`x_i`$ is the initial value of $`x`$. In reality, $`x`$ is determined by the instantaneous values of $`F`$ and $`T_b`$, so if we knew $`F(T_b)`$ we could map this onto $`T_b(x)`$. Instead, we specify $`T_b(x)`$ then as a check that our choice is reasonable, we can compute $`F(T_b)`$.
Figure 6 shows the results for the model with $`\dot{m}=0.1\dot{m}_{\mathrm{Edd}}`$, $`Z_{\mathrm{CNO}}=0.01`$ (mixed H/He ignition), with initial extent of the convective zone $`(y_by_c)/y_b=0.8`$, and base temperature $`T_b/10^9\mathrm{K}=1.5`$. As the convection zone retreats, we assume the base temperature falls to $`T_b/10^9\mathrm{K}=1.0`$. We show $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ profiles for the initial model (80% of mass convective), an intermediate case (50%) and when the convective zone has almost vanished (5%). These models have $`T_b/10^9\mathrm{K}=1.5,1.31,`$ and $`1.03`$ and $`F/F_{\mathrm{Edd}}=0.74,0.49,`$ and $`0.21`$ respectively. The jumps in $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ at the top and base of the convection zone persist as it shrinks. In addition, the matter that becomes radiative spins up as the physical thickness of the underlying convection zone decreases. This gives rise to an inversion in $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$, as shown in Figure 6. One might worry that such an inverted profile would be unstable to axisymmetric perturbations. However, all the profiles in Figure 6 have increasing specific angular momentum with radius. In addition, the Brunt-Väisälä frequency is so large, $`N^2\mathrm{\Omega }^2`$4.1), that we expect the Solberg-Hoiland criterion for stability (see e.g. Endal & Sofia 1978) will always be satisfied.
### 3.3. Spin Up in the Cooling Tail
As the atmosphere cools in the tail of a burst, its thickness decreases, and it spins back up. We compute simple models of the cooling atmosphere by assuming it carries a constant flux and has a fixed uniform composition. In Figure 7 we show the thickness $`\mathrm{\Delta }z(`$90%$`)`$ as a function of the flux $`F/F_{\mathrm{Edd}}`$. The solid curve is for the $`\dot{m}=0.1\dot{m}_{\mathrm{Edd}}`$, $`Z_{\mathrm{CNO}}=0.01`$ mixed H/He ignition model, for which we assume a composition of <sup>73</sup>Kr, since the rp process makes elements beyond the iron group (Schatz et al. 1998 and references therein; Koike et al. 1999). The dashed lines are for pure He ignitions, with $`Z_{\mathrm{CNO}}=0.01`$, $`\dot{m}=0.01\dot{m}_{\mathrm{Edd}}`$ (upper curve) and $`\dot{m}=0.015\dot{m}_{\mathrm{Edd}}`$ (lower curve), and for a composition <sup>56</sup>Ni. In each case, we show $`\mathrm{\Delta }z(90`$%$`)`$ just before ignition by a horizontal solid line.
Figure 7 shows that in the limit $`F/F_{\mathrm{Edd}}0`$, the thickness of the ashes is about half the pre-burst thickness for mixed H/He ignition, but similar to the pre-burst thickness for pure He ignition. This is because electrons (which provide most of the pressure support) are consumed in hydrogen burning, but not during helium burning. In the tail of a burst, when, for example $`F/F_{\mathrm{Edd}}0.1`$, the thickness of the atmosphere is $`1\mathrm{m}`$ different from the pre-burst thickness in both cases. Thus if the burning layers remain decoupled as they cool, we expect the spin frequency in the tail to be different from the stellar spin by a part in $`10\mathrm{km}/1\mathrm{m}10^4`$. We discuss the implications of this result in §5.3.
### 3.4. Heat Transport and Coherence of the Oscillation
We now turn to the vertical transport of heat, and its effect on the coherence and amplitude of the burst oscillations. Figures 2 and 4 show that in a radiative zone the magnitude of the spin down depends on depth. However, a single coherent frequency is observed during bursts. This implies that either convection or some other mechanism (see §4) is enforcing rigid rotation in the burning region, or, as pointed out by Miller (2000), the burning layers must be geometrically thin. Since the fractional change in spin across a height $`\mathrm{\Delta }z`$ in a radiative zone is $`2\mathrm{\Delta }z/R`$, an observed coherence $`Q`$ demands a burning layer thickness $`R/2Q`$. Strohmayer & Markwardt (1999) found $`Q4000`$ for bursts from 4U 1728-34 and 4U 1702-43 when the observed frequency evolution was accounted for, while Muno et al. (2000) found $`Q5000`$ for burst oscillations from KS 1731-26. These results imply a burning region thickness $`100\mathrm{cm}`$. This is perhaps more likely for pure He bursts, because of the temperature sensitivity of He burning reactions.
Also important is heat transport across the differentially-rotating atmosphere. The time for radial heat transport from the burning layers to the photosphere can be estimated from the entropy equation, $`c_PT/t=F/y`$, together with equation (5) for the flux $`F`$, which gives a timescale $`t_{\mathrm{therm}}=3\kappa c_Py^2/4acT^3`$. In our numerical calculations, we calculate the heat capacity at constant pressure $`c_P`$ exactly. We make an analytic estimate for a mixture of ideal gas and radiation using $`c_P=(5k_B/2\mu m_p)f(\beta )`$ where $`f(\beta )=(3224\beta 3\beta ^2)/5\beta ^2`$ (Clayton 1983), giving
$$t_{\mathrm{therm}}=0.6\mathrm{s}\frac{f(\beta )}{\mu }\left(\frac{\kappa }{0.08\mathrm{cm}^2\mathrm{g}^1}\right)\frac{y_8^2}{T_9^3}.$$
(22)
The factor $`f(\beta )`$ is unity for $`\beta =1`$ ($`P=P_g`$) and grows increasingly larger as $`\beta `$ approaches zero ($`P=P_r`$), for example $`f(0.85)=2.6`$ whereas $`f(0.4)=27`$ (this is because the internal energy of a photon gas at constant pressure is independent of temperature). Thus at a fixed pressure, the thermal time first decreases with increasing temperature (since $`t_{\mathrm{therm}}1/T^3`$) but starts increasing again once radiation pressure becomes important.
Equation (22) shows that the thermal time ($`1\mathrm{s}`$) is much longer than the time for hydrostatic readjustment ($`10^6\mathrm{s}`$), explaining why spin down is not seen in the beginning of a burst. By the time the heat released by nuclear burning reaches the observer, the layers have expanded and spun down.
The burning layers revolve around the underlying star in a time $`\mathrm{\Delta }\nu ^12\pi /\mathrm{\Delta }\mathrm{\Omega }1\mathrm{s}`$, comparable to the thermal time. The ratio of the shearing time to the thermal time is important. If the heat diffuses quickly, $`t_{\mathrm{therm}}\mathrm{\Delta }\nu 1`$, the atmosphere simply transmits the burning pattern from below, much like passing a flashlight behind a piece of paper. If, on the other hand, heat diffuses slowly, the oscillations will be smeared out. In the limit $`t_{\mathrm{therm}}\mathrm{\Delta }\nu 1`$, the atmosphere is heated uniformly by the quickly revolving burning layers.
Figure 8 shows $`t_{\mathrm{therm}}`$ (solid line) and $`\mathrm{\Delta }\nu ^1`$ (hatched region for $`\nu _s=300`$$`600\mathrm{Hz}`$) as a function of depth for the convective models of Table 3. For mixed H/He and pure He ignitions, the wrap around time $`\mathrm{\Delta }\nu ^1`$ increases for decreasing base temperature and physical thickness. The thermal time, however, decreases with decreasing temperature at a given depth because radiation pressure (eq. ) is less important. Figure 8 shows that the mixed H/He ignition bursts have $`t_{\mathrm{therm}}\mathrm{\Delta }\nu 1`$, whereas the pure He ignition bursts have $`t_{\mathrm{therm}}\mathrm{\Delta }\nu 1`$. This is because both the spin down and thermal time are less for the pure helium bursts ($`\mathrm{\Delta }z_c1/\mu `$ and $`t_{\mathrm{therm}}1/\mu `$). As we discuss further in §5.2, Figure 8 may explain the lack of burst oscillations seen in bursts with the characteristics of mixed H/He ignition. For these mixed H/He ignition bursts, $`t_{\mathrm{therm}}\mathrm{\Delta }\nu 1`$ and any oscillations in the flux from deeper regions can get washed out.
### 3.5. Radius Expansion Bursts
Many of the bursts which have shown oscillations are radius expansion bursts (Table 1), in which super-Eddington luminosities result in expansion of the photosphere to radii $`20`$$`100\mathrm{km}`$ (for a review see Lewin, van Paradijs, & Taam 1993). Burst oscillations have not been observed during the peak of radius expansion bursts, but are sometimes seen during the burst rise, and often during the tail of the burst, once the photosphere has fallen back to the stellar radius $`R`$ (for an example, see Figure 2 of Strohmayer et al. 1998c which shows a radius expansion burst from 4U 1636-54).
We did not model radius expansion in §3. However, several of our models have super-Eddington fluxes, and thus are appropriate for the early stages of these bursts. Because the opacity deep in the atmosphere is smaller than at the photosphere, the hydrostatic structure we calculate breaks down only in the very upper layers near the photosphere. Our radiative models are also appropriate for the cooling tail of these bursts, when the photosphere has fallen back to the stellar surface.
Whatever the cause of the asymmetry on the neutron star surface, the fact that oscillations are not seen during the peak of radius expansion bursts is as we would expect, given the discussion of §3.4. The long thermal time across the extended envelope during radius expansion (Paczynski & Anderson 1986), as well as the fact that the horizontal and vertical lengthscales are similar, will hide any asymmetry and wash out the oscillation.
## 4. Angular Momentum Transport and Recoupling
We have shown that if different layers of the atmosphere conserve their angular momentum, expansion during the burst results in differential rotation within the burning layers. Thus as we noted in §3.4, the fact that a single coherent frequency is observed implies that either the burning layers are geometrically thin, or that they must rotate rigidly. While convection may enforce rigid rotation in the early stages of a burst, we do not expect the atmosphere to be convective in the cooling tail (see §2.2). In this section, we investigate mechanisms which might transport angular momentum between the differentially-rotating burning layers. In addition, we investigate whether the burning layers can remain decoupled from the underlying cold ashes for the $`10\mathrm{s}`$ duration of the burst and the implied several phase wraps.
### 4.1. Kelvin-Helmholtz Instability
Shear layers are notoriously unstable to hydrodynamic instabilities, in particular the Kelvin-Helmholtz instability (Chandrasekhar 1961). The shear may be stabilized by buoyancy, however, if the work that must be done against gravity to mix up the fluid is greater than the kinetic energy in the shear. The importance of buoyancy is measured by the Richardson number, $`\mathrm{Ri}=N^2/(dU/dz)^2`$ (for example, see Fujimoto 1988), where $`N`$ is the Brunt-Väisälä frequency, a measure of the buoyancy. The Kelvin-Helmholtz instability occurs when $`\mathrm{Ri}<1/4`$ (Chandrasekhar 1961; Fujimoto 1988).
There are two sources of buoyancy in the atmosphere, thermal buoyancy and buoyancy due to the composition gradient. We write the Brunt-Väisälä frequency $`N`$ as
$$N^2=\frac{g}{H}\left[\frac{\chi _T}{\chi _\rho }\left(n\frac{d\mathrm{ln}T}{d\mathrm{ln}y}\right)+\frac{\chi _{\mu _e}}{\chi _\rho }\frac{d\mathrm{ln}\mu _e}{dz}+\frac{\chi _{\mu _i}}{\chi _\rho }\frac{d\mathrm{ln}\mu _i}{dz}\right]$$
(23)
(Bildsten & Cumming 1998), where $`\chi _Q\mathrm{ln}P/\mathrm{ln}Q`$ with the other independent thermodynamic variables held constant. The first term of equation (23) is the thermal buoyancy. For an ideal gas atmosphere carrying a constant heat flux in which the opacity is Thomson electron scattering opacity, this is $`N_{th}^2=3\mu m_pg^2/20k_BT`$ (Bildsten 1998), giving
$$\frac{N_{th}}{2\pi }\frac{3\times 10^4\mathrm{Hz}}{T_9^{1/2}}\left(\frac{\mu }{0.6}\right)^{1/2}\left(\frac{g_{14}}{2}\right).$$
(24)
We estimate the composition gradient terms as $`N_\mu ^2g\mathrm{\Delta }\mathrm{ln}\mu /H`$, giving
$$\frac{N_\mu }{2\pi }7\times 10^4\mathrm{Hz}\left(\mathrm{\Delta }\mathrm{ln}\mu \right)^{1/2}\left(\frac{10\mathrm{m}}{H}\right)^{1/2},$$
(25)
where we have taken the density scale height $`H`$ as the lengthscale over which the mean molecular weight changes by $`\mathrm{\Delta }\mathrm{ln}\mu 𝒪(1)`$.
We show the Brunt-Väisälä frequency during the burst for two fully-radiative models in the top panel of Figure 9. We show the mixed H/He ignition model with $`F=0.83F_{\mathrm{Edd}}`$, and the pure He ignition model with $`F=0.59F_{\mathrm{Edd}}`$ (see Table 4). We show the total buoyancy by a solid line, the thermal buoyancy as a dotted line and the composition piece as a dashed line. For the pure He ignition model, there is a peak in the buoyancy at the place where the hydrogen runs out ($`yy_d`$), at this depth the composition piece of the buoyancy dominates the thermal piece. Our models do not include the composition jump at the base between the burning layers and the ashes. The point marked by a cross at the base of each model shows an estimate of the buoyancy at the base, where we estimate the composition piece of the buoyancy using equation (25) and take the ashes to have $`\mu =2.1`$ (for a single species, $`\mu =A/1+Z`$).
The Richardson number $`\mathrm{Ri}=N^2/(dU/dz)^2`$ is plotted in the lower panel of Figure 9 for $`\nu _s=300`$$`600\mathrm{Hz}`$. Since $`\mathrm{\Delta }\mathrm{\Omega }(2z/R)\mathrm{\Omega }`$, we estimate $`dU/dz=Rd\mathrm{\Delta }\mathrm{\Omega }/dz2\mathrm{\Omega }`$, giving
$$\mathrm{Ri}7000\left(\frac{\nu _s}{300\mathrm{Hz}}\right)^2\left(\frac{N}{5\times 10^4\mathrm{Hz}}\right)^2,$$
(26)
which agrees well with Figure 9. As suggested by Bildsten (1998), the strong buoyancy in the atmosphere gives $`\mathrm{Ri}1`$. Thus we do not expect the Kelvin-Helmholtz instability to operate either within the burning layers or between the burning layer and the underlying ashes.
### 4.2. Ekman Pumping
We now investigate how fast viscosity acts to smooth out differential rotation. The molecular viscosity in the atmosphere is determined by ion-ion collisions, giving
$$\nu 88\mathrm{cm}^2\mathrm{s}^1\left(\frac{T_9^{5/2}}{\rho _5}\right)\left(\frac{A^{1/2}}{Z^4}\right)\left(\frac{8}{\mathrm{log}\mathrm{\Lambda }}\right)$$
(27)
(Spitzer 1962), where $`T_9T/10^9\mathrm{K}`$, $`\rho _5\rho /10^5\mathrm{g}\mathrm{cm}^3`$ and $`\mathrm{log}\mathrm{\Lambda }`$ is the Coulomb logarithm. Thus the time for molecular viscosity to transport angular momentum over a scale height by diffusion is
$$t_\nu 3.2\mathrm{h}\left(\frac{H}{10\mathrm{m}}\right)^2\left(\frac{\nu }{88\mathrm{cm}^2\mathrm{s}^1}\right)^1,$$
(28)
much longer than a rotation period. In this case, it is possible to exchange angular momentum on a faster timescale than $`t_\nu `$ by the process of Ekman pumping. This mechanism involves a secondary circulation in which fluid elements are exchanged between the bulk of the fluid and the thin viscous boundary layer. Ekman pumping is well-studied in fluid dynamics (Benton & Clark 1974), and Livio & Truran (1987) suggested that it may operate in accreting white dwarfs. For a non-stratified fluid, the Ekman spin up or spin down time is $`t_E(t_\nu /\mathrm{\Omega })^{1/2}`$ or
$$t_E2.5\mathrm{s}\left(\frac{t_\nu }{3.2\mathrm{h}}\right)^{1/2}\left(\frac{\nu _s}{300\mathrm{Hz}}\right)^{1/2}$$
(29)
(Greenspan & Howard 1963; Benton & Clark 1974). The thickness of the viscous boundary layer is $`(\nu /\mathrm{\Omega })^{1/2}0.2\mathrm{cm}(\nu /88\mathrm{cm}^2\mathrm{s}^1)^{1/2}(\nu _s/300\mathrm{Hz})^{1/2}`$.
Spin up in a stratified fluid is different since the buoyancy may inhibit vertical motion of fluid elements, limiting the extent of the secondary circulation. Spin up in a cylinder of radius $`a`$ with $`g`$ parallel to $`\mathrm{\Omega }`$ was first studied by Holton (1965; see Benton & Clark 1974 for a review). In this case, the secondary circulation is confined to a layer of vertical thickness $`a(\mathrm{\Omega }/N)`$ (Walin 1969). Sakurai, Clark & Clark (1971) found a similar result for a spherical stratified fluid, for which Ekman spin up occurs in a thin layer of radial extent $`R(\mathrm{\Omega }/N)`$. The neutron star atmosphere during a burst has $`H<R(\mathrm{\Omega }/N)60\mathrm{m}`$. Thus, if the spherical result is applicable to a thin layer on the surface of a sphere, we would not expect the buoyancy to inhibit Ekman pumping. In addition, we have $`t_{\mathrm{therm}}t_E`$, so non-adiabatic effects may be important, reducing the effect of the thermal buoyancy (Sakurai et al. 1971). Detailed calculations are needed to find whether Ekman pumping operates during the burst.
### 4.3. Baroclinic Instability
The baroclinic instability is a hydrodynamic instability in which the fluid motions are close to horizontal, so it occurs even when the Richardson number is large. It has been well-studied in geophysics (Pedlosky 1987) and in astrophysics because of its possible role in angular momentum transport in stellar interiors (Knobloch & Spruit 1982; Tassoul & Tassoul 1982; Spruit & Knobloch 1984) and accreting compact objects (Fujimoto 1988, 1993). We start by showing that in the presence of a vertical shear, the surfaces of constant pressure and density are inclined with respect to each other. This represents a store of gravitational energy which the baroclinic instability can tap. We then present the results of a stability analysis of a plane-parallel model, for which we find that the strong buoyancy limits the baroclinic instability to short vertical wavelengths ($`200\mathrm{cm}`$).
#### 4.3.1 The Nature of the Baroclinic Instability
We first study the misalignment of the constant density and pressure surfaces that arises in the presence of vertical shear on a rotating star. This result is well-known on the Earth, where differential heating between the equator and pole gives rise to the “thermal wind” (Pedlosky 1987). To investigate recoupling of the ashes and burning layers, we adopt a simple “two layer” model, in which we include only the buoyancy associated with the interface between the ashes and the burning layers. Thus, we take the upper (lower) layer to have constant density $`\rho _+`$ ($`\rho _{}`$). Both layers are in hydrostatic balance in the vertical direction. As elsewhere in this paper, we assume the burning front has spread over the whole surface.
We work in the rotating frame, in which the lower fluid is stationary. We assume the upper layer is rigidly rotating with angular velocity $`\mathrm{\Delta }\mathrm{\Omega }`$ (so that $`\mathrm{\Delta }\mathrm{\Omega }`$ is positive for a spun down shell). Since it is moving, the upper fluid feels a Coriolis force in the transverse direction. In a time $`\mathrm{\Omega }^110^3\mathrm{s}`$, a transverse pressure gradient will be established which balances the Coriolis force (geostrophic balance). However, in the lower fluid there is no transverse pressure gradient (it feels no Coriolis force). Since the pressure must be continuous across the boundary, it cannot be horizontal, but must slope.
Figure 10(a) illustrates this. We show a small section of the two layers at latitude $`\theta `$ ($`\theta `$ is the angle from the pole) in the $`(r,\theta )`$ plane. Because the star is rotating, the equilibrium isobars are not spherical, but nearly (the stellar radius is larger at the equator than the pole by $`R^2\mathrm{\Omega }^2/g100\mathrm{m}`$ for typical parameters). We adopt coordinates $`z`$ and $`y`$ perpendicular and parallel to surfaces of constant pressure (dotted line). The vertical component of the rotation vector is $`\mathrm{\Omega }\mathrm{cos}\theta `$, and the upper fluid is moving out of the page with velocity $`U=R\mathrm{\Delta }\mathrm{\Omega }\mathrm{sin}\theta `$.
Now consider the pressure changes as we move from point A to point B on the boundary. First we move through the lower fluid along the dashed line, horizontally a distance $`dy`$, then vertically upwards a distance $`dz`$. The pressure change along this path is $`dP_{}=\rho _{}gdz`$. Second, we move through the upper fluid along the dashed line. This time, there is a pressure change while moving horizontally also, giving $`dP_+=\rho _+gdz\rho _+2\mathrm{\Omega }U\mathrm{cos}\theta dy`$. Demanding $`dP_{}=dP_+`$, gives the slope
$$\frac{dz}{dy}=\frac{2\mathrm{\Omega }^2R}{g}\frac{\mathrm{\Delta }\mathrm{\Omega }}{\mathrm{\Omega }}\frac{\rho _+}{\mathrm{\Delta }\rho }\mathrm{cos}\theta \mathrm{sin}\theta ,$$
(30)
where $`\mathrm{\Delta }\rho =\rho _{}\rho _+>0`$. If we take the equilibrium equipotentials to be spherical ($`dyRd\theta `$) and integrate, we find the change in height of the boundary is $`z_b(\theta )=\overline{z}_b\mathrm{cos}2\theta +\mathrm{const}.`$, where $`\overline{z}_b=(\mathrm{\Omega }^2R^2/2g)(\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega })(\rho _+/\mathrm{\Delta }\rho )`$ or
$$\overline{z}_b=35\mathrm{cm}\left(\frac{\nu _s}{300\mathrm{Hz}}\right)^2\left(\frac{\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }}{4\times 10^3}\right)\left(\frac{\rho _+}{\mathrm{\Delta }\rho }\right),$$
(31)
where we have taken $`g_{14}=2`$ and $`R=10\mathrm{km}`$.
The sloping interface represents a store of gravitational potential energy. Assuming constant density, the displaced mass per unit area is $`\rho z_b(\theta )`$, and the gravitational energy per unit area required to displace the interface is $`\rho gz_b(\theta )^2/2`$. Integrating over a spherical surface, we find the gravitational energy stored in the interface is $`\mathrm{\Delta }E_g=(16/45)(4\pi R^2\rho _+g\overline{z}_b^2/2)(\mathrm{\Delta }\rho /\rho _+)`$, or
$`\mathrm{\Delta }E_g=6\times 10^{34}\mathrm{erg}\left({\displaystyle \frac{\mathrm{\Delta }\rho }{\rho _+}}\right)\left({\displaystyle \frac{\overline{z}_b}{35\mathrm{cm}}}\right)^2`$
$`\left({\displaystyle \frac{\rho _+}{10^5\mathrm{g}\mathrm{cm}^3}}\right)\left({\displaystyle \frac{g_{14}}{2}}\right).`$ (32)
The amount of displaced mass is $`4\pi R^2\rho \overline{z}_b4\times 10^{19}\mathrm{g}`$, so $`E_g`$ is only a few keV per nucleon, much less than the energy produced by nuclear burning. Thus establishing the sloping interface poses no energetic obstacle. For comparison, the energy in the shear is $`\rho H_bU^2/2`$ per unit area, giving $`E_U(8\pi R^2/3)\rho H_b(R\mathrm{\Delta }\mathrm{\Omega })^2/2`$ or
$`E_U3\times 10^{34}\mathrm{erg}\left({\displaystyle \frac{\rho }{10^5\mathrm{g}\mathrm{cm}^3}}\right)\left({\displaystyle \frac{H_b}{10\mathrm{m}}}\right)`$
$`\left({\displaystyle \frac{\nu _s}{300\mathrm{Hz}}}\right)^2\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }}{4\times 10^3}}\right)^2,`$ (33)
when integrated over the surface.
The baroclinic instability acts to release the gravitational potential energy stored in the misaligned pressure and density surfaces. Consider the fluid displacements shown in Figure 10(b). If the fluid element is displaced to point A, it is heavier than its new surroundings and is pushed back by the buoyancy of the interface. However, if it is moved to point B, it falls in the gravitational field, releasing energy. Any displacement within the so-called “wedge of instability” (Pedlosky 1987) is convectively unstable in this way.
#### 4.3.2 Results of Stability Analysis
We have carried out a linear stability analysis of a plane parallel two-layer model. We do not present our detailed calculations, rather we summarize our results and use simple physical arguments to understand them. We start by considering motions about a latitude $`\theta =\theta _0`$, and adopting a local cartesian coordinate system $`(x,y,z)`$ where the transverse coordinate $`x`$ ($`y`$) is in the $`\varphi `$ ($`\theta `$) direction. We include the effect of sphericity using the “beta-plane approximation” of geophysics (Pedlosky 1987), i.e. we write
$$2\mathrm{\Omega }\mathrm{sin}\theta 2\mathrm{\Omega }\mathrm{sin}\theta _0+(\theta \theta _0)\mathrm{cos}\theta _0,$$
(34)
where we assume $`(\theta \theta _0)=y/R1`$. We consider a channel centered on $`\theta _0=\pi /4`$ and of width $`\pi R/4`$, stretching from the pole to the equator. We write the fluid displacement as ($`\xi _x`$,$`\xi _y`$,$`\xi _z`$), and look for solutions $`\mathrm{exp}(\sigma t+ik_xx+ik_yy)`$ where $`k_x=2\pi /\lambda _x`$ and $`k_y=2\pi /\lambda _y`$ are the transverse wavenumbers, and $`\sigma `$ is the growth rate. Pedlosky (1987) performs a similar analysis, although restricted to the case $`\mathrm{\Delta }\rho /\rho 1`$.
We find that small wavelengths are stable because the fluid displacements lie outside the wedge of instability. To see this, assume that the perturbations are in geostrophic balance, in which case the continuity equation gives $`\xi _z/\lambda _z\mathrm{Ro}\xi _{}/\lambda _{}`$, where $`\mathrm{Ro}U/2\mathrm{\Omega }\lambda _{}R\mathrm{\Delta }\mathrm{\Omega }/2\mathrm{\Omega }\lambda _{}`$ is the Rossby number of the perturbation, and the transverse wavelength $`\lambda _{}=2\pi /k_{}`$, with $`k_{}^2=k_x^2+k_y^2`$. Then the angle of fluid displacement is $`\xi _z/\xi _y\xi _z/\xi _{}R\mathrm{\Delta }\mathrm{\Omega }H/\lambda _{}^2\mathrm{\Omega }`$, where the vertical wavelength is set by the vertical extent of the burning layers $`\lambda _zH`$. For instability, the angle of fluid displacement should be less than the slope of the interface (eq. ), requiring $`\lambda _{}^2(gH_b/\mathrm{\Omega }^2)(\mathrm{\Delta }\rho /\rho _+)`$, or
$`\lambda _{}>\lambda _{BC}=7.5\times 10^5\mathrm{cm}\left({\displaystyle \frac{H_b}{10\mathrm{m}}}\right)^{1/2}\left({\displaystyle \frac{g_{14}}{2}}\right)^{1/2}`$
$`\left({\displaystyle \frac{\mathrm{\Delta }\rho }{\rho _+}}\right)^{1/2}\left({\displaystyle \frac{\nu _s}{300\mathrm{Hz}}}\right)^1,`$ (35)
where we insert the correct prefactor. Thus the nearly horizontal unstable displacements require transverse wavelengths of order the stellar radius or greater.
However, very large transverse wavelength perturbations can be stabilized because the vertical component of $`\mathrm{\Omega }`$ changes significantly across a wavelength. The changing vorticity provides a restoring force in the $`\theta `$ direction (Pedlosky 1987). This is the same force which supports Rossby waves (Pedlosky 1987; Brekhovskikh & Goncharov 1994; Dutton 1995). We find that the convective instability overcomes the Rossby wave restoring force when the shear velocity is greater than the Rossby wave speed, $`U>U_{\mathrm{Ro}}2\mathrm{\Omega }\lambda _{}^2/R`$ (see also Pedlosky 1987). Thus we require $`\lambda _{}R(\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega })^{1/2}`$, or
$$\lambda _{}<\lambda _{\mathrm{Ro}}=3.3\times 10^5\mathrm{cm}\left(\frac{\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }}{4\times 10^3}\right)^{1/2}\left(\frac{R}{10\mathrm{km}}\right)$$
(36)
for instability, where we insert the correct prefactor.
Transverse wavelengths which satisfy equation (4.3.2) do not satisfy equation (36). Thus in the context of our plane-parallel model, we conclude that, when $`H=10\mathrm{m}`$, large scale mixing between the burning layers and ashes by the baroclinic instability does not occur. Instability can occur when $`\lambda _{\mathrm{BC}}<\lambda _{\mathrm{Ro}}`$, or $`\mathrm{\Delta }\rho /\rho _+(4\mathrm{\Omega }^2R/g)(\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega })(R/H)`$. Since $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }2H/R`$, this is
$$\frac{\mathrm{\Delta }\rho }{\rho _+}0.15\left(\frac{\nu _s}{300\mathrm{Hz}}\right)^2\left(\frac{g_{14}}{2}\right),$$
(37)
so mixing may occur as the burning layers cool and the density contrast with the ashes decreases. The growth rate of an unstable mode is $`k_x`$, so that writing $`k_xm/R`$ (so that $`e^{ikx}e^{im\varphi }`$), we estimate the fastest growing modes will have $`m20(\lambda _{}/3\times 10^5\mathrm{cm}`$). Several authors have studied two layer models on a sphere (Hollingsworth 1975; Hollingsworth, Simmons & Hoskins 1976; Simmons & Hoskins 1976, 1977; Moura & Stone 1976; Warn 1976), but for parameters of interest for the Earth, namely $`\mathrm{\Delta }\rho /\rho 1`$ and $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }0.01`$. They find good agreement with growth rates calculated with beta-plane models. However, further calculations are needed to extend our analysis to the spherical case, and determine the spherical eigenfunctions and growth rates.
#### 4.3.3 Short Wavelength Modes
So far we have considered only vertical wavelengths $`\lambda _zH`$, because of the simple vertical structure of the two layer model. However, if we allow perturbations with $`\lambda _z<H`$, we expect to find instability for small enough vertical wavelengths. This is because for small vertical wavelengths, the fluid displacement is able to lie within the wedge of instability (Figure 10) while still having a transverse wavelength which is unaffected by Rossy wave restoring forces (eq. ).
We make a simple estimate of the vertical wavelength at which instability occurs by repeating the arguments leading to equation (4.3.2), but this time allowing $`\lambda _z<H_b`$. In this case, $`\lambda _{BC}`$ is given by equation (4.3.2) with $`H_b`$ replaced by $`\lambda _z`$. For instability, we require $`\lambda _{\mathrm{BC}}<\lambda _{\mathrm{Ro}}`$ (eq. ) giving
$$\lambda _z200\mathrm{cm}\left(\frac{\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }}{4\times 10^3}\right)\left(\frac{\nu _s}{300\mathrm{Hz}}\right)\left(\frac{\rho _+}{\mathrm{\Delta }\rho }\right),$$
(38)
where we take $`g_{14}=2`$ and $`R=10\mathrm{km}`$. We expect modes with vertical wavelengths small enough to satisfy equation (38) to be unstable.
Fujimoto (1988, 1993) used the short wavelength baroclinic modes to define a turbulent viscosity $`\nu _{\mathrm{turb}}\lambda _z^2/t_{\mathrm{grow}}`$ where the growth time of the instability is $`t_{\mathrm{grow}}\nu _s^1\mathrm{Ri}^{1/2}`$, or
$$t_{\mathrm{grow}}0.25\mathrm{s}\left(\frac{\nu _s}{300\mathrm{Hz}}\right)^1\left(\frac{\mathrm{Ri}}{7000}\right)^{1/2}$$
(39)
(Pedlosky 1987; Fujimoto 1988). The time to transport angular momentum across a scale height is $`t_\nu H^2/\nu _{\mathrm{turb}}t_{\mathrm{grow}}(H/\lambda _z)^2`$, or
$$t_\nu 6\mathrm{s}\left(\frac{t_{\mathrm{grow}}}{0.25\mathrm{s}}\right)\left(\frac{H}{10\mathrm{m}}\right)^2\left(\frac{\lambda _z}{200\mathrm{cm}}\right)^2$$
(40)
Thus turbulent transport of angular momentum driven by the short wavelength baroclinic modes could be important during a burst. In particular, these could act to force the burning shell to rigidly rotate.
### 4.4. Magnetic Field Winding
The majority of neutron stars in LMXBs show no evidence for coherent pulsations in their persistent emission. This implies that these neutron stars are weakly magnetic ($`B10^9\mathrm{G}`$) so the accretion flow is likely not disrupted before it reaches the neutron star surface. We now ask what effect would a weak magnetic field have on the shearing atmosphere during a burst? The Ohmic diffusion time across a scale height is $`t=4\pi \sigma H^2/c^2`$, or
$$t_{\mathrm{diffuse}}\frac{10^7\mathrm{s}}{Z}\left(\frac{H}{10\mathrm{m}}\right)^2\left(\frac{T}{10^9\mathrm{K}}\right)^{3/2}\left(\frac{8}{\mathrm{log}\mathrm{\Lambda }}\right)$$
(41)
(Brown & Bildsten 1998) where we estimate the conductivity as
$$\sigma \frac{2(2k_BT)^{3/2}}{\pi ^{3/2}m_e^{1/2}Ze^2\mathrm{log}\mathrm{\Lambda }}$$
(42)
(Spitzer 1962), $`Z`$ is the mean ionic charge and $`\mathrm{log}\mathrm{\Lambda }`$ is the Coulomb logarithm. Since $`t_{\mathrm{diffuse}}10\mathrm{s}`$, the MHD limit applies during the burst and the shearing atmosphere will bend the magnetic field lines. An initially poloidal magnetic field would prevent the shearing if the energy density in the field is greater than the shear energy, $`B^2/8\pi >\rho (R\mathrm{\Delta }\mathrm{\Omega })^2/2`$, or
$$B>10^{10}\mathrm{G}\rho _5^{1/2}\left(\frac{R}{10\mathrm{km}}\right)\left(\frac{\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }}{4\times 10^3}\right).$$
(43)
This is a similar limit to those from lack of persistent pulsations and from spectral modelling (Psaltis & Lamb 1998). The observed shearing for ten seconds most likely rules out fields this strong.
What can we say about a weaker, initially poloidal field like that often discussed for the the progenitors of millisecond pulsars, $`B10^810^9\mathrm{G}`$? Fields this weak can also have a significant effect when in the ideal MHD limit. In this case, the horizontal shearing and dragging of the field lines generates a strong toroidal field. Spruit (1999) recently considered the interaction of magnetic fields and differential rotation in stellar interiors. He considered an initially poloidal field $`(B_z,B_\theta ,0)`$, and showed that $`B_\varphi `$ grows with time according to $`B_\varphi =NB_z`$ where $`N`$ is the number of windings, $`N=R\mathrm{sin}\theta |\mathbf{}\mathrm{\Omega }|𝑑t`$. We have $`\mathbf{}\mathrm{\Omega }\mathrm{\Delta }\mathrm{\Omega }/H`$, giving $`N(R\mathrm{\Delta }\mathrm{\Omega }/H)t2\mathrm{\Omega }t`$. Thus we find $`B_\varphi B_z`$ after only a few revolutions of the star, much less than the wrap around time of the shear ($`1\mathrm{s}`$). The toroidal field grows so quickly because the field is sheared on a vertical scale $`HR`$. Another way to see this is to consider a vertical tube of fluid in the shearing region (with height $`H`$ and cross-section area $`A_i`$) that encloses a magnetic flux $`\mathrm{\Phi }=B_zA_i`$. After one differential wrapping time (roughly one second), this fluid element will be stretched into a very thin tube of length $`2\pi R`$ and so will have shrunk in cross-sectional area to $`A_f=A_iH/2\pi R`$. If flux-freezing holds, then after one second, $`B_\varphi B_zA_i/A_fB_z2\pi R/H10^4B_z`$.
This simple argument says that an initial field of $`10^6\mathrm{G}`$ can become dynamically important during the burst. However, further theoretical investigations are needed to see if the implied wrapping of the field is possible without instabilities setting in. If correct, this result implies that the bursters with frequency drift have magnetic fields much weaker than the $`B10^8`$$`10^9\mathrm{G}`$ fields usually assumed.
### 4.5. Summary of Coupling Mechanisms
We have investigated a number of different mechanisms that could enforce rigid rotation within the burning layers, or act to couple the burning layers to the underlying ashes. We find that the strong buoyancy of the atmosphere prohibits the Kelvin-Helmholtz instability and mixing between the ashes and the burning layers by the baroclinic instability. Short vertical wavelength baroclinic modes within the burning region may be unstable and vigorous enough to force the burning layers to rigidly rotate. The timescale for viscosity to act via an Ekman pumping mechanism may be a few seconds or less (eq. ), but the effect of buoyancy, in particular from the composition gradients, is not clear. Finally, our initial estimates suggest that winding of a weak magnetic field by the differential rotation is important during the burst.
In summary, we have not found a robust hydrodynamical mechanism that recouples the burning layers to the star during the $`10\mathrm{s}`$ burst duration. It thus seems plausible that the burning layers remain decoupled throughout the burst. Short wavelength baroclinic instabilities may operate within the burning layers, although more work is needed to see whether they cause the burning shell to rigidly rotate.
## 5. Discussion
We now summarize the hydrostatic expansion of the atmosphere during a Type I X-ray burst, and the implied spin evolution of the decoupled burning layers. We find the expected spin changes are of the same order as the observed frequency drifts, thus supporting the picture proposed by Strohmayer et al. (1997a). We then use our results to address two observational issues: (i) why oscillations are not seen in all bursts, and (ii) the observed long term stability of the oscillation frequency.
### 5.1. Temporal Evolution of the Burning Shell
The atmospheric evolution during the burst is shown schematically in Figure 11 for (a) mixed H/He ignition and (b) pure He ignition. As revealed by one-dimensional simulations, the rapid energy release from helium burning reactions makes the initial stages of most bursts convective. Later in the burst, the convection zone retreats, leaving a purely radiative atmosphere. Figure 12 shows the thickness of the atmosphere which contains 90% of the mass, $`\mathrm{\Delta }z(90`$%$`)`$, as a function of the flux for constant flux radiative atmospheres (solid and dotted lines) and for the convective models of Table 3 (solid and open squares). For the radiative models, the upper curve is for a composition profile the same as at ignition, while the lower curve is for a composition of <sup>73</sup>Kr (<sup>56</sup>Ni) for the mixed H/He (pure He) case, chosen to represent the products of burning. Initially, the evolution of the thickness of the layer is along or above the upper curve as it ignites and heats up, depending on the extent of the convection zone. As nuclear burning proceeds, the mean molecular weight increases, and the thickness decreases, eventually moving back along the lower curve as the atmosphere cools. Just before ignition (upper curve, $`FF_{\mathrm{Edd}}`$), the thickness of the accumulated layer is $`5\mathrm{m}`$. During the burst, the atmosphere expands hydrostatically by $`\mathrm{\Delta }z10`$$`40\mathrm{m}`$ ($`5`$$`20\mathrm{m}`$) for mixed H/He (pure He) ignition. For a given flux, the convective models have a greater thickness than fully-radiative models because the temperature profile is steeper in the convection zone, giving a larger base temperature (compare Tables 3 and 4). As we showed in §2.2, radiation pressure plays an important role when $`FF_{\mathrm{Edd}}`$, acting to lower the density and increase the hydrostatic expansion.
We showed in §3 that in a radiative atmosphere, the magnitude of the spin down depends on depth. Thus to observe a single, coherent frequency requires either the atmosphere to be rigidly rotating, or the burning region to be vertically thin (Miller 2000), so the differential rotation across it is small. Convection may enforce rigid rotation; however, one-dimensional calculations of bursts show that convection persists only for the first $`1\mathrm{s}`$ of the burst when the energy production is dominated by temperature sensitive and rapid helium burning reactions. In §4, we found that short wavelength baroclinic instabilities may act to enforce rigid rotation within the burning layers. Theoretical studies of detailed burst models and further investigations of angular momentum transport mechanisms are needed to determine whether a coherent pulse is possible once convection has ceased.
To compare the predicted spin down with observed values, we assume that some mechanism operates to enforce rigid rotation within the burning layers, but that the burning layers remain decoupled from the underlying ashes, as suggested by the results of §4. Figure 13 shows the spin evolution of the atmosphere assuming that the whole atmosphere rigidly rotates. For both convective and fully-radiative models, we ignore any differential rotation, and assume $`I\mathrm{\Omega }`$ is constant, where $`I`$ is the moment of inertia of the atmosphere (eq. ). As in Figure 12, the evolution during the burst is shown by the arrows, first the atmosphere spins down as flux increases, then spins back up as burning proceeds, increasing the mean molecular weight, and as the atmosphere cools. To compare with observations, we also indicate the observed frequency shifts (Table 1). For those bursts in which the oscillation frequency was seen only in the tail (Table 1), we plot $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ as a lower limit, since some cooling will have occurred before the oscillation is first seen.
Figure 13 shows that the change in spin frequency for pure He ignition models is less than that for mixed H/He ignition, because of the smaller hydrostatic expansion. For the largest frequency shifts observed, $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }8\times 10^3`$, pure He ignition models require a base temperature very close to the limiting value from radiation pressure (eq. ) and $`F>F_{\mathrm{Edd}}`$, so that the density at the base is decreased, enhancing the expansion. It is much easier to get $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ values this large with mixed H/He ignition, although as we describe below, all bursts with oscillations so far have the characteristics of pure He bursts.
In §3, we found that the time for heat to diffuse from the burning layers to the photosphere is $`1\mathrm{s}`$. This delay explains why the spin down of the atmosphere is not observed at the beginning of a burst; the expansion occurs before the signal is seen. We have also shown that the vertical propagation of heat through the atmosphere can affect the oscillation amplitude. A requirement for a large amplitude oscillation is that the time to transport heat from the burning layers (or the top of the convection zone) to the photosphere must be small compared to the shearing time across the atmosphere $`\mathrm{\Delta }\nu ^1`$. If so, the shearing atmosphere transmits the burning pattern from below, much like passing a flashlight under a piece of paper. If not, the oscillation will be smeared out, the atmosphere being heated uniformly by the quickly revolving burning layers.
Our results show that there are important differences in spin evolution depending on whether the helium ignites in a pure helium or mixed H/He environment. Because of the greater mean molecular weight of pure helium as opposed to a solar composition, the expansion of the atmosphere in pure He bursts is roughly half that of mixed H/He bursts, giving a smaller $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ (Figure 13). Also, the ratio of thermal time to shearing time is smaller for a pure helium atmosphere, and the energy production more localized because of the temperature sensitivity of He burning reactions, increasing the likelihood of observing coherent oscillations for pure He bursts.
### 5.2. Why Do Some Bursts Show Oscillations, But Not All?
A puzzle is that oscillations are not seen in bursts from all sources, or in all bursts from a particular source. So far, there has been only one study of the relation between burst oscillations and properties of the bursts. Muno et al. (2000) studied nine X-ray bursts from KS 1731-26 which occurred during observations with $`\mathrm{𝑅𝑋𝑇𝐸}`$ at various times between 1996 July and 1999 February. They found that bursts which show coherent oscillations are of short duration ($`10\mathrm{s}`$), show radius expansion, and have high peak flux. These characteristics are typical of helium rich bursts. Bursts with longer durations typical of H burning showed no oscillations. More studies are needed to see if this is a general result for other sources. Certainly all reported burst oscillations that we can find appear to be in bursts of short duration typical of pure He ignition. As far as we are aware, oscillations have not been detected in bursts of long duration ($`>50\mathrm{s}`$) typical of H burning.
Why should only pure-He ignition bursts show oscillations? The differences we find between pure He and mixed H/He bursts may provide an explanation. As we described in §3.4 (see Figure 8), we find that the thermal time is less and the shearing time is greater for pure helium bursts, making it more likely that a large amplitude oscillation from deep regions may propagate outwards.
### 5.3. Long Term Stability of the Oscillation Frequency
The frequency evolution reported for six bursts is well fit by the model $`\nu (t)=\nu _0\mathrm{\Delta }\nu e^{t/\tau }`$, where $`\nu _0`$ is the frequency in the burst tail, $`\mathrm{\Delta }\nu `$ the amount by which the oscillation frequency changes and $`\tau `$ the decay time. The parameters of these fits are given in Table 1. The observed decay times $`\tau 2`$$`7\mathrm{s}`$ are similar to the expected cooling time of the burning layers (eq. ). However, we found in §4 that there are mechanisms that might couple the burning layers to the star on a similar timescale. Can we distinguish between these two possibilities observationally?
Figure 12 shows that if the burning layers remain decoupled, the frequency observed in the burst tail will not be that of the neutron star spin. Because of the greater mean molecular weight of the ashes, the thickness of the cooling atmosphere in the tail is different to the initial thickness by $`1\mathrm{m}`$. This change in thickness depends on how complete the burning was, so that variations in the energetics and burning in the burst would translate into a scatter of $`1`$ part in $`10^4`$ in the final frequency. If the oscillation frequency could be shown to be more stable than this from one burst to the next, it would imply that the ashes and burning layers must recouple during the burst decay. This is complicated due to Doppler shifts from the orbital motion of the neutron star, which change the observed frequency at a level $`2\times 10^3(v_{\mathrm{orb}}/600\mathrm{km}\mathrm{s}^1`$) across the orbit. Strohmayer et al. (1998b) showed that two bursts from 4U 1636-54 showed a frequency difference consistent with orbital Doppler shifts, while two bursts from 4U 1728-34 separated by more than one year showed the same asymptotic frequency to $`1`$ part in $`10^4`$ (see Table 1). If the orbital parameters are known and the orbital Doppler shift accounted for, more measurements such as these would indicate whether recoupling occurs in the burst tail.
## 6. Conclusions
We have shown that the hydrostatic expansion, spin down, and later spin up during contraction of the neutron star atmosphere naturally explains the magnitude and sign of the observed changes in the nearly coherent oscillations observed in Type I X-ray bursts. Our results support the simple picture proposed by Strohmayer et al. (1997a) to explain the frequency evolution seen in Type I bursts, and thus the identification of the burst oscillation frequency with the neutron star spin.
The amplitude of the oscillation is set not only by the lateral extent of the asymmetry on the neutron star surface, but also by the vertical propagation of heat through the shearing atmosphere. In addition, we find that the spin behavior differs depending on whether the burst results from pure He ignition or He ignition in a H-rich environment. The spin down is smaller by roughly a factor of two for pure He ignition as opposed to mixed H/He ignition, because the pure He models expand outwards less due to the larger mean molecular weight. In addition, the time for heat propagation through the atmosphere is smaller for pure He models.
This might explain why oscillations have not been detected during bursts of long duration ($`30\mathrm{s}`$) typical of H burning (for example, the bursts from GS 1826-24 shown in Figures 2 and 5 of Kong et al. 2000). In particular, the recent study of Muno et al. (2000) of bursts from KS 1731-26 found oscillations only during those bursts with short durations $`10\mathrm{s}`$, typical of pure He ignition. However, the largest values of $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ that are observed (for 4U 1702-43 and 4U 1728-34) require base temperatures very close to the limiting value from radiation pressure for pure He bursts, and this may prove to be a problem. More observational studies of the phenomenology of burst oscillations and their relation to burst properties and source properties, in particular accretion rate, are needed. With further theoretical work, these studies could give us important clues as to how the nature of nuclear burning during bursts depends on accretion rate, and help to resolve the disagreement with one-dimensional models first found with EXOSAT (Bildsten 2000 and references therein).
We have not found a robust hydrodynamical mechanism that recouples the burning layers to the underlying star in the $`10\mathrm{s}`$ burst duration. However, the coherence of the observed burst oscillations suggests that little differential rotation occurs within the burning layer itself. As we note in §4, this might be accomplished through short vertical wavelength baroclinic instabilities, though more theoretical work is needed.
We have not discussed the complex question of the lateral spreading of the burning front during the burst rise, or the cause of asymmetry in the cooling tail of bursts. In the simplest picture of a spreading burning front (Strohmayer et al. 1997b; Strohmayer et al. 1998c) it is a mystery why oscillations are seen in the burst tail, when the whole surface has ignited. One possible explanation is that an asymmetry arises because the atmosphere takes a finite time to cool once ignited, and a finite time to spread around the star. Thus the matter on the opposite side of the star from the ignition point will be a little hotter since it ignited later. If the flux at each point is $`e^{t/\tau }`$ once ignited, and the burning front propagates around the whole star (starting at the equator) in a time $`t_{\mathrm{burn}}`$, it is straightforward to show that the peak-to-peak luminosity amplitude in the tail for an observer looking down on the ignition point is $`\mathrm{\Delta }L/L=2(1\mathrm{exp}(t_{\mathrm{burn}}/\tau ))/(1+\mathrm{exp}(t_{\mathrm{burn}}/2\tau ))^2`$, or, for $`t_{\mathrm{burn}}\tau `$, $`\mathrm{\Delta }L/Lt_{\mathrm{burn}}/2\tau `$. Thus for a spreading time $`t_{\mathrm{burn}}0.5\mathrm{s}`$ and $`\tau 10\mathrm{s}`$, the peak-to-peak amplitude in the tail is a few percent, a little less than the 5–20% amplitudes observed (Smith, Morgan & Bradt 1997; Zhang et al. 1998; Strohmayer et al. 1998a). Another possibility is that the oscillations come and go because of changes in the vertical transport of heat rather than changes in the size of the asymmetry on the stellar surface. For example, Figure 8 shows that as the atmosphere heats up during the burst rise, the ratio of thermal time to wrap-around time $`t_{\mathrm{therm}}\mathrm{\Delta }\nu `$ increases. As this ratio increases, it becomes more likely that the oscillation will be washed out. It may be that at the burst peak, the oscillations disappear because of this effect.
We have not addressed what mechanism might cause frequency doubling during bursts, as seems to be the case for at least one source (4U 1636-56; Miller 1999). There are two processes that we have discussed in this paper that may play a role. The first is the recoupling mechanism. If the burning layers recouple to the neutron star during the burst, they may do so in an $`m=2`$ pattern, leading to frequency doubling. Our linear instability analysis of the two-layer beta-plane model (§4.3) suggests that the baroclinic instability would have $`m20`$ when it goes unstable. However, we have not calculated the correct spherical eigenfunctions. The second possible process is the azimuthal wrapping of the the magnetic field. If an initially dipolar field is wound up, an $`m=2`$ pattern might emerge in the flow.
A test for the rotational modulation hypothesis might be the distribution of observed $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ values. For sources undergoing X-ray bursts of the same type (for example, similar accretion rates and same type of ignition) on similar mass neutron stars, we would expect the distribution of $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ values to be the same for each source, and independent of the spin frequency or number of hotspots on the surface. In fact, the distribution of $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ so far (Table 1) looks bimodal, though we should be wary of small number statistics. Two sources, 4U 1728 and 4U 1702, which have $`\nu _0350\mathrm{Hz}`$ show frequency changes $`\mathrm{\Delta }\nu /\nu 5`$$`8\times 10^3`$, whereas the other four objects which have $`\nu _0550\mathrm{Hz}`$ show frequency changes about half as large, $`\mathrm{\Delta }\nu /\nu 1`$$`4\times 10^3`$. If further observations show this bimodal distribution to be true, this might indicate that the expansion of the atmosphere is less for the $`\nu _0550\mathrm{Hz}`$ objects, for example, because of pure He rather then mixed H/He ignition. Alternatively, we note that the $`\mathrm{\Delta }\nu `$ values are very similar for all six objects, and this might point to another explanation for the observed frequency evolution. Care must be taken that the observed $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ represents the change in burst oscillation frequency from the burst rise to the burst tail. For example, if the oscillation is seen only during the burst tail, the observed $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ does not represent the full spin down of the atmosphere, as in this case some cooling will have occurred before the oscillation is seen.
We have not discussed other explanations for the frequency drifts. One possibility is that the burst oscillation is a non-radial oscillation (NRO) on the neutron star surface (Livio & Bath 1982; McDermott & Taam 1987; Bildsten & Cutler 1995; Strohmayer & Lee 1996). Non-radial oscillations of the frequency observed would be plausible on a rapidly rotating star, as the Coriolis force takes the typically much lower-frequency modes and makes them mugh higher (Bildsten, Ushomirsky and Cutler 1996). The observed frequency would be close to the spin frequency of the star or some multiple $`m\mathrm{\Omega }`$, and might change by a $`1`$$`2\mathrm{Hz}`$ during the burst as the atmosphere cooled. However, it is not clear why only a single mode would be selected, and why that frequency would be stable on timescales of a year or more, for which conditions in the surface layers of the star would most likely change significantly.
The physics of expansion-induced shear during shell flashes may prove important in other contexts. For example, mixing due to shear instabilities during a classical nova explosion could provide the enrichment of CNO elements necessary to achieve a fast nova (see Truran 1982 for a review of the enrichment problem). Previous authors have investigated shear mixing during accumulation of the fuel (see Livio & Truran 1990 for a review), but not during the burst itself when shear is induced by expansion of the atmosphere as in a Type I X-ray burst. A similar situation arises in AGB stars during He shell flashes, where mixing of carbon into the proton rich envelope is necessary to allow production of s-process elements (Iben 1991; Sackmann & Boothroyd 1991; Langer et al. 1999), and perhaps could be achieved by shear instabilities during shell flashes.
We thank Derek Fox, Daniel Holz, Yuri Levin, Mike Muno, Tod Strohmayer, Chris Thompson, Greg Ushomirsky, Marten van Kerkwijk and Ellen Zweibel for valuable discussions and comments. We thank the Astronomical Institute, “Anton Pannekoek” of the University of Amsterdam for hospitality and for supporting L. B.’s visit there as the CHEAF Visiting Professor. This research was supported by NASA via grants NAG 5-8658 and NAGW-4517 and by the National Science Foundation under Grants No. PHY94-07194 and AST97-31632. L. B. is a Cottrell Scholar of the Research Corporation.
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# Optimizing Conflicts in the Formation of Strategic Alliances
## 1 Introduction
Physical concepts might prove useful in describing collective social phenomena. Indeed models inspired by statistical physics are now appearing in scientific literature .
The process of aggregation among a set of actors seems to be a good candidate for a statistical physics like model. These actors might be countries which ally into international coalitions, companies which adopt common standards, parties that make alliances, individuals which form different interest groups, and so on.
Given a set of actors, there always exists an associated distribution of bilateral propensities towards either cooperation or conflict. The question then arises as to how to satisfy such opposing constraints simultaneously. In other words, what kind of alliances, if any, will optimize all actor bilateral trends to respectively conflict or cooperation.
It turns out that a similar problem does exist in spin glasses. For these systems, magnetic exchanges are distributed randomly between ferro and antiferromagnetic couplings. Indeed such an analogy has been used in the past in a few models.
The first one is by Axelrod and Bennett (hereafter denoted as AB) . They apply the physical concept of minimum energy to single out the stable coalitions within an configuration landscape. Later Galam demonstrated that the AB model is not fully consistent with its physical content. He then suggested a second model (hereafter denoted as G) based on both random bond and random site exchanges .
We propose here a qualitative extension of defining the coalition problem. While all previous approaches used implicitly a 2-side coalition dynamics, we allow multi-side coalitions. We indeed define the actual number of alliances as an internal parameter of the problem to be determined by simultaneously optimizing all bilateral propensities. While an Ising like variable may be appropriate to a huge number of problems both in physics and outside physics, there is no a priori reason to consider coalition dynamics restricted to a bimodal distribution.
We then apply our model to study the fragmentation of former Yugoslavia. Results are compared to the actual situation. We also revisit previous studies along this new scheme of not setting a priori the number of coalitions. As soon as the artificial bimodal constraint is relaxed, their results are no more consistent with reality.
Moreover, in physics models the exchange couplings and their distributions are known. In contrast, in the coalition problem the quantitative determination of bilateral propensities to either cooperation or conflict is a major challenge. We discuss this matter and argue that previous propensity calculations were too simplistic.
The rest of the paper is organized as follows. In the second part we review the AB and G models. The qualitative extension to multimodal coalitions is discussed in Section 3. On this basis, we revisit in Section 4 the cases of the Second World War and the Unix standard setting (studied in Ref. and respectively ). The case of the fragmentation of former Yugoslavia is analyzed in Section 5. Concluding remarks are contained in last Section.
## 2 Former models
### 2.1 The Axelrod-Bennett model (AB)
Axelrod and Bennett first attempted to explain the composition of coalitions using some characteristics of the involved actors (the elements of the system). The relative affinity between actor $`i`$ and actor $`j`$ is measured by a pairwise propensity $`p_{ij}`$, which is negative when they are in conflict and positive when they like to cooperate (with $`p_{ii}=0`$). Differences in actor sizes are expressed by assigning to each one a weight factor $`s_i`$ (a positive quantity). It may be a demographic, economic or military factor, or an aggregate parameter.
The creation of coalitions among the actors introduces a distance $`d_{ij}`$ between each pair $`(i,j)`$ of actors. It is 0 if $`i`$ and $`j`$ belong to the same coalition and 1 when they are in different coalitions. A given partition $`X`$ of the actors into coalitions is equivalent to knowing all $`d_{ij}`$.
A “frustration” of each actor $`i`$ is introduced to measure how much a configuration $`X`$ satisfies its propensities. It writes,
$$F_i=\underset{j=1}{\overset{n}{}}s_jp_{ij}d_{ij}(X),$$
(1)
where we considered a group of $`n`$ actors.
Adding all actor frustrations, respectively weighted by their sizes, results in an “energy” of the system,
$$E(X)=\frac{1}{2}\underset{i=1}{\overset{n}{}}s_iF_i,$$
(2)
which is,
$$E(X)=\underset{i>j}{\overset{n}{}}s_is_jp_{ij}d_{ij}(X).$$
(3)
It is then postulated that the actual configuration of the system is the one which minimizes the energy. The path followed by the system into the coalition landscape space from an initial configuration follows the direction of the greatest gradient of energy. Once a minimum is reached the system does not change.
The AB model has been applied to the study of the alliances of the Second World War . It was also used to study the standard setting coalitions formed by the companies which developed the Unix operating system .
Later Galam has demonstrated several inconsistencies of the AB model . Here we found additional setbacks (see Section 4).
### 2.2 Galam reformulation of the AB model
Galam has shown that in case of bimodal coalitions ($`A`$ and $`B`$), the AB model is totally equivalent to a finite spin glass at zero temperature.
Accordingly, configurations can be expressed by the spin variables $`\eta _i`$, where the spin is +1 if the actor $`i`$ belongs to coalition $`A`$ and it is -1 if the actor belongs to $`B`$. The distances can be rewritten as $`d_{ij}=\frac{1}{2}(1\eta _i\eta _j)`$. The energy becomes,
$$E(X)=E_0\frac{1}{2}\underset{i>j}{\overset{n}{}}J_{ij}\eta _i(X)\eta _j(X),$$
(4)
where $`J_{ij}=s_is_jp_{ij}`$ and
$$E_0=\frac{1}{2}\underset{i>j}{\overset{n}{}}J_{ij}$$
(5)
is a constant which depends only on the given sizes and propensities, but not on the configurations.
### 2.3 The Galam model (G)
Galam introduced a new model for the case of bimodal coalitions where the two alliances have clear characteristics that would allow to define for each actor a natural belonging $`ϵ_i`$ to each alliance. It is +1 if the actor $`i`$ should be in coalition $`A`$, according to its own characteristics compared with those of $`A`$. It is -1 for $`B`$ and 0 if there is no natural belonging to neither $`A`$ or $`B`$.
In the G model, the benefit $`J_{ij}`$ of cooperation or conflict between a pair of countries in either the same or opposite alliance is accounted for in addition to the local bond propensities $`G_{ij}`$. Both add together to a total propensity $`p_{ij}=G_{ij}+\eta _i\eta _jJ_{ij}`$.
In parallel, forces (like military or economic mechanisms) by which each coalition as a whole couples to the orientation of a given actor are expressed in terms of an external field $`\beta _i`$. It contributes to the overall conflict through the product $`\beta _ib_i\eta _i`$, where $`\beta _i=\pm 1`$ represents the direction of the force that acts on actor $`i`$ (towards $`A`$ or $`B`$), while $`b_i>0`$ is the amplitude of this force. The total energy of the system sums up to
$$E=\frac{1}{2}\underset{i>j}{\overset{n}{}}(G_{ij}+ϵ_iϵ_jJ_{ij})\eta _i\eta _j\underset{i}{\overset{n}{}}\beta _ib_i\eta _i.$$
(6)
This model has been used to explain the stability of alliances during the cold war, with opposing NATO and Warsaw pacts. The fragmentation of Eastern Europe which resulted from the Warsaw pact dissolution as well as the simultaneous western stability was recovered within the model. The European construction versus stability and the Chinese stability were also analyzed within the same model .
## 3 Beyond previous models
### 3.1 A new approach: going multimodal
In physics, the symmetry of the internal degrees of freedom of a given system is clearly determined from specific measurements on the studied material. In particular, most magnetic systems may be modeled using a spin variable. It can be, among others, a 2-state Ising variable, a q-state Potts model, a XY spin with planar symmetry, or a Heisenberg spin with continuous symmetry.
In the case of human systems, the a priori restriction of using an Ising variable has no ground. Some known coalitions exhibit more than 2 simultaneous alliances, though not much more. For example, the users and the producers of personal computers can be divided in at least three categories, with regard to the operating system they use (Windows, Mac OS, Linux $`\backslash `$ Unix); there were three groups that fought each other between 1941 and 1945 in the space of ex-Yugoslavia: the Chetnik guerilla groups, the Communist bands and the Nazis ; and in most democratic countries there are more than two independent political parties.
We extend previous approaches by allowing for multimodal coalitions. We assume the number $`q`$ of coalitions to be an internal degree of freedom to vary from 1 up to $`n`$. Before, $`q`$ was arbitrarily fixed to 2. The spin variable $`\eta _i`$ can thus take $`q`$ different values. Distances are expressed as $`d_{ij}=1\delta (\eta _i,\eta _j)`$ and the energy is
$$E(X)=2E_0\underset{i>j}{\overset{n}{}}J_{ij}\delta (\eta _i(X),\eta _j(X)),$$
(7)
where $`E_0`$ is given by equation (5). This corresponds to the $`q`$ state Potts model . The initial expression of the energy, equation (3), can still be used.
To demonstrate that the increase in the number of allowed coalitions may reduce the energy of the system, we will discuss the example of a system with $`n=3`$ actors. We assume $`J_{12}=J_{23}=J_{13}=1`$.
If we impose only one coalition, the energy is 0. If we impose two coalitions, the minimum energy is -2 and is three times degenerated (the system is frustrated). If we don’t impose a fixed number of coalitions, the system sets itself in a configuration with three different coalitions (all the actors are independent), with an energy -3, the lowest energy possible.
### 3.2 Gauge transformations
Let us consider a transformation of the distances as $`d_{ij}^{}=ad_{ij}+b`$. The energy becomes
$$E^{}(X)=\underset{i>j}{\overset{n}{}}J_{ij}d_{ij}^{}(X)=aE(X)+2bE_0.$$
(8)
If $`a>0`$ this transformation keeps unchanged the dynamics of the system.
An interesting transformation is given by $`d_{ij}^{}=2d_{ij}1`$. That is $`d_{ij}^{}=1`$ if $`i`$ and $`j`$ are in the same coalition and $`d_{ij}^{}=1`$ if not. For bimodal coalitions, we have $`d_{ij}^{}=\eta _i\eta _j`$. The energy expressed with this distance accounts equally for conflict and cooperation.
Another interesting transformation is $`d_{ij}^{\prime \prime }=d_{ij}1`$, which gives $`d_{ij}^{\prime \prime }=\delta (\eta _i,\eta _j)`$. In this case the energy is exactly the Potts energy,
$$E^{\prime \prime }(X)=\underset{i>j}{\overset{n}{}}J_{ij}d_{ij}^{\prime \prime }(X)=\underset{i>j}{\overset{n}{}}J_{ij}\delta (\eta _i(X),\eta _j(X)).$$
(9)
The minimization of this energy maximizes the cooperation.
Because all these distances yield the same dynamics of the system, it doesn’t need to be specified exactly if the actor aim is, psychologically, the minimization of conflicts, the maximization of cooperation, or both, and in which proportion. However, this prevents us from comparing the energies of different systems (including the same actors with some propensities changed) until this is established. In our study we will stick to the original AB form of the distance (which minimizes conflicts).
### 3.3 Neutrality
The introduction of spin variables allows us to consider the possibility of neutrality by letting the spins to have also the value $`\eta _i=0`$. In the case of bimodal coalitions, this is straightforward. In the case of $`n`$ coalitions, the energy (7) has to be rewritten as
$$E=2E_0\underset{i>j}{\overset{n}{}}J_{ij}\delta (\eta _i,\eta _j)[1\delta (\eta _i\eta _j,0)],$$
(10)
where the last factor accounts for the case when both $`i`$ and $`j`$ are neutral.
## 4 Application to real cases
### 4.1 The case of the Second World War
Axelrod and Bennett have applied their model of aggregation to explain the composition of the opposing alliances during the Second World War . The actors are 17 European countries involved in the war. Country sizes are measured with a national capabilities index which combines six components of demographic, industrial and military power. Propensities are computed from 1936 data. The criteria used are: ethnic conflicts, religion, border disagreements, type of government and history of war. They are combined with equal weights.
After numerical minimization of the energy, two different minima were found. The absolute one corresponds to respectively Britain, France, Soviet Union, Czechoslovakia, Yugoslavia, Greece and Denmark in one coalition versus Germany, Italy, Poland, Romania, Hungary, Portugal, Finland, Latvia, Lithuania and Estonia in the other one. After the criteria of Axelrod and Bennett (who measure the alignment by whether a country was invaded by another country, or had war declared against it), this corresponds to the historical reality, with the exception of Poland, which is on the wrong side, and of Portugal, which was neutral. There exists another local minimum with Soviet Union, Yugoslavia and Greece versus all the others.
While Axelrod and Bennett exhibit this result as a validation of their model, they hardly comment a crucial assumption they made. They supposed that the 17 countries can be partitioned in only 2 coalitions. Even if this was the historical reality, this artificially added constraint does not fit to the principles of the model. The final configuration should be determined only by the actors’ pairwise interactions, via the minimization of the energy. Nothing restricts a priori the existence of more than 2 coalitions. A partition in two coalitions should be indeed a result of the model and a proof of its predictiveness.
### 4.2 Revisiting the Second World War case
We have redone the AB computation using the same data set for sizes and propensities (which is available on Internet ).
We first confirmed the AB results once the number of coalitions is restricted to 2. We then introduced the possibility of spin 0 in the hope of capturing the neutrality of Portugal while keeping the other conditions unchanged. The results stay unchanged with no country neutral.
We then allowed an a priori unlimited number of coalitions (from 1 up to $`n=17`$). Exploring then the full energy landscape, we found that there exists one single minimum associated to 3 alliances. They are respectively Soviet Union and Greece versus Germany, Italy, Estonia and Latvia versus Britain, France, Czechoslovakia, Denmark, Yugoslavia, Poland, Romania, Hungary, Portugal, Finland and Lithuania. This minimum has an energy $`E=132.36`$. It is much less than the energies of the two AB configurations which are respectively -94.23 and -91.06.
We also tried to find a 2 coalition configuration imposing geographical constraints, but the result stays a three alliance configuration. This shows that, in fact, the real prediction of the AB model, using the same propensities, is rather far from the historical reality.
### 4.3 The Unix case
Axelrod et al. applied the AB model of aggregation to explain also the formation of standard setting alliances, like for the Unix operating system . The actors are 9 companies involved in the development of Unix (AT&T, Sun, Apollo, DEC, HP, Intergraph, SGI, IBM and Prime). Sizes are given by the firms’ share in the technical workstation market or by expert estimates (for AT&T).
If a firm $`i`$ belongs to alliance $`L`$, its utility (the satisfaction of the economic agent, given by its profit) is expressed by
$$U_i(L)=\underset{jL}{}s_j\left[\alpha \underset{jL^{}}{}s_j+(\alpha +\beta )\underset{jL^{\prime \prime }}{}s_j\right],$$
(11)
where $`s_j`$ is the size of firm $`j`$, and $`L^{}`$ and $`L^{\prime \prime }`$ form a partition of alliance $`L`$ into distant and close rivals of $`i`$ ($`L=L^{}L^{\prime \prime }`$ and $`L^{}L^{\prime \prime }=\mathrm{}`$). The parameters $`\alpha `$ and $`\beta `$ are positive. The identification of close and distant rivals is done from the degree of specialization of the firms in the production of Unix-based workstations.
We may rewrite the utility as
$$U_i(L)=\underset{jL}{}s_jp_{ij}=\underset{j=1}{\overset{n}{}}(1d_{ij})s_jp_{ij},$$
(12)
where the last expression gives the utility of $`i`$ without having to specify its alliance. This leads to the propensities $`p_{ij}=1\alpha `$, if $`i`$ and $`j`$ are distant rivals, and $`p_{ij}=1(\alpha +\beta )`$, if $`i`$ and $`j`$ are close rivals.
The energy (3) equals
$$E=2E_0\frac{1}{2}\underset{i=1}{\overset{n}{}}s_iU_i,$$
(13)
so a minimization of the energy yields a weighted maximization of the utilities.
For the choice of parameters $`\alpha =\beta =1`$ (and in all cases for which $`0.8\alpha 1.5`$, $`0.7\beta 1.5`$), Axelrod et al. found 2 minimums of the energy with the same value. One configuration is Sun, DEC and HP versus AT&T, Apollo, Intergraph, SGI, Prime and IBM. The other one is Sun, AT&T, Prime and IBM versus DEC, HP, Apollo, Intergraph and SGI. The latter configuration corresponds to Unix International vs. Open Software Foundation, and only IBM is incorrectly assigned.
Axelrod et al. advocate this result in favor of the effectiveness of their methodology. However, they have again imposed the artificial constraint of allowing only a maximum of 2 coalitions. They motivate this choice by the fact that “the positive externality created by standardization declines as the number of standards increases and thereby reduces the principal advantage of setting standards, which is a larger post-standardization market” (, p. 1498).
But this condition is in fact included in the formula of the utility $`U_i(L)`$, which grows linearly with the total size of coalition $`L`$. Therefore, the aggregation of firms into a small number of alliances should be only the result of the energy minimization. If the utility grows faster than linearly with the coalition size (which may be the case in reality), and the linear approximation is not good enough, then even if we impose a maximum of 2 coalitions the result would be wrong, because of the discrepancy between the real and the modeled utility.
### 4.4 Revisiting the Unix case
We checked the above model with no conditions imposed on the coalition number. For the case $`\alpha =\beta =1`$ there exist 150 different minimums of the energy, with an average of 6.57 coalitions per configuration. This is definitely a bad choice of parameters because all the propensities are negative or 0, so there is little incentive for aggregation.
The condition for the beginning of aggregation is $`\alpha <1`$. We sampled the interval $`0\alpha 0.9`$, $`0\beta 2`$, with a step of 0.1. For $`\alpha +\beta <1`$, all firms form a single block (all propensities are positive). Configurations with 2 coalitions appear in an interval given roughly by $`1\alpha <\beta <2(1\alpha )`$.
There are 14 types of configurations of 2 coalitions. We ranked their importance from the number of pairs ($`\alpha `$, $`\beta `$) for which they are realized, weighted with the size of their basin of attraction. The most preferred configurations are (i) Sun, AT&T and IBM versus the others (29.6 %), (ii) AT&T and Apollo versus the others (28.4 %), (iii) IBM, DEC and Apollo versus the others (12.5%), the rest having less than 10 % each.
The “best” configuration predicted by Axelrod et al. is not realized at all, while the other configuration they predicted has rank 4 (5.5%). The real empirical configuration Sun, AT&T and Prime versus the others has rank 9 (2.7%) and doesn’t have a basin of attraction greater than 6% for any combination of parameters.
Again, as for the Second World War, the prediction of the non constrained AB model, using the same propensities, doesn’t fit the empirical reality.
## 5 A new application: the case of Yugoslavia
We have studied a new problem, the fragmentation of the former Yugoslavia. The actors are its 8 administrative entities (provinces or republics): Serbia, Croatia, Bosnia, Slovenia, Macedonia, Vojvodina, Kosovo and Montenegro. Here a coalition means a federation of entities, or an independent state if it has only one member.
In the initial configuration all the actors are in the same coalition - the former Yugoslavian federation. The system is then left to evolve in the direction of the greatest gradient of energy down to the minimum which should correspond to the stable configuration.
To implement our model we needed first to evaluate all the propensities among the set of the 8 entities. The ethnic group diversity of the whole set is a major ingredient which we considered in evaluating those propensities. We also took into account differences in religion and language. Entity sizes are taken proportional to population sizes.
We used the 1981 census results , . We considered eight major ethnic groups: Serbs, Croats, Muslims, Slovenes, Macedonians, Montenegrins, Albanians and Hungarians. We neglected in our study the influence of other ethnic groups. They accounted for less than 1% each of the total population of Yugoslavia in 1981. We also neglected the 5.4% of the population which label themselves Yugoslavian. We assumed this later group to be perfectly tolerant of the others (zero propensity).
Due to lack of accurate data for all the entities, we considered all the Serbs, Macedonians and Montenegrins to be Orthodox Christians; all the Croats, Slovenes and Hungarian to be Catholic Christians; and all the Muslims and Albanians to be Muslim. For 1999, there exist differences up to 25% between the figures for nationality and religion, but in most cases they are much less . We classified the ethnic groups, with regard to their language, as Serbo-Croats (Serbs, Croats, Muslims and Montenegrins), other Slavs (Macedonians and Slovenes) and non Slavs (Albanians and Hungarians) .
The propensities between pairs of entities are computed as follows:
$$p_{ij}=\underset{k,l}{\overset{8}{}}q_{ik}q_{jl}w_{kl},$$
(14)
where $`q_{ik}`$ represents the percentage of ethnic group $`k`$ in entity $`i`$ and $`w_{kl}`$ represents the pairwise propensity between ethnic groups $`k`$ and $`l`$.
For $`k=l`$, $`w_{kk}=+1`$. For $`kl`$, the $`w_{kl}`$’s are computed as the sum of 2 terms. One stands for religion and the other for language: $`w_{kl}=\omega _{religion}(k,l)+\omega _{language}(k,l)`$.
We used the hypothesis that Christians are more tolerant of other Christian religion members than for Muslims. We also assumed that the Serbo-Croats are more tolerant of other Slavs than of non Slavs.
For religion, the factor $`\omega _{religion}`$ is positive and equals $`+\omega _1`$ for pairs of ethnic group with the same religion. It is negative and equals $`\omega _2`$ for pairs of Catholic Christian and Orthodox Christian ethnic groups. For pairs of Christian and Muslim groups, the factor is $`\omega _3`$.
For language, the factor $`\omega _{language}`$ is $`+\omega _4`$ in the case of two Serbo-Croatian speaking groups. For pairs of two different Slav groups, the factor is $`\omega _5`$. For pairs that include at least a non Slavic language, the factor is $`\omega _6`$. All the $`\omega _i`$’s are positive.
The graduation of tolerance previously described yields the following conditions: $`\omega _2<\omega _3`$ and $`\omega _5<\omega _6`$. We also have the condition $`\omega _1+\omega _4<1`$ to prevent a pairwise propensity between 2 different ethnic groups to be greater than the propensity within the same group, which is the reference factor for other propensities.
The parameters $`\omega _i`$ are unknown. We varied those parameters in the domain $`0<\omega _{1,4}<0.5`$; $`0<\omega _{3,6}<1`$; $`\omega _2<\omega _3`$; $`\omega _5<\omega _6`$, with a step of 0.05 and checked the results predicted by the minimization of the energy. There are 28 resulting configurations (from a total of 4140 possibilities).
All these configurations respect the geographical connectivity. We ranked them from the number of cases they are realized, weighted with the relative size of their basin of attraction in each case. Only three configurations appear in more than 10% of cases each and 11 other appear in more than 1% of cases each.
The main three configurations are: (i) A federation including Serbia, Croatia, Bosnia, Montenegro and Vojvodina, with the other entities being independent, is obtained in 42.4% of cases. (ii) In 12.1% of cases the result is Serbia, Montenegro and Vojvodina, and the others independent. (iii) In 10.9% of cases, only Serbia and Vojvodina stay together.
The above results yield the real configuration which resulted from the fragmentation of Yugoslavia for the second largest set of parameters. We considered Kosovo practically separated from Serbia because of its special situation after the 1999 war, under the control of NATO peacekeepers.
For this final configuration, in 31% of cases (for example for a choice of parameters $`\omega _i`$ like (0.2, 0.5, 0.6, 0.1, 0.3, 0.5)), the order of fragmentation is the real one: Croatia and Slovenia - June 1991, Macedonia - September 1991, Bosnia - April 1992, with the general exception of Kosovo. Our model predicts its splitting from Yugoslavia on the first or the third step, due to its mostly Albanian population. Within our model and the G model we might say that Kosovo was artificially kept inside Serbia using an external field. The field was then destroyed by the NATO bombing.
## 6 Conclusion
The introduction of analogies inspired from spin glass models into the study of social systems makes possible the prediction of the dynamics of macroscopic alliances formation within a given system of actors. While pairwise interactions are clearly instrumental in this approach, there is no scientific method to date to select the factors to be accounted for in their evaluation.
We have shown in particular that previous studies were not consistent with their own principles. Using the propensities and sizes computed from their methodology and then computing the resulting configurations, following strictly and solely the principle of the minimum energy, we have shown that the results don’t fit any more to reality. Nevertheless we have also shown that there exist some parameter ranges which do yield the real situation.
We have also emphasized that analogies with physics should not mean a straight mapping. A major difference here is to consider the number of alliances to be a free internal parameter of the system. In physics it is predetermined.
In conclusion our approach validates the feasibility of modeling strategic international behavior, but at the same time it demonstrates also the dangers of taking it too simplistically. Within our model, we have also reported configurations which don’t fit to reality. At this stage, this shows that more work is necessary to single out a feasible scheme for the evaluation of parameters.
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# Detection of a Far IR Excess with DIRBE at 60 and 100 Microns
## 1 INTRODUCTION
The extra-galactic background light (EBL), from optical to sub-millimeter wavelengths, records the energetics history of galaxy formation. This background is the cosmic relic of star formation, AGN, and black hole formation. The existence of such a background was discussed first in the optical and near IR (Partridge & Peebles (1967)) and then at other wavelengths (Low & Tucker (1968); Peebles (1969); Harwit (1970); Kaufman (1976); Dube et al. (1979)). However, measurement of an unresolved background is very difficult at most wavelengths because of numerous foregrounds which may be many times brighter. Direct measurement of the individual background sources by deep, high resolution imaging has only recently become possible, and only at selected wavelengths.
### 1.1 Current Knowledge of Extragalactic Background
The Infra-Red Astronomical Satellite (IRAS) in 1983 obtained the first full-sky census of far infrared (FIR) point sources. Some $`300,000`$ point sources, including $`20,000`$ galaxies, were detected in four bands from $`12100\mathrm{\mu m}`$ (see Beichman et al. (1988)). Optical follow-up indicated that typically 30% of the bolometric luminosity of these galaxies is radiated in the FIR (Soifer & Neugebauer (1991)), presumably as thermal radiation from dust heated by optical/UV radiation. In the case of ultra-luminous IR galaxies (ULIRGs) up to 95% of the bolometric luminosity is radiated in the FIR (e.g., Sanders & Mirabel 1996), suggesting that optically obscured galaxies might produce a substantial fraction of the extragalactic background light. Mid IR ($`1225\mathrm{\mu m}`$) emission from most ULIRGs is centrally concentrated, consistent with AGN activity (Soifer et al. (1999)).
More recently, Puget et al. (1996) used the COBE Far InfraRed Absolute Spectrophotometer (FIRAS) data to constrain the FIR/SMM background at longer wavelengths. They found that the integrated energy of the EBL in the $`200\mathrm{\mu m}2`$mm window is comparable to that emmited at optical/near IR wavelengths. This picture was confirmed by Guiderdoni et al. (1997) who concluded that the majority of high-$`z`$ star formation may be hidden by dust. The FIRAS measurement was greatly refined by Fixsen et al. (1998) who obtained a fit to the CIB of the form
$$I_\nu =(1.3\pm 0.4)\times 10^5(\nu /\nu _0)^{0.64\pm 0.12}B_\nu (18.5\pm 1.2\mathrm{K})$$
(1)
in the interval $`150<\nu <2400\mathrm{GHz}`$ ($`2000125\mathrm{\mu m}`$), where $`\nu _0=3000\mathrm{GHz}`$ and $`B_\nu (T)`$ is the Planck function. The integrated intensity observed in this frequency interval is $`14\pm 4\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$. This curve is the dotted line shown in Figure 1.
Recent efforts to resolve this background into discrete sources with the SCUBA detector on the James Clerk Maxwell Telescope have been very successful (Smail et al. (1997); Hughes et al. (1998); Barger et al. (1998); Eales et al. (1999)). By extending the counts to a limit of 0.5 mJy in cluster-lensed fields, enough counts are found to account for most of the expected $`850\mathrm{\mu m}`$ background of $`\nu I_\nu =5\pm 2\times 10^1\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ (Blain et al. 1999b ).
At shorter wavelengths, the background is not resolvable with current instruments, and zodiacal emission from interplanetary dust (IPD) hampers detection of the unresolved background. The zodiacal emission peaks at $`25\mathrm{\mu m}`$ and dominates any expected signal from the CIB in most DIRBE channels. Nevertheless, the background has been measured at $`140\mathrm{\mu m}`$ and $`240\mathrm{\mu m}`$ by Schlegel, Finkbeiner, & Davis (1998). A more thorough analysis by the DIRBE team (Hauser et al. (1998)) improved those results, and also provided upper limits in the other 8 DIRBE channels (see Figure 1). The FIRBACK survey at $`175\mathrm{\mu m}`$ (Puget et al. (1999)) using the *Infrared Space Observatory* (ISO) detected point sources to a flux limit of $`200\mathrm{mJy}`$, yielding an integrated flux significantly higher than a simple extrapolation of IRAS counts would have predicted, but still much lower than the observed background. Puget et al. (1999) claim that a plausible extrapolation of the counts down to $`10\mathrm{mJy}`$ would account for the entire background at $`175\mathrm{\mu m}`$.
At wavelengths $`\lambda 3.5\mathrm{\mu m}`$, emission from Galactic stars dominates that from zodiacal dust. Using ground-based $`2.2\mathrm{\mu m}`$ counts to remove the stellar foreground, and adopting a value of $`7.4\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`2.2\mathrm{\mu m}`$, Dwek & Arendt (1998) measured the background at $`3.5\mathrm{\mu m}`$ to be $`\nu I_\nu 10\pm 3\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$. Gorjian, Wright, & Chary (1999) make different assumptions about the stellar foregrounds, yet arrive at a similar result: $`16.2\pm 6.4\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`2.2\mathrm{\mu m}`$ and $`9.3\pm 3\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`3.5\mathrm{\mu m}`$.
At optical wavelengths the Hubble Space Telescope can resolve a large fraction of the extragalactic background. By integrating the light of the resolved galaxies in the HDF, Pozzetti et al. (1998) find that the brightness of the extragalactic sky is $`2.1_{0.3}^{+0.4}\times 10^{20}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1\mathrm{Hz}^1\mathrm{sr}^1`$ in I-band ($`0.8\mathrm{\mu m}`$), which is equivalent to $`\nu I_\nu =8\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$. In the other Hubble broadband filters, this background obeys $`\nu I_\nu \nu ^1`$ for 2000 - 8000Å. The HDF counts in each filter appear to flatten at the faint end, indicating a possible convergence, although other authors find hints of an optical/UV background as much as twice this large (Bernstein (1997)). Comparing the latest FIR background with the optical/UV background has supported the view that $`2/3`$ of the starlight in high-$`z`$ galaxies is reprocessed by dust into FIR radiation, a much higher percentage than in local galaxies. Where, then, are the sources of the FIR background?
It is well established that the ULIRGs are typically interacting galaxies (see Sanders & Mirabel 1996), and Williams et al. (1996) find that most high-$`z`$ objects in the HDF are interacting. What is not very well established is the dominant energy source in these objects, especially at the high-luminosity end. Certainly, many of these objects are powered by starbursts; only the very brightest objects so far observed are primarily powered by AGN (Genzel et al. (1998); Lutz et al. (1998)). But this scenario may change at high redshift. Measurements of the X-ray background provide some constraint on AGN activity, unless the very brightest AGN are so obscured that even $`10\mathrm{keV}`$ photons cannot escape.
### 1.2 Interpretation
The above observational evidence allows for a coherent interpretation. Interacting ULIRGs at high-$`z`$ undergo violent episodes of dusty star formation, with AGN as a minor contributor to the energy. Roughly 2/3 of this energy is reprocessed to the FIR. Accounting for present day metallicity of galaxies and the IGM enrichment, the energy from stars allows for an integrated extragalactic background light (EBL) of $`I_{\mathrm{bol}}=50/(1+z_f)\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ (see §5.1), which is roughly consistent with current measurements if the formation redshift $`z_f2`$.
The IR excess measurements presented in this paper, if they are interpreted as an extragalactic background, would push this paradigm to the limit. They would indicate an integrated EBL flux in the far IR of $`40\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ with a hot spectrum, possibly suggesting that AGN dominate the energy input in early galaxies, at least at short wavelengths. This interpretation might also imply the presence of up to $`0.15\%`$ of all baryons in black holes, and possibly violate constraints on the X-ray background. Furthermore, recent measurements of high energy gamma rays place limits on the opacity of the IGM, which is primarily due to pair production on CIB photons. The interpretation of the $`60100\mathrm{\mu m}`$ excess as CIB raises serious problems; however, we have been unable to identify any alternative source in the Galaxy, local bubble, or Solar system that can account for the emission.
### 1.3 Organization of Paper
We review previous foreground models in §2, and present our two procedures for extracting the IR excess signal in §3. Our assessment of various systematic errors is given in §4. In §5 we discuss the energy crisis resulting if this emission is interpreted as an extragalactic background. In addition, we discuss the constraint on the density of CIB photons implied by the observation of TeV photons from nearby AGN. These difficulties, as well as general conclusions and prospects for the future appear in §6.
## 2 FOREGROUNDS
### 2.1 Zodiacal Emission
The main difficulty in measuring the CIB at $`60240\mathrm{\mu m}`$ is contamination from zodiacal light, which is thermal emission from the interplanetary dust (IPD). This emission is brightest in the DIRBE 12 and $`25\mathrm{\mu m}`$ bands, falling approximately as a blackbody at longer wavelengths. The emission at the ecliptic poles is $`\nu I_\nu 260`$ and $`50\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ in the DIRBE 60 and $`100\mathrm{\mu m}`$ bands respectively (Fig. 1). Emission in the ecliptic plane is $`34`$ times brighter. This emission must be carefully removed in order to measure the much fainter extragalactic background.
The IPD cloud is difficult to model, especially at low ecliptic latitude. At high latitude, one looks through dust in the neighborhood of the Earth, but at low latitude, the situation is more complex: the dust density and temperature vary significantly along a line of sight, several distinct dust rings are seen, and a density resonance in the Earth’s orbit is observed. These factors make it nearly impossible to model the zodiacal light at low latitudes from $`5100\mathrm{\mu m}`$. Furthermore, temperature and density variations appear even in the near-Earth dust, beyond the expected variation due to changes in the Earth-Sun distance through the year.
#### 2.1.1 SFD98 Zodiacal Emission Model
A simple approach to the problem of separation of zodiacal light from other emission is presented in Schlegel, Finkbeiner, & Davis (1998; hereafter SFD98). They use the $`25\mathrm{\mu m}`$ map as a spatial template of the zodiacal light. At high ecliptic latitude, most of the dust is less than 0.4 AU from Earth and has a fairly uniform temperature. Near the ecliptic plane, one sees dust of varying temperature out to several AU. Therefore, this $`25\mathrm{\mu m}`$ template does not extrapolate to longer wavelengths in a linear way, and another level of detail must be added.
To good approximation, the error made by a linear model is a function only of ecliptic latitude. Therefore, one must modulate the $`25\mathrm{\mu m}`$ map by some reasonable function of ecliptic latitude. Rather than choose standard basis functions such as a series of ($`m=0`$) spherical harmonics, SFD98 instead employs a function determined by the dust itself - simply the $`25\mathrm{\mu m}`$ flux binned in ecliptic latitude. This modulation is adequate with only one term in the expansion, and comprises the so-called “quadratic zodiacal light model” (see SFD98, equation 3).
The success of such a method depends on a few assumptions. One is that the Galactic dust cirrus is negligible at $`25\mathrm{\mu m}`$. SFD98 showed that, in any place where the $`100\mathrm{\mu m}`$ flux is less than $`100\mathrm{MJy}\mathrm{sr}^1`$, the $`25\mathrm{\mu m}`$ flux due to cirrus causes at most a 1 percent error in the final map. Another assumption is that there is no CIB in the $`25\mathrm{\mu m}`$ map. This approach will not be sensitive to any CIB component that has the same spectral shape as the zodiacal emission. Current models of the CIB predict no substantial flux at $`25\mathrm{\mu m}`$, and certainly less than the flux at $`140\mathrm{\mu m}`$ (in $`\nu I_\nu `$ units). If the CIB flux at $`25\mathrm{\mu m}`$ were as high as it is at $`140\mathrm{\mu m}`$ ($`25\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$), the offset introduced in the 60 and $`100\mathrm{\mu m}`$ cirrus maps would be negligible ($`3\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`60\mathrm{\mu m}`$ and $`1\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`100\mathrm{\mu m}`$).
This “quadratic model” was used by SFD98 to separate cirrus (Galactic ISM) emission from CIB and zodiacal light. At 140 and $`240\mathrm{\mu m}`$, the zodiacal emission is faint enough that its latitude dependence provides enough information to separate it from the CIB; at shorter wavelengths, more sophisticated methods are required.
#### 2.1.2 Goddard Zodiacal Emission Model
The model invented by the Goddard team (Kelsall et al. (1998)) is an ambitious attempt to parameterize the full spatial-temporal dependence of the IPD emission. It contains six components, a smooth cloud, circumsolar ring, a density enhancement following the Earth and in resonant lock, and also three dust bands near the ecliptic at 3 AU. Certainly, the complexity of the dust cloud justifies the 46 model parameters, and a sophisticated model is required to establish the isotropy of the background. However, even though the parameterization is physically well motivated, a simple subtraction of the model from the data leaves significant unexplained residuals at $`\lambda 100\mathrm{\mu m}`$. These residuals may result from detector gain drifts, or small variations in dust density and temperature. For this reason, a more robust approach is required.
The two methods given in §3 each construct a dimensionless parameter that is robust (at first order) with respect to gain drifts and dust density and temperature variation. This parameter is evaluated for each of the 40 weeks of the DIRBE mission, using the DIRBE weekly skymaps with Galactic emission subtracted. Before discussing these methods in detail, we next consider the Galactic (ISM) emission removal from each weekly skymap.
### 2.2 Galactic ISM dust emission
Emission from Galactic dust in the ISM (cirrus) is typically comparable in brightness to the IR excess at $`100\mathrm{\mu m}`$, and about 1/4 as bright at $`60\mathrm{\mu m}`$ (Fig. 1). As pointed out in the “zodiacal light” fits in SFD98, it is easier to separate the Galactic cirrus from the sum of the other emission components than it is to isolate the CIB. SFD98 describe FIR maps that have a “quadratic zodiacal light model” (actually zodiacal light plus CIB) removed, which we call the “cirrus-correlated maps.” These maps are constructed by removing a quadratic zodiacal emission model from the DIRBE annual average maps so as to maximize the correlation of the cirrus map with the Leiden-Dwingeloo H I survey at high galactic latitude (see SDF98 for details). Because only 3 degrees of freedom are used to fit the HI data to the full-sky DIRBE maps, the structure of the Galactic ISM is derived directly from the DIRBE maps, not from H I data. These annual-average DIRBE maps are subtracted from each weekly average DIRBE map before the processing described in §3 is performed. The structure of the Galactic cirrus and extragalactic point sources are therefore almost perfectly subtracted, with remaining residuals attributed to detector gain drift and shifted effective pixel centers from week to week.
The uncertainty in our final IR excess measurements resulting from the cirrus emission is discussed in §4.2. There may also be substantial emission from dust in the warm ionized medium (WIM), as traced by H$`\alpha `$ and pulsar dispersion measures. However, the WIM is well correlated with the CNM on spatial scales of interest, and this contribution should not result in a significant error. Section 4.2 gives a complete discussion of systematic uncertainty associated with the WIM dust.
The next section addresses the brightest foreground, the zodiacal light.
## 3 ZODIACAL EMISSION REMOVAL: TWO METHODS
### 3.1 Method I
Our first method makes use of the north-south variation of the zodiacal emission observed by DIRBE as a function of time. The Interplanetary Dust (IPD) cloud is inclined $`2^{}`$ with respect to the ecliptic, resulting in a north-south asymmetry with a one year period (see Figs. 2 and 3). At first, it appears that this temporal variation is an undesirable complication in an already complex problem. However, this temporal variation allows us to probe the dust cloud in the $`z`$ direction (normal to the dust plane) and separate zodiacal emission from the other components. Because the temperature and density of the IPD change throughout the year, it is useful to consider a dimensionless parameter that corresponds to our $`z`$ position in the IPD cloud, a number for which uncertainties in the overall IPD density and temperature cancel out. Therefore, it is convenient to define a dimensionless ratio,
$$R_b\frac{N_bS_b}{N_b+S_b},$$
(2)
where $`N_b`$ ($`S_b`$) is the total DIRBE flux at the north (south) ecliptic pole in band $`b`$. In bands where the zodiacal light overwhelms emission from cirrus, CIB, or other contamination, this quantity is related only to the $`z`$ position of Earth in the dust cloud - not to density or temperature. Because variations in dust temperature, density, and detector gain drifts cancel out in $`R_b`$, it is far more robust than any absolute measurements. In fact, $`R_b`$ is nearly independent of IPD model. We are assuming that the $`z`$-dependence of IPD emission per volume has the same functional form for each waveband, which should be true for dust near the Earth. This requirement is not strictly satisfied by the IPD, but is justified in section 4.
Because it is exceedingly difficult to model the IPD cloud along all lines of sight at all times, we restrict ourselves to analysis of patches within $`5^{}`$ of the ecliptic poles. These regions contain data for every week of the DIRBE mission, and effects dependent on solar elongation angle cancel out to first order. We have prepared a mask for pixels in this region, excluding those pixels not present in symmetric combinations for each week. Details are given in §4.5.
Our resulting measurements of the IR excess at 60 and $`100\mathrm{\mu m}`$ depend upon two main assumptions: (1) that the amplitude of the annual variation in $`R`$ would be nearly constant in all wavebands in the absence of the excess, and (2) that our $`100\mathrm{\mu m}`$ cirrus map (SFD98) is correctly zero-pointed. Both of these assumptions have been explored thoroughly, and are investigated in detail in §4.
Plots of $`R_{12}`$ and $`R_{25}`$ are shown in Figures 2b and 3b. The line is a simple sinusoidal model - the best fit for
$$R_{model}=A\mathrm{sin}(\nu \nu _0)+C$$
(3)
where $`A`$ is the amplitude, $`\nu `$ is the true anomaly (i.e. the angle, measured at the Sun, between perihelion and the Earth $``$ true longitude minus $`102.8^{}`$), $`\nu _0`$ is a phase angle, and $`C`$ is a small constant. The true anomaly, $`\nu `$, is used instead of mean heliocentric longitude because it corresponds more directly to Earth’s $`z`$ position in the dust cloud. The deviations between the model and data are too small to see in Figures 2 and 3, and therefore are plotted in Figure 4. The RMS dispersions of $`R_{12}`$ and $`R_{25}`$ relative to this simple model are 1% of the amplitude of the sine wave - and 0.1% of the total emission at the poles. Furthermore, there is a strong correlation between the residuals at 12 and $`25\mathrm{\mu m}`$ \- suggesting that the residuals are physical. This is a powerful test of the relative photometric stability of the DIRBE 12 and $`25\mathrm{\mu m}`$ detectors, and allows us to proceed to the next step.
#### 3.1.1 Scale-Height of the Dust
A common *Ansatz* for $`z`$-dependence of density in a disk is $`\mathrm{exp}(|z|/h)`$ where $`z`$ is the height above the dust plane and $`h`$ is some scale height. Such a model has a density cusp at $`z=0`$, but because we travel only $`0.03`$AU in the $`z`$ direction, the effect of the cusp is negligible. Another reasonable guess might be a gaussian
$$\rho =\rho _0\mathrm{exp}(z^2/2\sigma _z^2).$$
(4)
This distribution has the convenient property that $`\rho (r,z)`$ is smooth near $`z=0`$, consistent with our sinusoidal fit. A third model is the Goddard widened fan model (Kelsall et al. (1998), Eq. ), with a vertical profile
$$f(\zeta )=\mathrm{exp}(\beta g^\gamma )$$
(5)
where $`\zeta |Z_c/R_c|`$,
$$g=\{\begin{array}{cc}\zeta ^2/2\mu \hfill & \hfill \mathrm{for}\zeta <\mu \\ \zeta \mu /2\hfill & \hfill \mathrm{for}\zeta >\mu \end{array}$$
(6)
and $`\beta =4.14\pm 0.067`$, $`\gamma =0.942\pm 0.025`$, and $`\mu =0.189\pm 0.014`$ are the best-fit values of the parameters.
Results presented in this paper are derived using the gaussian model and are indistinguishable from the other models. All we require is that the density distribution be reasonably smooth near $`z=0`$ and not extend too high above the plane.
How much north-south variation is expected? Symmetry considerations suggest that the dust plane is approximately aligned with the invariable plane of the solar system, which is perpendicular to the total angular momentum vector of the solar system. The north pole of this plane is at $`\alpha _{2000}=273.85`$, $`\delta _{2000}=66.99`$ or $`\lambda =17.8`$, $`\beta =88.42`$ in ecliptic coordinates. (We denote the heliocentric mean ecliptic longitude of Earth with $`L`$ hereafter in this paper, and reserve the usual symbol $`\lambda `$ for wavelength.) This difference between the ecliptic plane and invariable plane is critical for our method, but the actual inclination angle is nearly degenerate with the dust cloud height in our fits. Our results are the same whether the $`1.58^{}`$ angle to the invariable plane, or the $`2.03^{}`$ inclination of the Kelsall et al. (1998) smooth cloud is used.
The perihelion (minimum Earth-Sun distance) is at $`L102.8^{}`$. The ascending node of the putative dust plane is at $`107.8^{}`$. Because they are near each other, the time of maximal excursion from the dust midplane occurs when the Sun-Earth distance is approximately 1 AU. Therefore, the extreme values of $`z`$ are $`\pm 0.0276`$ AU. Within this range (according to the fit in Figure 3) resides $`10.4`$% of the dust. Simple algebra then gives a scale height of $`\sigma _z=0.20\pm 0.01`$ AU for the gaussian. This is very similar to the FWHM of the Kelsall model. Of course, the density profile need not be a gaussian in $`z`$; if the emissivity of the IPD can be written as $`I_b(z,r)n(z,r)`$ then we have the model-independent constraint that
$$_{low}^{hi}I_b(z,r)n(z,r)𝑑z=0.104_{\mathrm{}}^{\mathrm{}}I_b(z,r)n(z,r)𝑑z$$
(7)
We only use this to justify our assumption that the dust along a line of sight is near the Earth and of uniform temperature. None of the conclusions in this paper depend on the functional form of the $`z`$-dependence.
#### 3.1.2 Model Fit
Now that we have established a reasonable model for the $`z`$ dependence of dust emission, let us formalize it a bit. Suppressing the $`b`$ subscript for convenience, let us define N (S) to be the total emission observed near the north (south) pole of the invariable plane:
$$NZ_N+B_N$$
(8)
$$SZ_S+B_S$$
(9)
where $`Z_N`$ ($`Z_S`$) is the zodiacal light in the north (south) and $`B_N`$ ($`B_S`$) is the time-independent background in the north (south), including cirrus, CIB, Reynold’s Layer, and any other unknown backgrounds such as halo dust. Let us also define
$$ZZ_N+Z_S;BB_N+B_S$$
(10)
where $`Z`$ is the total column emission through the IPD plane, while $`B`$ is *twice* the average background.
For simplicity, we assume that the emission per volume is constant near the ecliptic plane, such that $`R`$ depends only on $`z=r\mathrm{sin}\theta \mathrm{sin}i`$, with $`r`$ and $`\theta `$ suitably defined for an elliptical orbit. The position of Earth is parameterized by the true anomaly, $`\nu `$, given in the almanac as (Astronomical Almanac (1991), p. E4):
$$\nu =M+(2ee^3/4)\mathrm{sin}M+(5e^2/4)\mathrm{sin}2M+(13e^3/12)\mathrm{sin}3M+O(e^4)$$
(11)
where $`e`$ is the eccentricity of Earth’s orbit ($`e0.0167`$), $`M=L\stackrel{~}{\omega }`$, $`L`$ is the mean longitude and $`\stackrel{~}{\omega }102.8^{}`$ is the mean longitude of perihelion. The true Earth-Sun distance is then
$$r=(1e^2)/(1+e\mathrm{cos}\nu )$$
(12)
where $`r`$ is in AU. We let $`\mathrm{\Omega }`$ denote the nominal longitude of the ascending node of the dust plane (approximately $`L=77^{}`$) and define $`\theta `$ as
$$\theta =\nu +\stackrel{~}{\omega }\mathrm{\Omega }$$
(13)
The height of Earth above the dust plane is given by
$$z=r\mathrm{sin}\theta \mathrm{sin}i$$
(14)
where $`i`$ is the inclination of the dust plane ($`i1.58^{}`$). The zodiacal emission at the ecliptic poles is given by
$$Z_N=\frac{Z}{2}\left(1+Ar\mathrm{sin}\theta \right)$$
(15)
$$Z_S=\frac{Z}{2}\left(1Ar\mathrm{sin}\theta \right)$$
(16)
and the ratio $`R`$ by
$$R=A^{}r\mathrm{sin}\theta +C.$$
(17)
with the definitions
$$A^{}A\left(\frac{Z}{Z+B}\right);C\frac{B_NB_S}{Z+B}$$
(18)
Here the constant factor $`\mathrm{sin}i`$ is absorbed in the $`A`$ coefficient to emphasize that $`A`$ and $`\mathrm{sin}i`$ are degenerate parameters in this model. The results of this paper do not depend on the value of $`\mathrm{sin}i`$ in detail, only that it be small and constant. Physically, $`A`$ is the amplitude of the annual variation in $`R`$ due to the zodiacal emission and $`A^{}`$ is the observed amplitude.
The total background may now be expressed in terms of the observables, $`N`$ and $`S`$:
$$B=(N+S)\left(1\frac{A^{}}{A}\right),$$
(19)
and the difference,
$$B_NB_S=C(N+S),$$
(20)
which leads immediately to expressions for the north and south:
$$B_N=\frac{B+C(N+S)}{2}$$
(21)
$$B_S=\frac{BC(N+S)}{2}$$
(22)
#### 3.1.3 Results
Results of the fits in the DIRBE 5-$`240\mathrm{\mu m}`$ bands are presented in Table 1. The best-fit model parameters $`A^{}`$, $`\nu _0`$, and $`C`$ are given for each waveband, as well as the derived parameters $`Z`$, $`B_N`$, and $`B_S`$. The value of $`\nu _0`$ is determined from the 12 and $`25\mathrm{\mu m}`$ bands and adopted for the others. Values of $`A(\lambda )`$ must be assumed in order to derive $`B`$ from $`A^{}`$, and because the $`25\mathrm{\mu m}`$ band is dominated by zodiacal emission, $`AA^{}`$. As can be seen from the large errors for the 140 and $`240\mathrm{\mu m}`$ bands, this method breaks down at long wavelengths where the zodiacal emission is weaker and the S/N of the detectors is much lower. Although the actual CIB level at 12 and $`25\mathrm{\mu m}`$ is unknown, it must be non-negative. The assumption of a significant CIB in these wavebands would only increase the amount deduced for the longer wavelengths. Therefore, we assume the CIB at $`25\mathrm{\mu m}`$ to be zero to make a conservative assessment of the emission at $`60240\mathrm{\mu m}`$. The systematic errors introduced by this assumption are determined by computing the change in the measured excess if the $`25\mathrm{\mu m}`$ CIB has the same $`\nu I_\nu `$ as $`140\mathrm{\mu m}`$. The large excess at $`12\mathrm{\mu m}`$ is very uncertain, because of the larger model dependence at short wavelengths. The excess at $`5\mathrm{\mu m}`$ is sufficiently model dependent that it should not be taken seriously.
Values for $`B_N`$ and $`B_S`$ are shown separately in Table 1 to demonstrate that the $`NS`$ difference due to cirrus has been adequately removed from the weekly maps. We take the $`A`$ dependence calculated for the Goddard widened-fan model, but consider other models in §4 to estimate systematic errors. This slight model dependence will propagate into our final systematic errors, but for now we retain the assumption that $`A`$ is constant in every DIRBE band.
#### 3.1.4 Model Refinements
The residual seen in Figure 4 suggests a residual with a 1/3 yr period at both 12 and $`25\mathrm{\mu m}`$. We have added a few parameters to the model to allow for two dust disks of different thicknesses, inclinations, and ascending nodes (similar to the “circumsolar ring” in the Goddard model), and this improvement removes the 1/3 yr period signal. Another improvement was made to account for possible emission from interstellar dust grains focussed by the Sun into a cone in the downstream direction of the Sun’s motion relative to the local ISM. Although inclusion of these effects improves the $`\chi ^2`$ of the fit, the change is not substantial, and no significant variation in the derived IR excess results. They are therefore not considered further in this paper.
### 3.2 Method II
In this section, we present a method to disentangle the CIB from zodiacal emission by making use of the ecliptic latitude dependence of the latter. We define a dimensionless statistic, $`\mathrm{\Xi }`$ that can be measured independently for each of the 41 weeks of the DIRBE mission. This quantity is designed to be insensitive to the zodiacal emission model parameters. At the same time, this statistic provides a sensitive measure of any isotropic background. Although this statistic involves less robust assumptions about the IPD than method I, we test its dependence on the parameters of the Kelsall et al. (1998) model. The $`\mathrm{\Xi }`$ statistic allows for a wider range of statistical tests and addresses questions about isotropy. Furthermore, the assumption of a small CIB value at $`25\mathrm{\mu m}`$ is not required.
In SFD98, the CIB was measured at $`140\mathrm{\mu m}`$ and $`240\mathrm{\mu m}`$ by fitting and removing a $`\mathrm{csc}|\beta |`$ slab component to the annual-average DIRBE skymaps at high Galactic latitude. Of course, the annual-average maps combine data from many solar elongation angles averaged over 41 weeks (not 1 yr) and contain the resulting artifacts. Moreover, there are theoretical and observational reasons to suspect that the IPD is a “modified fan” and not a slab at all (see Kelsall et al. (1998), §4.2 for discussion). In the limit where the zodiacal emission is small compared to the CIB, this method gives reasonable results, but fails badly at $`\lambda 100\mathrm{\mu m}`$. The basic idea can be used successfully, however.
By working only with solar elongation $`90^{}`$ data in each weekly map, the problem is conceptually simpler. In Figure 5 we show the volume emissivity density contours of the $`e=90^{}`$ plane, i.e., the plane containing Earth and perpendicular to the Earth-Sun line. The axes are labeled with Cartesian ecliptic coordinates, in AU. The Earth is in the middle (at a time of year when $`x=1`$, $`y=z=0`$), with lines of sight to the ecliptic poles labeled NP and SP, and the “forward” direction of Earth’s orbit to the right. Four other lines are drawn, all of which are at latitude $`|\beta |=45^{}`$. They are labeled NF for “North-Forward,” NB for “North-Backward” and so on. A convenient dimensionless ratio to define is:
$$\mathrm{\Xi }(\beta )\frac{NF+NB+SF+SB}{2(NP+SP)}$$
(23)
This quantity is almost completely insensitive to vertical position in the dust cloud, to a small inclination of the dust cloud with respect to the ecliptic, dust temperature, or to nearly any other parameter in the Goddard (Kelsall et al. 1998) model. In fact, when $`\mathrm{\Xi }(\beta )`$ is computed for the Goddard model (the most realistic model to date), its annual variation is negligible ($`<0.001`$) for the latitudes of interest ($`35^{}<\beta <50^{}`$).
If the zodiacal emission were approximately a slab, one would expect the functional dependence $`\mathrm{\Xi }(\beta )=\mathrm{csc}(\beta )`$. It is notable, however, that the $`\mathrm{\Xi }`$ “measured” from the Goddard model at long wavelengths is significantly greater than $`\mathrm{csc}|\beta |`$ in the range of interest ($`35^{}<\beta <50^{}`$). This is because the Rayleigh-Jeans emissivity density contours shown in Figure 5 (*solid contours*) follow the volume density contours closely, and the “fan” nature of the model causes an upward curvature of the contours, resulting in $`\mathrm{\Xi }(\beta )>\mathrm{csc}|\beta |`$.
At short wavelengths, the situation is reversed. The emissivity is so temperature sensitive that distance from the Sun is the overriding concern. In this case, the emissivity contours curve downward, yielding $`\mathrm{\Xi }(\beta )<\mathrm{csc}|\beta |`$ (Figure 5, *dashed contours*).
These deviations from the geometry of a uniform slab are fine points, and do not affect the measurements at $`60`$ and $`100\mathrm{\mu m}`$, as we shall see, but one must account for them carefully in order to measure the background at $`3.512\mathrm{\mu m}`$, which is beyond the scope of the current paper.
For simplicity, we now return to the approximation that $`\mathrm{\Xi }_0=\mathrm{csc}|\beta |`$. Now let us consider the effect of an isotropic background, $`B`$, on observed values of $`\mathrm{\Xi }`$.
$$\mathrm{\Xi }=\frac{B+Z\mathrm{csc}|\beta |}{B+Z}=1+\frac{\mathrm{csc}|\beta |Z}{B+Z}$$
(24)
where $`Z`$ is $`NP+SP`$. Solving for $`B`$ gives
$$B=\left(1\frac{\mathrm{\Xi }1}{\mathrm{csc}|\beta |1}\right)I$$
(25)
where $`I`$ is the observed flux at the poles ($`Z+B`$) and $`B`$ is the twice the value of the CIB, as it was in method I. Figure 6 contains plots of $`\mathrm{\Xi }`$ at $`35^{}`$ for 60 and $`100\mathrm{\mu m}`$ applied to the cirrus-subtracted weekly DIRBE maps. . Unsurprisingly, the fit residuals are correlated from week to week. To account for this correlation, we estimate there are no more than 4 independent measurements within the 41 weeks, and thus the standard deviation of the mean of $`\mathrm{\Xi }`$ is reduced from the rms scatter by only $`\sqrt{4}`$.
Results for method II are shown in Table 2. For each latitude bin and each band $`b`$, $`\mathrm{\Xi }_{K_b}`$ is calculated from the Kelsall IPD model and compared with the measured $`\mathrm{\Xi }_b`$. The background, $`B/2`$, is then determined from eq. (25). A weighted average gives $`28.0\pm 1.9\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`60\mathrm{\mu m}`$ and $`21.9\pm 2.7\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`100\mathrm{\mu m}`$.
This agreement between the two methods is encouraging, and suggests that the observed excess is not coming from within the solar system, at least it does not vary spatially or temporally in the way the IPD is expected to.
## 4 SYSTEMATIC ERRORS
The DIRBE detector noise is small enough that measurement errors in the determination of the FIR excess are modest, but several distinct systematic errors contribute to the uncertainty in the final result. The IPD model dependence is discussed in §4.1 and possible emission correlated with the WIM is considered in §4.2. Section 4.3, more generally rules out emission from a dust slab aligned with the Galactic plane, and possible emission from the Galactic halo is addressed in §4.4. Final uncertainty estimates are presented in §4.5.
### 4.1 IPD Model Dependence
A fully self-consistent model of the IPD emission has not yet been found, probably because of the large number of dust components whose temperature and density may vary spatially and temporally. The Goddard model (Kelsall et al. (1998)) is certainly the most complete, but it still must resort to fudge factors to explain the emissivity function of the IPD, and it assumes that the emissivities of the spatially separate components are identical. With currently available data, it is not economical to introduce still more parameters in order to solve this problem, so no CIB measurement that depends in detail on the zodiacal emission model can be trusted. However, the Kelsall et al. model is used as a reference model in the following.
Neither method described in the previous section is strongly influenced by this choice of IPD model. Our analyses rely only on data at high ecliptic latitudes, where many of the zodiacal components, such as the dust bands at $`\beta <15^{}`$ can be safely ignored. Furthermore, the dimensionless parameters, $`R`$ and $`\mathrm{\Xi }`$, can be predicted in a nearly model-independent way and readily compared with the data. The advantage of this approach is that the results depend on relative measurements made on short timescales, and are almost independent of the choice of IPD model. What little dependence there is enters the two methods in different ways.
#### 4.1.1 Method I
For most reasonable models of zodiacal emission, the expected amplitude $`A`$ in method I should be a weak function of wavelength, not a constant as we assumed. Figure 7a, shows the function $`A(\lambda )`$ predicted by the Goddard (Kelsall et al. (1998)) model (*dot-dashed line*) compared to the observed $`A^{}(\lambda )`$ values. The solid line is the Goddard model evaluated at the poles of the dust plane instead of the ecliptic poles, for $`T_0=286`$K. In this model, the dust temperature depends on distance from the Sun, and $`T_0`$ is the temperature at 1 AU. The other lines are for $`T_0`$ higher and lower by a factor of two. Of course this is an absurdly large range of temperatures, but we use it to illustrate that for single emission components, the amplitude $`A(\lambda )`$ is nearly constant on the Rayleigh-Jeans side of the IPD spectrum. The dependence of $`A`$ on other model parameters is similarly weak.
Two-component models of IPD can also be considered. If there are two distinct grain sizes with different $`T_0`$, it is still impossible to fit the observations. In the worst-case scenario, the dust grains are segregated into two populations, one in a dust layer near the ecliptic plane, and the other away from it. The spectrum of $`A`$ is then simply the ratio of the spectra of the two components - both of which have identical Rayleigh-Jeans tails.
In order to fit the observed $`A^{}(\lambda )`$ values with no background, one requires two populations of dust grains with different emissivity laws. An example of this is shown in Figure 7b - but this is an extreme case of grain properties and geometry constructed to give the desired result. In this model, the IPD away from the midplane is very similar to the Kelsall model: $`T_0=286\mathrm{K}`$, blackbody dust. However, the slice of dust within $`|z|<0.03`$AU (the extent of Earth’s annual excursion from the midplane) has $`T_0=236\mathrm{K}`$ and $`\nu ^{0.5}`$ emissivity. These parameters seem unlikely, but they greatly reduce the excess $`60\mathrm{\mu m}`$ flux, and affect the $`100\mathrm{\mu m}`$ excess slightly.
We considered a class of two-component models with different emissivity power laws, different temperatures, with one in a much thinner disk than the other. Exploring such models in detail, we found that it to be impossible to reduce the $`60100\mathrm{\mu m}`$ background substantially without either producing excessive emission in the $`12\mathrm{\mu m}`$ zodiacal signal or going to unreasonable parts of the model parameter space.
We therefore proceed by assuming that an extrasolar, isotropic excess at $`60\mathrm{\mu m}`$ is physically more acceptable than a contrived IPD emission model. The reasonable range of IPD models do not alter the derived $`60`$ and $`100\mathrm{\mu m}`$ excesses by more than $`5\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ (95% confidence).
#### 4.1.2 Method II
Because method II is sensitive to the shape of the $`e=90^{}`$ emissivity contours (as seen in Fig. 5), it is immune to a two-component model of the type contrived above, as long as the components are layered with azimuthal symmetry. However, method II is in general more dependent upon the zodiacal light model used. This dependence is calculated by using the Kelsall et al. model as a reference IPD model. Table 3 displays the error in the IR excess measurements introduced by $`1\sigma `$ and $`3\sigma `$ changes in the main 6 parameters of the Kelsall model. Only the parameter $`\alpha `$, (the density $`\rho R^\alpha `$) affects method II results significantly. This is not surprising, because $`\alpha `$ affects the shape of the IPD cloud more than most of the other parameters. Uncertainties in the Kelsall model propagate into our results only at the level of $`3\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`60\mathrm{\mu m}`$ and $`1\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`100\mathrm{\mu m}`$ (95% confidence). If the Kelsall model provides a reasonable description of the shape of the IPD cloud, then the model uncertainties cannot be much bigger than this. If, however, the Kelsall model is missing a substantial component that changes the emissivity contours’ shape in the $`e=90^{}`$ plane, then the uncertainty could be much larger.
### 4.2 Dust Emission from the ISM
In this section we assess the uncertainty in the IR excess described above due to a zero-point error in the Galactic ISM (or cirrus) emission maps. According to SFD98, Table 2, the largest formal uncertainty in the cirrus emission is less than $`1\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$. However, there is a larger systematic uncertainty resulting from neglect of dust in the warm ionized medium (WIM). The issue is not whether there is dust in the WIM; all ISM dust emission, described by the SFD98 cirrus maps, is subtracted from the DIRBE weekly maps before $`R`$ and $`\mathrm{\Xi }`$ values are fit. However, any zero point error in the SFD maps will propagate directly into the measurement of the CIB. In fact, the flux measured at the poles contains no cirrus and are essentially a full-sky zodiacal light fit, evaluated at the poles (method I) and at other latitudes (method II). This is important for isotropy considerations, because the results we obtain effectively use data from the entire high-latitude sky via the SFD98 fit to the Leiden-Dwingeloo H I survey.
#### 4.2.1 Correlation with H-alpha
One tracer of the WIM that can be used to constrain the cirrus zero-point is H$`\alpha `$ emission. A recent paper by Lagache et al. (1999) addresses this question. Using the high-quality H$`\alpha `$ data of Reynolds et al. (WHAM; Haffner et al. (1998)) to trace the WIM, they claim to find significant WIM emission at $`1001000\mathrm{\mu m}`$, using 2% of the sky at high galactic and ecliptic latitude, and H$`\alpha `$ emission between 0.2 and 2 R. We have repeated their analysis with the same data in the same regions of the sky and find no WIM dust emission uncorrelated with H I emission. This high-latitude analysis provides no indication that our zero-point is incorrect.
Although we are interested in zero point problems in the diffuse cirrus at high galactic latitude, where the SFD98 zero point was determined, we must resort to a different analysis that makes use of H$`\alpha `$ data closer to the galactic plane where the signal is strong. This will provide a “worst case” result. Heiles et al. have shown that a simultaneous fit of H I and H$`\alpha `$ yields only a modest shift in zero point (Heiles, Haffner, & Reynolds (1999)). In several regions near Eridanus, they perform a fit of the form
$$I_\lambda =A+BN(\mathrm{H}\mathrm{I})+CN(\mathrm{H}\alpha )$$
(26)
and find that $`A`$ varies by approximately $`6\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ RMS at $`100\mathrm{\mu m}`$ in the various regions, with a central value of $`2\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$. These fits are performed with the SFD98 $`100\mathrm{\mu m}`$ map which already has a model of zodiacal light+CIB removed, i.e., it is zeroed to H I. The fact that these offsets are so near zero indicates that the dust correlated with H$`\alpha `$ might possibly explain as much as $`6\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ (the extreme case found by Heiles et al.) of our $`30\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ background measurement at $`100\mathrm{\mu m}`$, but is not likely to significantly alter the result.
Although the H$`\alpha `$ result is encouraging, it suffers from a few weaknesses. The sky coverage is small (2% in Lagache et al. $`10\%`$ in Heiles et al.), and future analyses using the entire WHAM data set may provide more concrete answers. Also, the H$`\alpha `$ emission is not proportional to N(H II) but rather to $`n_p^2𝑑l`$ and is also weakly dependent on temperature. This means that a correlation of H II density and dust/gas ratio could contrive to produce $`A0`$ even though the derived zero point is incorrect. Therefore, an alternative method is desirable as a confidence check.
#### 4.2.2 Pulsar Dispersion Measures
The pulsar dispersion measure is a straightforward determination of the column density of electrons $`N(e^{})`$ along the line of sight to a pulsar. It does not depend on the temperature or density of the ionized gas, but does rely on the pulsar being far enough away to give a fair assessment of the Galactic H II. For this work, we have made use of the pulsar catalogue assembled by J. Taylor at Princeton (Taylor et al. (1993) <sup>1</sup><sup>1</sup>1The latest version is available at http://pulsar.princeton.edu/ftp/pub/catalog/.).
The catalogue contains 707 pulsars, 146 of which are at $`|b|>20^{}`$ and have distance quality codes of “a” or “b”. Of these, 108 are in the region covered by the Leiden-Dwingeloo survey, allowing a comparison of $`N(e^{})`$ and H I with $`100\mathrm{\mu m}`$ flux. A further requirement that the pulsars be out of the plane ($`|z|>400`$pc) reduces the list to 46. As one can see in Figure 8, the zero point of the H/dust regression changes only modestly when the pulsar data are included, even though the slope changes by roughly 1/3. This indicates that H I and H II are correlated, but perhaps are no more correlated than any other $`\mathrm{csc}|b|`$ mechanism. Because the scatter is no tighter with H II included, one might conclude that there is little dust associated with H II. On the other hand, the poor assumption that each pulsar is behind all Galactic dust may add noise, cancelling out the improvement.
Curiously, the $`y`$-intercept in 8(a) is not zero, even though it is forced to be (by construction) over the average high-latitude sky (see SFD98). This nonzero intercept may indicate variation in the gas/dust ratio - even at high $`|b|`$, or may also reflect large-scale gradients in the SFD98 dust model. Whether the SFD98 temperature correction is used or not, we find that the largest zero-offset induced by such changes is at most $`0.1\mathrm{MJy}\mathrm{sr}^1`$ ($`3\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$) at $`100\mathrm{\mu m}`$. If we assume that perhaps half of the dust emission emanates from above the 400pc pulsar cut, and double this effect, it is still negligible. Furthermore, if we consider the change in zero point due to the use of a temperature correction as a systematic error that propagates directly to the IR excess, we still find a systematic uncertainty of only $`6\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$. Because this is an uncertainty similar to that obtained from the H$`\alpha `$ analysis of Heiles above, we adopt $`6\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ as the (95% conf.) uncertainty associated WIM-correlated dust emission, and add this to our systematic error budget (Table 4) at 60 and $`100\mathrm{\mu m}`$.
### 4.3 Ruling out a dust slab
There is still a chance that a diffuse layer of dust more than 400pc above the disk of the Galaxy could be responsible for the emission. Such a layer, if behind most of the pulsars, could be either uncorrelated with H$`\alpha `$ emission, or could be associated with H II so diffuse that H$`\alpha `$ emission is effectively suppressed. This sort of a foreground would be indistinguishable from the IR excess using the methods described in this paper, but would reveal itself by a dependence on Galactic latitude.
Because the DIRBE data do not extend over a full year, and because of drifts in the zodiacal light intensity and detector gain with time, the annual average maps contain unphysical gradients that may confuse a direct fit of a $`\mathrm{csc}|b|`$ component. In order to test for the presence of such a component, we again introduce a dimensionless parameter.
For each weekly DIRBE map, we construct a dimensionless parameter $`\chi `$ from the flux in four patches on the sky, always placed at solar elongation $`|e|=90^{}`$ and $`|\beta |=75^{}`$.
The mean flux values in the four patches of sky are designated $`I_{NF}`$ for “North-Forward”, $`I_{NB}`$ for “North-Backward”, and likewise $`I_{SF}`$, and $`I_{SB}`$ for the south, just as in the definition of $`\mathrm{\Xi }`$. In this case, “North forward” refers to a direction on the sky $`(\lambda ,\beta )=(L_{true}+90,+75)`$ where $`L_{true}`$ is the true heliocentric longitude of Earth. The $`L_{true}+90`$ direction does not correspond precisely to the direction of Earth’s velocity around the Sun because of eccentricity; rather it is at solar elongation $`90^{}`$. These values are computed for each week of data, except when any of the patches is at low Galactic latitude ($`|b|<10^{}`$)
The four lines of sight used each week form an “X” in space. Two useful combinations are $`I_A=I_{NF}+I_{SB}`$ and $`I_B=I_{NB}+I_{SF}`$. We then define the dimensionless ratio
$$\chi =\frac{I_AI_B}{I_A+I_B}$$
(27)
in which gain drifts and dust variations nearly cancel out. Because $`I_A`$ and $`I_B`$ are measured at solar elongation $`90^{}`$, $`\chi `$ would be zero in the absence of ISM emission, if the IPD were aligned with the ecliptic plane. Misalignment with the ecliptic plane will produce a periodic signal in $`\chi `$, revealing the inclination of the dust plane if the ascending node is known. Likewise, a uniform slab of emitting material will contribute another $`\mathrm{csc}|b|`$ periodic term. Unfortunately, the Galactic plane and dust midplane have similar ascending nodes so that their signal in this statistic is nearly degenerate, making a simultaneous fit impossible. Fortunately, an error of $`1^{}`$ in the inclination of the dust plane would result in an error of 1.5 and $`0.3\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at the poles in the 60 and $`100\mathrm{\mu m}`$ channels respectively. An isotropic background does not contribute to $`\chi `$.
As an example of the power of this technique, we display in Figure 9 $`\chi _{60}`$ and $`\chi _{100}`$ for the 60 and $`100\mathrm{\mu m}`$ channels after removal of the H I correlated component. For these fits, we used the inclination angle $`i=2.03^{}`$ obtained by the DIRBE team (Kelsall et al. (1998)). Before H I removal (not shown in figure) the $`\mathrm{csc}|b|`$ term is strong, with values of 40 and $`100\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at the Galactic poles for 60 and $`100\mathrm{\mu m}`$ respectively. The flux is twice that at the ecliptic poles at $`|b|=29.8^{}`$. After removal of the best fit H I coefficient, only a small signal is left, at the level of 2.5 and $`1.2\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$. Even though we have not explicitly removed a $`\mathrm{csc}|b|`$ component, the H I fit appears to have done so. The same procedure is difficult in the other wavebands. The noisy 140 and $`240\mathrm{\mu m}`$ channels do not give a meaningful measurement of $`\chi `$, and at 12 and $`25\mathrm{\mu m}`$ the technique is much more sensitive to the inclination angle of the dust plane.
This small residual signal at 60 and $`100\mathrm{\mu m}`$ may reflect a small error in the H I subtraction, or may actually be an emission component not correlated with H I. Whether real or not, this component is very small compared to the IR excess measured in this paper, and in fact is small compared to the WIM error derived above. Therefore, we consider the error due to a Galactic slab to be already included in the WIM error adopted above.
### 4.4 Galactic Halo Dust
In this section we consider another potential component of Galactic emission which would not have appeared in this other tests: dust emission in the halo of our galaxy. We have ruled out dust correlated with H$`\alpha `$, dust correlated with H II within 400pc of the Galactic plane, and dust in a roughly $`\mathrm{csc}|b|`$ distribution. However, a very diffuse component of dust mixed into the Galactic halo – associated with ionized H or perhaps no gas at all – has not been strictly ruled out. The marginal detection of reddening along lines of sight passing near spiral galaxies (Zaritsky (1994)) implies that there may be dust at $`r>60`$kpc, although it is another question whether such dust is warm enough to mimic the measured excess. Alton et al. (1998) have observed cold (15K) dust 2 kpc off the disk of NGC 891 with SCUBA, but such dust would not give the observed signal at 60 and $`100\mathrm{\mu m}`$ (see also Howk 1999).
Mechanisms that transport dust from the Galactic disk into the halo by radiation pressure have been proposed (e.g. Ferrara et al. (1991)), but contain only crude approximations to the Galactic magnetic field. Consideration of a more realistic magnetic field would increase the diffusion time and increase the probability of dust destruction. The dust which makes it into the halo, if any, might be expected to be hotter than the disk dust, as long as the radiation field is similar. The grains should be smaller, and thus attain a higher equilibrium temperature; and the radiation pressure mechanism preferentially transports grains with a high optical/UV cross section to mass ratio - and these may also tend to be hotter. However, to explain the observed spectrum, they must be quite hot. Taking the warm $`\nu ^{2.6}`$ component from Finkbeiner et al. (1999), a temperature of $`28\mathrm{K}`$ must be maintained to explain the observed spectrum - in contrast to an average interstellar dust temperature of $`18\mathrm{K}`$ for single-component $`\nu ^2`$ models or $`16\mathrm{K}`$ for the warm component of the Finkbeiner et al. two-component model.
Such a halo must also be nearly isotropic. The halo proposed by Ferrara et al. (Fig. 2 in Ferrara et al. (1991)) is unlikely, as it would have been detected by the method described in §4.3. In fact, a uniform spherical dust halo bright enough to explain the IR excess would appear isotropic enough for $`r_{halo}20`$kpc that it would not be noticed in any conceivable analysis of the DIRBE data. The theoretical complications of such a halo are many, but are perhaps no more distasteful than the conflict this IR background causes with TeV gamma ray observation, as we discuss in §5.2. However, the idea appears to be ruled out by the observational non-detection of a $`60\mathrm{\mu m}`$ halo in M31 at the level of $`0.1\mathrm{MJy}\mathrm{sr}^1`$ ($`5\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$). In fact, Rice et al. (1988) found no evidence that the IRAS $`60\mathrm{\mu m}`$ emission of optically large galaxies extends into the halos. It is difficult to see why our Galaxy should be any different.
To summarize this section, if there is indeed a component of dust in our Galaxy that accounts for the IR excess, it must be uncorrelated with H$`\alpha `$, more than 400pc off the plane, and differ substantially from $`\mathrm{csc}|b|`$. This forces us to seriously consider the possibility that the observed excess is mostly extragalactic in origin.
### 4.5 Final Results
In this section, we average the results from methods I and II and state the final results for the DIRBE 60 and $`100\mathrm{\mu m}`$ channels, including systematic uncertainties. The casual reader should skip directly to §4.5.3
#### 4.5.1 Method I
In method I, several computational choices introduce systematic uncertainties. One choice is whether to take a mean, median, or mid-average of the pixels at the ecliptic poles in each week. Some outlier rejection is necessary because of imperfect point source removal. However, because there is a gradient across the polar cap (due to solar elongation dependence) a mean is more stable. All results in this paper were derived with a mid-average, in which the highest and lowest 10% of the values are discarded, and the remaining values averaged. This proves to be more robust than either a mean or a median. Use of a mean increases our results by $`5\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`60\mathrm{\mu m}`$ and twice that at $`100\mathrm{\mu m}`$. However, a straight mean also results in an unacceptable $`NS`$ asymmetry (formally $`25\sigma `$ at $`25\mathrm{\mu m}`$) which is not present when a mid-average is used.
Another choice is the size of the polar region. Ecliptic latitude cuts of $`|\beta |=85^{}`$ and $`|\beta |=87^{}`$ were tested. The $`85^{}`$ cut provides adequate signal, but excludes the brightest parts of the LMC. The variation introduced by this choice is negligible at 60 and $`100\mathrm{\mu m}`$, but large at 140 and $`240\mathrm{\mu m}`$. This is because the uncertainty in those channels is dominated by measurement noise, and a $`|\beta |=87^{}`$ leaves much less signal to work with.
The use of the ecliptic pole is motivated by the obvious symmetries. Another potential choice is the apparent pole of the dust plane, which is inclined $`2^{}`$ with respect to the ecliptic. This choice can modify results by $`5\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$.
A number of pixels are discarded for lack of coverage. Some weeks of DIRBE data, particularly “science” weeks 20 and 35 (numbered 0 to 40) contain very little sky coverage and are rejected. Other weeks with more than 10% bad pixels at the poles are also discarded so as not to mask out too many pixels in the remaining weeks.
The masks are generated independently for each waveband, which may raise questions about comparing the derived $`A^{}`$ and $`C`$ parameters for each waveband. However, the amplitude $`A`$ determined at $`12,25\mathrm{\mu m}`$ is quite insensitive to the mask used. One can take the union of multiple masks and apply that “master” mask to each waveband and still obtain the same results, though with lower signal-to-noise. We have no indication that such a procedure is necessary, so we simply apply to each waveband its own mask. Table 5 shows the fraction of pixels lost due to the mask in each waveband, and the number of weeks of good data used.
Even the method of pixel mask symmetrization calls for some judgment. For each bad pixel $`(\lambda ,\beta )`$, it is necessary to mask $`(\lambda ,\beta )`$, or else the solar elongation gradient will be aliased into the time domain, and appear as signal in $`R`$. In fact, masking $`(\lambda ,\beta )`$ and $`(\lambda ,\beta )`$ corrects this problem to a higher order. This 4-fold mask symmetrization was followed for all Method I results. Failure to symmetrize the mask in this way adds $`2\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`60\mathrm{\mu m}`$ and $`6\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`100\mathrm{\mu m}`$.
In light of these systematic errors, we assign a systematic uncertainty of 7 (11)$`\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at 60 (100)$`\mathrm{\mu m}`$ to method I (95% conf.), in addition to the formal statistical errors. At 140 and $`240\mathrm{\mu m}`$, measurement noise dominates, so no such systematic errors are added.
The intent of method I is to make reliable IR excess measurements at $`60240\mathrm{\mu m}`$, but the less reliable results obtained at shorter wavelengths are interesting as well.
The flux measured at $`12\mathrm{\mu m}`$ ($`2\%`$ of zodiacal light) is almost certainly an artifact of the greater model dependence at short wavelengths. The $`5\mathrm{\mu m}`$ emission shown in Table 1 is relatively stronger ($`10\%`$ of zodiacal light) but is subject to extreme model dependence. It is interesting that this flux level of $`23\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ is similar to the $`5\mathrm{\mu m}`$ excess found by Dwek & Arendt (1998). They declined to call this a CIB because of anisotropy. This agreement between the two numbers may only be a curious coincidence.
#### 4.5.2 Method II
In method II, there are fewer choices to make. Each line of sight uses pixels in a $`5^{}`$ diameter patch, as this is the largest patch that will be statistically independent from week to week and latitude bin to latitude bin. This method is sufficiently robust that the choice of the ecliptic as the symmetry plane is unproblematic: the Kelsall model gives less than 1 part in 1000 variation in $`\mathrm{\Xi }`$ due to this choice. In fact, the only significant systematic error this method has in common with method I is the mask-symmetry error, which we take to be the same as in method I. We therefore adopt a systematic uncertainty of 2 (6)$`\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at 60 (100)$`\mathrm{\mu m}`$ for method II.
#### 4.5.3 Combined Results
The two methods are complementary: the first uses time-variability at the ecliptic poles for the analysis, and the second uses the spatial morphology of the data in each week. Method II is superior in that it samples a larger fraction of the data set and achieves a much higher S/N, but the use of different regions of the sky in each week results in systematic errors that are difficult to understand in detail. Therefore, there is no *a priori* reason to prefer one method over the other. We combine the results with a weighted average, using the formal random errors and systematic errors that apply to each method. Such errors average down when results of the two methods are combined (see Table 4. The remaining systematic errors that apply to *both* methods, such as the uncertainty in the ISM emission, are added after this averaging.
In Table 6 the results for method I, method II, and the average are stated. These numbers represent the FIR excess in 60 and $`100\mathrm{\mu m}`$ DIRBE filters, which are calibrated assuming a flat spectrum in $`\nu I_\nu `$. Color corrections for plausible FIR background spectra do not change these results by more than $`1020`$% (see DIRBE Exp. Supp. (1995)).
## 5 DISCUSSION
Because we are unable to find any emission component within the Galaxy to explain the observed IR excess, we tentatively interpret the emission as an isotropic extragalactic background (CIB). In this section we discuss the consequences of such an interpretation.
### 5.1 The EBL Energy Crisis
The integrated energy in the CIB for $`60\mathrm{\mu m}<\lambda <1`$mm is $`40\pm 12\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ according to our measurements and those of Fixsen et al. (1998) - roughly twice the FIR energy content inferred from Hauser et al. (1998), and much higher than that predicted by Malkan & Stecker (1998). Where is this energy coming from? A comprehensive analysis by Bond, Carr, & Hogan (1986) reviews several possible sources for the FIR background, including primeval galaxies, pregalactic stars, black hole accretion, and decaying particles. Nuclear fusion within stars at epoch $`z`$ contributes a radiation energy density $`\mathrm{\Omega }_R`$ today (in units of the critical density) of
$$\mathrm{\Omega }_R0.007\frac{\mathrm{\Delta }Z}{1+z}\mathrm{\Omega }_B$$
(28)
where $`\mathrm{\Omega }_B`$ is the baryon density of the Universe today (e.g. $`\mathrm{\Omega }_Bh^2=0.019\pm 0.001`$; Burles & Tytler (1998)), $`\mathrm{\Delta }Z`$ is the mean metallicity of all baryons in the Universe, and $``$ is the fraction of the emitted radiation which is reprocessed into the FIR.
Black hole accretion is another source of energy for the CIB and yields a radiation density
$$\mathrm{\Omega }_R=\left(\frac{ϵ}{1ϵ}\right)(1+z)^1\mathrm{\Omega }_{BH,acc}$$
(29)
where $`ϵ0.1`$ is the efficiency of rest mass to energy conversion and $`\mathrm{\Omega }_{BH,acc}`$ is the mass density of accreting black holes.
With the definition of the critical density, $`\rho _{crit}=1.9\times 10^{29}h^2\mathrm{g}\mathrm{cm}^3`$ and $`h=H_0/`$(100 km s<sup>-1</sup> Mpc$`{}_{}{}^{1})`$, the detected background of $`40\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ becomes $`\mathrm{\Omega }_{CIB}=9.8\times 10^7h^2`$. If this radiation was generated within stars and mostly reprocessed into the FIR at $`z2`$, then $`\mathrm{\Delta }Z.02`$, (regardless of $`h`$) averaged over the entire Universe. This seems rather large, and would indicate that most star formation must be well hidden from view.
Another possibility is generation of the CIB by accretion onto black holes, again at $`z2`$, with most of the radiation emerging in the FIR. If the black holes were formed by the accretion of baryonic matter ($`\mathrm{\Omega }0.019h^2`$), with an efficiency $`ϵ=0.1`$, then the cosmological density of black holes today must be $`\mathrm{\Omega }_{BH}2.9\times 10^5h^2`$ and the fraction of all baryons that must have fallen into massive black holes (assuming no accretion of dark matter) is approximately 0.15%. While this number is large, it is smaller than the black hole mass fraction of 0.5% in galactic bulges (Magorrian et al. (1998)).
It is also possible that energy extraction from black hole accretion is more efficient than assumed above. Recent work by Gammie (1999) indicates that for rotating black holes with a magnetic field, accretion efficiencies of $`ϵ0.5`$ are possible. This reduces the requirement to 0.015% of baryon mass in black holes, consistent with the Magorrian et al. measurement and a disk/bulge mass ratio of $`30`$.
Recent work (Almaini et al. (1999)) shows that AGN observed by Beppo-Sax can explain the 30 keV X-ray background, and concludes that most of the energy generation takes place in obscured AGN. However, those AGN have only been shown to contribute 10-20% of the extragalactic background at 240 and $`850\mathrm{\mu m}`$, and are unlikely to produce the measured flux at 60-$`100\mathrm{\mu m}`$. A recent model by Fabian (1999) suggests that the majority of black holes undergo a highly obscured growth phase, and predicts a population of objects at $`z>1`$ which emit predominantly hard ($`E>30\mathrm{keV}`$) X-rays and FIR/submillimeter photons. The Chandra telescope has recently resolved the majority of the 2-10 keV background, finding two new classes of objects: 1) optically “faint” galaxies ($`I>>23`$) with very high X-ray to optical ratios, and 2) point-like, luminous hard X-ray sources in the nuclei of normal bright galaxies showing no other sign of activity (Mushotzky et al. (2000)). The former group is consistent with either early quasars, or extremely dust enshrouded AGN at $`z>2`$. Conceivably some combination of accretion onto black holes and powerful starbursts within the core regions of merging galaxies is capable of explaining the IR excess reported here.
The spectral shape of the IR excess also has interesting implications. Because the $`60/100\mathrm{\mu m}`$ ratio is $`1`$ (in $`\nu I_\nu `$ units) the source must have a high dust temperature of $`T(1+z)28\mathrm{K}`$. For $`z23`$ this temperature indicates a violent process is at work, and argues in favor of a few hot, bright sources. SIRTF will provide valuable information about the luminosity function of such sources at 24, 70, and $`160\mathrm{\mu m}`$.
### 5.2 TeV Gamma Crisis
Another observational conflict is inevitable if the observed IR excess is of extragalactic origin: the observation of gamma rays by HEGRA (Konopelko et al. (1999), Aharonian et al. (1999)) at energies up to $`E20`$TeV from Mkn 501, and by Whipple (Samuelson et al. (1998)) up to $`E10`$TeV. The $`\gamma `$-ray opacity on CIB photons may be approximated by
$$\tau _{\gamma \gamma }(E_\gamma )0.24\left(\frac{E_\gamma }{1\mathrm{TeV}}\right)\left(\frac{u(ϵ_{})}{10^3\mathrm{eV}\mathrm{cm}^3}\right)(z_S/0.1)h_{60}^1$$
(30)
where $`u(ϵ_{})=ϵ_{}^2n(ϵ_{})`$ is the typical energy density in an energy band centered on $`ϵ`$, $`h_{60}`$ is the Hubble constant, and $`z_s`$ is the source redshift (see Coppi & Aharonian (1999)). Integrating the exact cross section over the DEBRA (Diffuse Extragalactic Background Radiation, including CMBR, CIB and optical/UV EBL), Stecker & de Jager (1998) find an optical depth of 2.5 at $`20\mathrm{TeV}`$ photons using a CIB prediction (Malkan & Stecker 1998) based on an extrapolation of IRAS counts and other data. The HEGRA data are consistent with this value at 20 TeV (Konopelko et al. (1999)). However, using the currently accepted measurements of the CIB, Coppi & Aharonian (1999) obtain an optical depth $`\tau 5`$ for $`20\mathrm{TeV}`$ photons. If we modify the CIB spectrum to take account of the 60-$`100\mathrm{\mu m}`$ CIB measurements given in this paper, $`\tau `$ doubles to $`10`$. Even without this doubling, the TeV observations imply that the intrinsic Mkn 501 spectrum is concave upward in the 10-30 TeV range, contradicting the synchrotron self-Compton emission model which requires it to be concave downward. Including the current CIB measurements would imply the intrinsic Mkn 501 spectrum increases a factor of 1000 in $`\nu F_\nu `$ from 10 to $`20\mathrm{TeV}`$. Although very little is yet known about the true intrinsic spectrum of these blazars, this seems an unlikely explanation (however, see Mannheim 1999 for an alternative emission model).
There has also been speculation (Harwit, Protheroe, & Biermann (1999)) that multiple TeV photons may be emitted coherently by blazars such as Mkn 501, and might arrive in Earth’s atmosphere so close in time and space that they are confused with a single-photon event. Such coherent emission seems implausible, given the very large phase space available to TeV photons, but some similar mechanism might yet resolve the apparent conflict between TeV gamma observations and the expected opacity of the CIB.
Still more radical explanations have been proposed. Coleman & Glashow (1997, 1999) have proposed that quantum gravity effects cause a small violation of the invariance principle for very high energy particles. This effect may be large at the Planck scale ($`10^{19}\mathrm{GeV}`$) but even at 20 TeV could have measurable consequences. Kifune (1999) has shown that one possible effect of such a violation is a sudden drop in the effective $`\gamma `$-ray / IR cross section for $`E>10\mathrm{TeV}`$. The resulting dispersion relation may also have observable consequences for $`\gamma `$-ray bursts at cosmological distances (Amelino-Camelia (1998)). Future TeV $`\gamma `$-ray data may rule out these quantum gravity theories, or perhaps demand a further exploration.
### 5.3 Future Data
Data to be gathered within the next few years should be adequate to resolve the problems discussed in this paper. The Space Infra-Red Telescope Facility (SIRTF) will obtain deep number counts at 25, 70, and $`160\mathrm{\mu m}`$ with the Multiband Imaging Photometer (MIPS) (Rieke et al. (1996)). If the integrated flux from measured sources sums to the observed IR excess, then the TeV $`\gamma `$-ray measurements must be reinterpreted. However, if the source density of objects generating the CIB is sufficiently high, then the deep SIRTF images will be confusion limited, and constraints on this diffuse emission will have to come from a fluctuation analysis, in which separation of the extragalactic component from the diffuse cirrus emission of our galaxy will be a limiting factor. Non-detection by SIRTF would confine the emission to our Galaxy, indicating either Galactic halo dust emission, a serious systematic calibration problem in the COBE/DIRBE instrument, a new emission component, or a serious flaw in the present analysis.
Another test is to observe more blazars at still higher energies. By combining data from blazars at different distances and different energies, one can place reliable limits on the DEBRA intensity, as long as the intrinsic blazar spectra are at all similar, or can be predicted by their relation to the X-ray emission from these objects. Work by the HEGRA and Whipple collaborations is underway.
## 6 SUMMARY
Previous attempts by the DIRBE team to measure the cosmic IR background made use of a sophisticated model of the Inter-Planetary Dust (IPD) (Kelsall et al. (1998)). Such detailed modeling was necessary to establish the isotropy of the CIB signal detected at 140 and $`240\mathrm{\mu m}`$. The DIRBE team found excess emission at $`100\mathrm{\mu m}`$ also, but doubted the isotropy of the emission and declined to call it a measurement. In order to recover the valuable information about galaxy formation and evolution contained in the CIB, we have measured the FIR excess at 60 and $`100\mathrm{\mu m}`$.
We analyze the excess DIRBE emission using two different methods. Each of these methods uses a dimensionless parameter derived from the DIRBE data in each week of the mission, parameters that are robust with respect to dust temperature and density variation and detector gain drift. These statistics are nearly insensitive to details of the IPD model. It is not necessary to know the IPD emission for every line of sight at every time, although we use the Kelsall et al. (1998) model as a reference. Method I uses time variation observed in the flux at the ecliptic poles to measure the background at the poles. Method II uses the spatial “shape” of the $`e=90^{}`$ data for each week to remove it and yields an independent measurement of the background in each week. Results derived from these two methods are consistent with each other, giving a background of $`\nu I_\nu =28.1\pm 1.8\pm 7\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`60\mathrm{\mu m}`$ and $`24.6\pm 2.5\pm 8\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ at $`100\mathrm{\mu m}`$.
A variety of arguments rule out alternative sources of emission. Analyses of pulsar dispersion measures and WHAM H$`\alpha `$ data demonstrate that signal from the WIM-correlated dust is already accounted for in the cirrus zero point determined from H I. After removing the cirrus, there is no additional component of this emission correlated with $`\mathrm{csc}|b|`$, arguing against any additional dust slab aligned with the Galaxy. The absence of extended far-IR emission halos around nearby galaxies (e.g. M31) rules out dust emission from extended halos. Unable to find an alternative emission mechanism, we cautiously consider the implications if this excess is an extragalactic background.
The energy required to produce a $`44\pm 9\mathrm{nW}\mathrm{m}^2\mathrm{sr}^1`$ integrated FIR background is large compared to the energy expected from stellar fusion. If the observed flux is indeed of extragalactic origin, then stellar fusion is probably not the dominant source of energy in the universe. Alternative sources such as highly obscured AGN at moderate redshift are a possibility, but would predict that 0.15% of all baryons are in black holes at the present time (although this figure is uncertain by a factor of $`10`$ because of uncertainty in accretion efficiency). A large X-ray background would also be predicted, unless the obscuration is of sufficient optical depth ($`N(H)>10^{24}`$) to block it. The most serious problem with an extragalactic origin of the IR excess is the observation of TeV gamma rays. The opacity of the measured CIB to 20 TeV photons coming from Mkn 501 is 10 optical depths, much greater than the apparent absorption measured by HEGRA (see §5.2). Because of these inconsistencies, there is currently no satisfactory explanation for the observed excess, especially at $`60\mathrm{\mu m}`$. We continue to search for possible sources of emission in the solar system or Galaxy that could account for the observed emission, and urge caution in the use of these results. We eagerly await source counts from SIRTF and the X-ray observatories that might help to solve this mystery.
## 7 ACKNOWLEDGMENTS
We would like to thank Carl Heiles, Chris McKee, Bill Reach, Rick Arendt, and Eli Dwek for helpful discussions. Computers were partially provided by a Sun AEGP Grant. DJS is partially supported by the MAP project and by the Sloan Digital Sky Survey. This work was supported in part by NASA grants NAG 5-1360 and NAG 5-7833. The COBE datasets were developed by the NASA Goddard Space Flight Center under the guidance of the COBE Science Working Group and were provided by the NSSDC.
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# Identification and astrometry of variables in M3
## 1 Introduction
Starting with the discovery of a W Vir star in M3 \[Pickering(1889)\], which was the first pulsating variable to be observed in any globular cluster, M3 deserved the reputation of being the richest globular cluster in variables. The pioneering work of Bailey at the turn of the century yielded more than hundred variables, mostly RR Lyrae stars \[Bailey(1902)\]. By 1939, following discoveries of Shapley, Larink and Müller increased the number of catalogized variables in this cluster to 201, though variability of several candidates was later refuted, and misidentifications were prevalent \[Sawyer(1939), and references therein\]. Observations of Sandage, Kurochkin, Kukarkin, Kholopov, Russev, Meinunger and Kaluzny led to the extension of later editions of the Variable Stars in Globular Clusters to 225 variables for M3 \[Sawyer(1955), Sawyer(1973), and references therein\], and recently to 238 entries \[Clement(1998)\].
However, erroneous identifications are present, partly because the coordinates respect to the cluster’s (not uniformly accepted) central position were taken from various sources, partly due to the severe crowding conditions. Moreover, in some cases astrometry was carried out on photographic plates exposed under imperfect seeing conditions, and by manually centering on the stars’ profiles. \[Evstigneeva et al.(1994)\] compiled a precise list of positions for all known or suspected variables using a homogeneous reference grid. As several new variable discoveries have sprung up since 1994, compilation of a new astrometric list with cross references became substantial.
This work is based on an observing campaign of M3 with a 0.9m Schmidt and a 1m RCC telescope, consisting of several observing runs in the spring seasons of 1998 and 1999. We primarily lay emphasis on identification and astrometry, but variability of the candidates is also confirmed in numerous cases. Photometric aspects and results on individual variables will be dealt with in forthcoming papers \[Benkő et al.(2000), Bakos et al.(2000)\].
## 2 Observations
As concerns the peripheral variables, we used CCD observations obtained with the 60/90/180cm Schmidt telescope at the Piszkéstető Mountain Station of Konkoly Observatory. Observations in V-band were carried out on 11 nights evenly scattered in three months of the the spring season of 1998. The Kodak KAF-1600 $`1024\times 1536`$ chip yielded a $`21^{}\times 28^{}`$ field of view (FOV), with $`1^{\prime \prime }/\mathrm{pixel}`$ resolution. Only the best seeing ($`1.5^{\prime \prime }`$) image was used for astrometry, but all of them were employed so as to verify the variability of the sources.
While identification of outer variables is straightforward in principle, crowding conditions impose difficulties in the inner few arcminutes. A subset of observations carried out with the 1m, f/13.6 RCC telescope at the same observatory was used to patch astrometry in the dense, inner region. The TH7896M UV-coated $`1024\times 1024`$ CCD chip attached to the telescope yielded a FOV of $`5^{}\times 5^{}`$, with $`0.288^{\prime \prime }\mathrm{pixel}`$ resolution. The observing run in March 1999 produced time-series data for 5 nights with strongly varying seeing conditions (ranging from $`1.2^{\prime \prime }`$ to $`2.5^{\prime \prime }`$). Only the best night, and only V-filter observations were used in the present work.
As the innermost region of M3 was still not clearly resolved, we used archive Hubble Space Telescope (HST) Wide Field Planetary Camera 2 (WFPC2) observations<sup>1</sup><sup>1</sup>1 Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute. STScI is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract NAS 5-26555. of M3 (proposed by Fusi Pecci, 1994, cf. Ferraro et al. 1997). This is the only set of publicly available WFPC2 observations where frames have short enough exposure times ensuring that RR Lyrae variables are not saturated. We used an F814W mosaic image with 3s integration time designated as “u2li010dt” in the archive.
The standard iraf/ccdred package was used for overscan and flatfield correction of the Schmidt and RCC images. The HST archive image was already calibrated by the HST pipeline, and while we were not performing direct photometry on the image, recalibration was needless.
## 3 Astrometry
All astrometric work was based on the USNO-A2.0 catalogue, which had been derived from PMM (Precision Measuring Machine) scans of POSS-I O and E plates \[Monet(1996), Monet(1998)\]. The typical astrometric error is about $`0.15^{\prime \prime }`$ for individual stars away from the plate corners and not in crowded fields. After thorough examination of crowding as a function of radial distance from the center of M3, we selected USNO-A2.0 stars with radial distance $`r>6^{}`$ and brighter than $`M_{red}=17.8^m`$ (hereafter grid0). Uncertainties in USNO astrometry of stars closer to the center might have led to considerable errors, ie. it was not possible to use USNO stars for reference in the RCC and HST fields. By careful choice of overlapping regions on the Schmidt, RCC and HST frames, the astrometric reference grid was extended towards the center. Equatorial coordinates given in the International Celestial Reference Frame (ICRF) were transformed to relative coordinates ($`\mathrm{\Delta }\alpha `$ in $`15\times \mathrm{seconds}`$<sup>2</sup><sup>2</sup>2This unit was chosen for convenience, as it has the same order of magnitude as arcseconds for small declinations. Values in Table 1 are given in normal second units., $`\mathrm{\Delta }\delta `$ in arcseconds) respect to the center<sup>3</sup><sup>3</sup>3Note that this center slightly differs from the one accepted in \[Clement(1998)\] precessed for J2000.0 of M3 (J2000.0): $`\alpha _{M3}=13^\mathrm{h}42^\mathrm{m}11.2^\mathrm{s}`$, $`\delta _{M3}=28^{}22^{}32^{\prime \prime }`$ \[Harris(1996)\].
Using our self developed iraf/gastro astrometric package<sup>4</sup><sup>4</sup>4All self-written software are available from the first author via request in e-mail., we selected isolated, not saturated stars with $`r>6^{}`$ on the Schmidt image (grid1a). The centers of these sources were determined by a script (gast\_mkgrid) built up from the iraf/daofind, phot, pstselect programs.
Using a best estimate of the transformation between grid0 and grid1a (initially a manual guess by gast\_crude), stars were cross-identified, and transformations were refined. This iteration was done several times with variable rejection thresholds and fitting parameters, see-sawing back and forth between identification and refinement (gast\_prec). Finally we fitted a second-order transformation between grid0 (ICRF) and grid1a (Schmidt) of the form
$`X_{ICRF}`$ $`=`$ $`A_{11}+A_{21}X_{Sch}+A_{12}Y_{Sch}+A_{31}X_{Sch}^2+A_{22}X_{Sch}Y_{Sch}+A_{13}Y_{Sch}^2`$
$`Y_{ICRF}`$ $`=`$ $`B_{11}+B_{21}X_{Sch}+B_{12}Y_{Sch}+B_{31}X_{Sch}^2+B_{22}X_{Sch}Y_{Sch}+B_{13}Y_{Sch}^2,`$
which adequately conforms to the large-scale distortion of the Schmidt-camera field. The typical scatter of least-squares fits based on $`150`$ evenly scattered stars was lower than $`0.13^{\prime \prime }`$ rms, which is even better than the USNO-A2.0 precision itself.
As USNO-A2.0 astrometry is not reliable in RCC fields covering the dense center of the cluster ($`5^{}\times 5^{}`$), we extended our Schmidt grid to the innermost $`3^{}>r>1.5^{}`$ ring (grid1b), and in a similar manner, established a grid on the RCC fields with $`r>1.5^{}`$ (grid2a). The transformation between the two grids, based on 250 common stars, was derived similarly as previously mentioned, and had an overall scatter less than $`0.04^{\prime \prime }`$.
Finally, transformations were calculated between the RCC images (grid2b: $`r<1.5^{}`$) and the HST PC, WF2, WF3, WF4 chips, using 70, 150, 150 and 200 stars, respectively. Scatter in the transformations was less than $`0.04^{\prime \prime }`$.
On the whole, it became possible to transform any astrometric position from the HST, RCC and Schmidt images to equatorial system (ICRF) with $`0.15^{\prime \prime }`$ precision.
## 4 Identification of variables
### 4.1 Tools
Four basic tools were applied to identify variables: coordinate lists, finding charts, variability images and light-curves. Relative coordinates for 238 variables are given in \[Clement(1998)\]. Unfortunately these coordinates are collected from several sources, thus far from being homogeneous, and the precision is only $`\pm 0.11.0^{\prime \prime }`$, which is not eligible for unambiguous identification in the center. Relative coordinates in right ascension are (traditionally) given in arcseconds as opposed to seconds, which involves a projection of the spherical equatorial coordinate system. \[Evstigneeva et al.(1994)\] give a more accurate list with proper equatorial coordinates. To be on the safe side, we used both lists to filter out misidentifications. New HST discoveries from \[Guhathakurta et al.(1994)\], properly transformed to equatorial system by \[Goranskij(1994)\] were also used. Finding charts are enclosed only in a few sources, such as \[Kholopov(1963), Kholopov(1977)\] and \[Kaluzny et al.(1998)\].
As concerns the variability images, we applied the Image Subtraction Method (ISM) for the best night on the RCC frames, and for all 11 nights on the Schmidt frames, as described in \[Alard and Lupton(1998)\] and \[Alard(1999)\]. After registering and convolving the images, then subtracting them from the reference frame with the aid of isis2.0, the resultant difference images were almost blank, only variable stars having negative or positive profiles. The variability image \[Olech et al.(1999)\] was the average of the absolute values of several difference frames, containing accumulated contributions from all variations respect to the reference frame, that is showing all sources variable on the timescale of the observations ($`40`$ and $`25`$ images for the Schmidt and RCC observations, respectively). In the case of the RCC observations, we restricted our investigation to variability detectable in one single night (mostly RR Lyraes and perhaps SX Phe stars), data reduction of subsequent frames will improve identification of both longer period (RGB) and fainter variables. On the other hand, the Schmidt observations spanning three months were also convenient for detecting long period irregular variables. It has to be emphasized that the primary goal was identification of the known variables and suspected ones, and not a complete variable search, although six new variables were found. In some complex cases (mostly close doubles) we also made use of scrutinizing the dependence of the light-curves on the (fixed) center of the photometry; the smaller scatter of the curve implies that the center of photometry is closer to the real position of the variable, which can serve as a criterion for selecting the truly varying star.
We implemented our iraf/gisis/gvarfind task, which blinks the normal and the variability images, thus identification can be judged not only upon the initial coordinate-estimate but also on the variability of the source. This method seemed to be extremely useful in the central areas, where sometimes several candidates were equidistant from the corresponding rough position. The aforementioned task also performed fitting of an elliptical gauss profile to the selected source, either on the normal image or on the variability image. If variability was sufficient, the latter case was preferred, as non-variable neighbors completely disappeared even in the dense regions, where profile-fitting would have been biased by overlapping profiles. It is worthy to note that HST variability image was not available, as we had only a single image with not saturated RR Lyrae stars.
In all cases identification meant determining the position of the variable by precise psf-fitting, and adequately transforming the coordinates to the ICRF.
Photometric accuracy is an important aspect when confirming variability or invariance of a star. Investigation of light curves derived by aperture photometry on the Schmidt and RCC frames yielded a crude estimate of the photometric errors being smaller than $`\delta V0.05^m`$ and $`\delta V0.03^m`$, respectively. However, significantly smaller accuracy could have been achieved in the most crowded fields.
As the Image Subtraction Method is less sensitive to crowding conditions, the accuracy of the (ISM) light curves we checked for variability were in the order of hundreds of mags in all cases. Variability of each star was checked and commented only if no light-curve had been previously given.
### 4.2 Finder charts
We present finder charts for 286 variables and suspected variables within the boundaries of the RCC frame (Figs. 1 to 10), or within the Schmidt frame (Figs. 11 to 15), to avoid further misidentifications and assist in any variability study of M3. Variables on the very edges of the Schmidt frame were omitted (V17, V92, V117), as well as variables denoted by “E” in Table 1 (V82, V91, V112, V113, V114, V115, V123, V141, V205, V206, V230). Although variable vZ1283 was identified on the edge of the RCC frame, the finder chart is given on the Schmidt image. Two HST variables (GU9016, GU9025) were not included as they were not readily visible on the RCC images (both identified from the HST frame). Several candidates revealed as doubles on the HST images were not resolved on the RCC frame, in which cases only one chart was given (V159a,b, V192a,b, KG8a,b, X14a,b). Though good quality finder charts for V237, V238 has been recently published by \[Kaluzny et al.(1998)\] they are repeated here for completeness.
#### Figures 1.-10
Caption: Finding charts for variables on the RCC frames. Boxes are approximately $`15^{\prime \prime }`$ in width, North is up, East is to the left. Boxes from very sparse regions labeled as “WF” are $`30^{\prime \prime }`$ wide. Different intensity scales are choosen to reach the clearest resolution of each frame. (Note: figures are enclosed in png format for compression)
#### Figures 11.-15.
Finding charts for variables on the Schmidt frames. Boxes are approximately $`80^{\prime \prime }`$ in width, North is up, East is to the left. (Note: figures are enclosed in png format for compression)
The revised positions and some additional information on the variables are given in Table 1. The IDs of the variables were extended in a similar manner to the \[Clement(1998)\] catalogue in case of all confirmed variables (column 1.). The last column gives all the crucial cross-references. The type of variability indicated in the second column of Tab. 1 was determined from our observations if light variation was detected. For some of the HST variables classification given by \[Guhathakurta et al.(1994)\] was adopted. Columns 3–5 and 6 give the relative coordinates and the observation used for the identification, respectively. Any comments of the individual stars are indicated in column 7.
Whenever possible, variables were identified on the HST image, as this yielded the most precise astrometry due to the $`0.05^{\prime \prime }`$ and $`0.1^{\prime \prime }`$ per pixel resolution for the PC and WF chips, respectively. Furthermore, due to the enhanced resolution, several variables turned out to be close double stars, in which cases original notation was split (eg. 122a, 122b)<sup>5</sup><sup>5</sup>5Notation was also splitted in case more than one candidates were equidistant from the position given in previous references, and identification was dubious.. As we could not check variability on the HST frames, we tried to identify the variable component using the profile on the RCC variability frame. If a star was not within the WFPC2 boundaries, RCC images were checked, and in case the star was beyond these frames, Schmidt images were used. In eleven cases when variables were even outside the Schmidt frames, we accepted positions from \[Evstigneeva et al.(1994)\]. Still, our grid of variables remained self-consistent, as our coordinates are compatible with \[Evstigneeva et al.(1994)\] in the sense that no discernible offset is present (see later).
Kholopov variables (X) were manually selected with gvarfind using finder charts in \[Kholopov(1963), Kholopov(1977)\] and our variability image. Some of the Kholopov variables were already included in \[Clement(1998)\], without firm constraints on their light-curves. All these cases were checked for variability. We identified the remaining 11 suspected candidates, and confirmed variability of five of them.
\[Kadla and Gerashchenko(1980)\] variables (KG) were identified from the list in the original paper, and from \[Evstigneeva et al.(1994)\]. Variability of several candidates is confirmed, and new ID numbers are given.
All von Zeipel (vZ) candidates listed in \[Evstigneeva et al.(1994)\] were identified on the RCC and Schmidt images. As all of them were suspected as long-period variables \[Welty(1985), Meinunger(1980)\], variability was checked on the Schmidt frames. Variability is ascertained only in one case (vZ297), while six other candidates remain suspected variables. The remaining stars are constant within the photometric errors of the Schmidt measurements.
Sandage variables (S-) were identified with the aid of finder charts in \[Sandage(1953)\]. Except for S-AQ, S-I-II-52, S-I-II-54, these stars were suspected as bright, long-period variables \[Meinunger(1980)\], however all, but one (S-I-VI-65) proved to be constant on the Schmidt frames. Bao-An et al. (1993a,b) and \[Bao-An et al.(1994)\] drew attention to three new short-period, low-amplitude variables (S-I-II-52, S-I-II-54, S-AQ). These stars were identified from charts in \[Sandage(1953)\], but low amplitude variability has not yet been confirmed on the Schmidt images.
\[Guhathakurta et al.(1994)\] gave a list of variable stars (GU) in the center of M3 using HST/WFPC observations, giving unique ID numbers, relative coordinates respect to star AC999, rough average magnitudes, and classification. The list was appropriately transformed by \[Goranskij(1994)\] to match \[Evstigneeva et al.(1994)\], and cross-identification showed that only 11 out of the 40 variables were new discoveries. By performing the same analysis, we reduced the number of new HST discoveries to 9 (confer revised cross-references in Table 1). Variability is confirmed except for GU1711 (long-period suspected variable) and GU9016 (faint SX Phe star).
Three prominently variable RR Lyrae stars were found on the RCC images, and three long-period variables on the Schmidt images, without any previous reference. These supposedly new variables are listed at the end of Table 1. Preliminary light-curves for the newly discovered variables 269–274 are presented in Fig. 16.
Coordinates in the resultant list were correlated with the original list of \[Clement(1998)\], and with the exception of few misidentifications and defective coordinates in the above reference (V146, V151, V165, V181, V196, V199, V204), the overall scatter of the fits was $`1^{\prime \prime }\mathrm{rms}`$. Correlation with the list of \[Evstigneeva et al.(1994)\] yielded considerably more precise fits with $`0.2^{\prime \prime }\mathrm{rms}`$. This reinforces the absolute precision of our astrometry being better than $`0.2^{\prime \prime }`$, even with taking the uncertainties rising from the USNO-A2.0 catalogue into account. As concerns relative astrometry, our grid is self-consistent within the error of $`0.15^{\prime \prime }`$.
## Acknowledgments
We are grateful to Dr. Béla Szeidl for his persevering help and expertise in M3 variables. This project was supported by OTKA grants T-24022 and T-30954.
## Appendix A Notes on individual stars
V2: Two closeby candidates in \[Evstigneeva et al.(1994)\] (hereafter ESTS94); none of them is variable. From charts in \[Kholopov(1963), Kholopov(1977)\] (hereafter K63 and K77) it is not obvious which star is marked as “2”. Original positions from \[Sawyer(1973)\] (hereafter SH73) and \[Clement(1998)\] (hereafter CC98) are closer to the west candidate.
V29: Positions of V29 and V155 are interchanged in ESTS94.
V122: Two candidates; V122a is more likely the true variable, based upon its light-curve. Merging with V229.
V127: CC98 position merges with position of V222. K77 gives correct identification.
V146: CC98 position is $`4^{\prime \prime }`$ away from the variable. As V127 and V222 are closeby, this position error might be misleading. Identification confirmed from K77 charts.
V151: Quite far ($`3^{\prime \prime }`$) from the CC98 position, the same star is marked in K77. No other candidate.
V154: W Vir star, merging with an RR Lyrae V268.
V155: Positions of V29 and V155 are interchanged in ESTS94.
V156: Formal position of GU9019 is $`1.3^{\prime \prime }`$ away from V156, not coinciding with any star. It is likely to be the same variable as V156.
V158: Despite the entry in CC98, it doesn’t seem to be variable on the RCC images (one night), neither on the Schmidt images.
V159: Double on HST/WF2 images: V159a, V159b. Not possible to select the variable component.
V161: Misidentified in ESTS94 (with non-variable star in $`2^{\prime \prime }`$ distance).
V163: Star from \[Müller(1933)\]; identification is dubious, position is equidistant from a non-variable and a variable star (later V180). K63 and K77 marks the non-variable candidate as 163. CC98 says V163 is not var. The not variable candidate is accepted as V163, in order to avoid confusion.
V165: The coordinates in SH73, CC98 are erroneous (sign on declination offset is “$``$” instead of “+”)
V170: Elongated profile on HST/WF2 image. Maybe double?
V179: This candidate is very far from the cluster, only crude position is known, charts are not available. Three stars are equidistant on the POSS images from the quoted position. CC98 notes that this star is non-variable.
V180: Two equidistant candidates close to SH73 (Sh) position: one is the same as X13=(V239), which is variable indeed. K77 marks the other star as 180, and announces X13 as a new variable. Shapley says that V180 is non-variable. For convenience, the candidate shown in K77 is chosen to be V180. Variability is confirmed, as opposed to CC98.
V181: Appropriately transformed CC98 position coincides with a variable star. However, CC98 notes that this star is not variable. K77 finder charts mark a closeby, non-variable star as 181. ESTS94 identifies the same star as V181. We propose that both K77 and ESTS94 identifications are erroneous, and V181 is variable, identical to GU9008. Faint companion to the NW on HST/WF2 image.
V185: No variable star in the vicinity of the position in CC98. V185, as identified in ESTS94 shows no variability.
V192: Obviously double on HST/PC image; not possible to select the variable.
V194: Two candidates on HST/PC frame. Merging with other stars to the N in a NS chain, all of them showing variability? Time-series observations with better seeing would be needed for firm identification. Formal position of KG10, as derived from \[Kadla and Gerashchenko(1980)\] (hereafter KG80) and ESTS94, is merging with V194.
V196: CC98 position is close to a constant star, but it is not the same as the one marked in K77 finder charts. The latter is also constant (196c). The star close to the CC98 position, is a double star on the HST/WF4 image (196a, 196b), both components are constant.
V198: Despite the entry in CC98 and the short period ($`0.2797^d`$), it doesn’t seem to be variable on the RCC images (one night).
V200: Misidentified in ESTS94, merging with a brighter star.
V204: Roughly $`6^{\prime \prime }`$ away from the position in CC98; identification from K77 chart. No significant variability detected on the RCC frames (one night), neither on the Schmidt frames.
V209: ESTS94 lists two stars at the position. Unfortunately the area is not within the HST fields. Variable is more likely to be the southern one (from variability image).
V211: Obviously double on HST images.
V229: Merging with V122a and V122b.
V237, V238: Too faint stars for variability to be confirmed on Schmidt images.
V239=X13: Confer comments at V180.
V240=X14: Obviously double on HST/WF2 images: X14a, X14b. Maybe both components are variable, as the source is elongated on the variability image.
V241=X17: See comment at V262.
X18: K77 marks a relatively bright star as X18, which shows no significant variability on the RCC frames spanning one day. Very close to formal position of V251 (=KG3).
V242=X20: Merging with a bright star, and close to V257.
V244=X23: Two candidates on HST/WF2 image: the true variable was selected using the RCC variability image. X23 is not identical with GU9012, as proposed by \[Goranskij(1994)\].
V248=X36: Elongated profile on HST/WF2 image.
V249=KG1: Identification with the help of KG80 positions adequately transformed to our system.
V250=KG2: Two close candidates on HST/WF2 image.
V251=KG3: Merging with X18, only resolved on HST/PC frames. Variability (due to the extreme crowding) not detected unambigously.
KG5: Bright star, identified using both ESTS94 and KG80: no significant variability.
V254=KG8: Two candidates on HST/WF4 image, merging on gound-based images. Light curves are almost the same for the two slightly different positions. Further investigation in good seeing would be needed to draw firm conclusion.
V256=KG11: Two candidates on HST/WF4 images; the true variable was selected using the RCC variability image.
KG13: Identification dubious; two bright stars and a fainter companion close to the position showing no significant variability.
V259=KG15: Two candidates on HST/WF4 images; the true variable was selected using the RCC variability image.
vZ380: Star in ESTS94, variability from \[Meinunger(1980)\] – obj. a.
vZ1066: Star in ESTS94, variability from \[Meinunger(1980)\] – obj. b.
vZ1283: Star in ESTS94, variability from \[Meinunger(1980)\] – obj. c.
S-AQ, S-I-II-52, S-I-II-54: Low amplitude variability \[Bao-An et al.(1993a)\], \[Bao-An et al.(1993b)\], \[Bao-An et al.(1994)\] is not detected on the Schmidt images.
V262=GU552: This might be a very close companion to V241, only resolved on HST images. Light curve is almost identical to that of the closeby V241, due to merging. No trace of RRc variation is detectable in the light-curve, in spite of comments in \[Guhathakurta et al.(1994)\]. Higher resolution time-series observations would be needed to conclude more on its variability.
V263=GU576: SX Phe, barely visible on the 1m RCC images. Position from HST data. Elongated on HST/WF2?
V265=GU1489: RR Lyrae, merging with V154 and V268.
GU1711: Suspected variable in \[Guhathakurta et al.(1994)\], constant on the RCC images.
V267=GU9016: SX Phe star merging with V224, resolved on the HST/WF2 frame.
V268=GU9025: RR Lyrae star merging with V154 (W Vir) and V265.
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# 1 Introduction:
## 1 Introduction:
To understand the nature of the dark matter is one of the pressing unsolved mysteries in cosmology and particle physics nowadays. The dark matter problem is the fact that the matter in luminous forms (stars and gas) is inferred to account only for $`\mathrm{\Omega }_{lum}0.007`$, i.e. for less than one percent of the critical density. On the other hand, dynamical estimates based on the gravitational influence of the mass on test bodies (gas orbiting galaxies, galaxies moving inside clusters) imply that there is at least ten times more ‘dark’ mass than luminous one on galactic scales, and even more ($`\mathrm{\Omega }0.3`$) on cluster scales. The issue on even larger scales is still unsettled, existing at present indications in favor of some dark energy (cosmological constant) accounting for $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$, i.e. filling the gap up to the inflation preferred value $`\mathrm{\Omega }=1`$ corresponding to a flat universe.
The possible constituents of the dark matter naturally split into baryonic and non-baryonic candidates. The last ones would be some kind of weakly interacting particle permeating the universe, such as supersymmetric neutralinos, axions or massive neutrinos, while the first ones would be made of just ordinary protons and neutrons hidden in some non-luminous forms. The simplest examples to achieve this would be to have them in MACHOs (the acronym for massive astrophysical compact halo objects), such as stellar remnants (white dwarfs, neutron stars, black holes), brown dwarfs (stars with mass $`m<0.1M_{}`$, which never become hot enough to start sustained nuclear fusion reactions) or even planets. Also cold gas clouds have been suggested as dark baryonic constituents of the galactic halos.
The main constraints on the amount of baryons in the universe come from the theory of primordial nucleosynthesis, which requires $`\mathrm{\Omega }_b0.01`$$`0.05/(H_0/70`$ km/s/Mpc)<sup>2</sup>, with the lower (higher) values corresponding to the high (low) primordial deuterium abundance determinations from QSO absorption lines . Hence, we see that for a Hubble constant $`H_050`$–70 km/s/Mpc nucleosynthesis indicates that dark baryons should exist and that their total amount is probably insufficient to account for dark galactic halos ($`\mathrm{\Omega }0.1`$). Clearly to account for the observations on cluster scales non-baryonic dark matter is also required, and it would then be natural to expect to have some amount of it also at galactic scales.
Although dark compact objects are, by definition, very hard to be seen directly, they may reveal themselves through the gravitational lensing effect they produce on background stars. Indeed in 1986 Paczynski showed that by monitoring a few million stars in the Large Magellanic Cloud (LMC) during a few years would allow to determine whether the dark halo of the Galaxy consists of MACHOs with masses in the range $`10^7`$$`10^2M_{}`$, covering essentially all the range of suggested candidates. Soon after, with the advent of CCDs this program became feasible and several groups started the searches at the beginning of the nineties.
## 2 Microlensing expectations:
The gravitational deflection of light by massive bodies is one of the predictions of general relativity which is now tested to better than the percent level. This phenomenon leads to several beautiful effects such as the multiple imaging of QSOs by intervening galaxies, giant arcs around clusters or weak lensing distortions of faint background galaxies, which have the potential of giving crucial information for cosmology (measurement of $`H_0`$ from the time delays between different QSO images, estimates of cluster masses, etc.). When the deflector is a compact star like object, a MACHO, two images of the background sources are produced, one on each side of the deflector. In the particular case of perfect lens-source alignment, the image is a ring with angular size (Einstein’s angle)
$$\theta _E\mathrm{mas}\sqrt{\frac{m}{M_{}}\frac{(1x)10\mathrm{kpc}}{D_o\mathrm{}}},$$
where taking $`D_{os}`$ as the distance from the observer to the source, $`D_o\mathrm{}xD_{os}`$ is the distance to the lens. Hence, for lenses at a few kpc distance with masses in the range of interest for the MACHO searches this angle is much smaller than the typical telescope resolution ($`40`$ mas for HST!) and hence cannot be resolved. However, the important effect is that the brightness of the source is magnified due to a lensing effect (for reviews see refs. ) by a total amount
$$A=\frac{u^2+2}{u\sqrt{u^2+4}}$$
where $`u^2=\xi ^2/R_E^2`$, with $`R_E=\theta _ED_o\mathrm{}`$ being the radius of the Einstein ring in the lens plane and $`\xi `$ is the distance between the lens and the line of sight to the source. Due to the lens relative motion orthogonally to the l.o.s. with velocity $`v^{}`$, one has $`\xi ^2=b^2+v^2(tt_0)^2`$, with $`b`$ the impact parameter of the lens trajectory and $`t_0`$ the time of closest approach. Hence, the amplification will vary with time in a very specific form and it is sizeable ($`A>1.34`$) as long as the lens distance to the line of sight is not larger than $`R_E`$. Hence, one can estimate the optical depth, i.e. the probability that a given star is magnified significantly for a given lens distribution with density $`n_{\mathrm{}}`$, as the number of lenses within a distance $`R_E`$ from the l.o.s.. If the lenses are assumed to have a common mass $`m`$, one has $`n_{\mathrm{}}=\rho _{\mathrm{}}/m`$, and hence
$$\tau =\frac{4\pi }{m}_0^{D_{os}}dD_o\mathrm{}R_E^2\rho _{\mathrm{}}.$$
In the case of a halo consisting fully of MACHOs, the lens density will be just the halo density, which from the observed rotation curve of the Galaxy should be
$$\rho _H\rho _0^H\frac{r_0^2+a^2}{r^2+a^2},$$
with $`r_0=8.5`$ kpc the distance to the galactic center, the core radius $`a`$ is of a few kpc and the local halo density is $`\rho _0^H10^2M_{}`$/pc<sup>3</sup>. This would lead to a predicted optical depth for LMC stars of $`\tau ^H4\times 10^7`$.
The timescale of the events, defined as $`TR_E/v^{}`$ (and determined from a fit to the light curve with the theoretical expression for $`A(u)`$) has an average value $`T65\mathrm{d}\sqrt{m/M_{}}`$ given the typical velocity dispersion expected for halo objects ($`\sigma 155`$ km/s). One can show that for 100% efficient searches the number of expected events is $`N=(2/\pi )(\tau /T)\times \mathrm{Exposure}20\mathrm{events}\sqrt{M_{}/m}(\mathrm{Exposure}/10^7\mathrm{stars}\mathrm{yr})`$. Actually, the efficiencies are $`30\%`$ and are time dependent, so that a detailed prediction requires to convolute the differential rates with the efficiencies (and also with the mass distribution of the lenses).
## 3 First microlensing results and their interpretation
By 1992 the observational searches towards the LMC by the french EROS and australo-american MACHO collaborations had started and already after one year the first candidate events had showed up. It soon became apparent however that the number of events observed were a factor $`5`$ below the expectations from a halo made of objects lighter than a few solar masses. The estimate of the optical depth to the LMC from the first 3 MACHO events was indeed $`\tau =8.8_5^{+7}\times 10^8`$ and the one from the first two EROS events was similar, $`\tau =8.2_5^{+11}\times 10^8`$. The typical event durations of a few weeks were also indicative of lens masses in the ballpark of 0.1 $`M_{}`$, i.e. in the limit between brown dwarfs and very light main sequence stars. Furthermore, the non observation of short duration events ($`T`$ less than a few days) by the EROS CCD program and by MACHO set stringent bounds on light lenses ($`10^7<m/M_{}<10^3`$) which could contribute no more than $`25`$% to the overall halo mass .
Although the number of events were few, they were significantly larger than the backgrounds expected from the known stellar populations. For instance, the faint stars in the galactic disc would contribute little to the optical depth towards the LMC, because although the local disc density is an order of magnitude larger than the halo one ($`\rho _0^D10^1M_{}/`$pc<sup>3</sup>) its scale height is very small ($`h100`$–300 pc). In general, adding an exponential disc like population with local column density $`\mathrm{\Sigma }`$ leads to an optical depth to the LMC of $`\tau ^D10^8(\mathrm{\Sigma }/30M_{}/`$pc$`{}_{}{}^{2})(h/300`$ pc) . The column density in the thin disc is $`50M_{}/`$pc<sup>2</sup>, of which 20% is in gas. Furthermore taking into account the measured mass function of disc stars it can be shown that a large fraction of $`\tau ^D`$ would be due to stars too bright to go undetected , and hence the thin disc predicted depth is actually below $`10^8`$.
The Galaxy is also known to have a thicker disc of stars, with $`h1`$ kpc, probably originating from a past merger of a satellite which somehow heated the disc stars. The column density associated to the observed thick disc stars is typically taken as $`\mathrm{\Sigma }^{TD}4M_{}/`$pc<sup>2</sup>, and hence the expected optical depth is tiny ($`0.5\times 10^8`$). However, in ref. it was noticed that the maximum allowed column density consistent with dynamical observations, $`\mathrm{\Sigma }_{max}^{TD}50M_{}/`$pc<sup>2</sup>, would lead to $`\tau 5\times 10^8(h/`$kpc), and hence for the scale height adopted there (1.4 kpc) it would already be close to the observed value.
Another known stellar galactic population is the spheroid (or ‘stellar halo’), observed through the old metal poor stars at high latitudes and also as high velocity stars nearby (the observed velocity dispersion is $`\sigma 120`$ km/s). It was formed in the first Gyrs of the Galaxy lifetime, has a probably slightly flattened spherical shape ($`c/a0.7`$–1)with density profile $`r^\alpha `$, with $`\alpha 2.5`$–3.5 . The predicted optical depth for this population, adopting $`c/a=1`$ and $`\alpha =3.5`$, is $`\tau ^S0.3\times 10^8(\rho _0^S/10^4M_{}/`$pc<sup>3</sup>). This value would increase by 50% for $`\alpha =2.5`$ and decrease by $`25`$% for $`c/a=0.7`$. The spheroid density inferred from star counts, $`\rho _0^S5\times 10^5M_{}/`$pc<sup>3</sup> leads then to a tiny optical depth. However, there have been since many years suggestions that the spheroid density should be larger by an order of magnitude, based on global fits to the galactic rotation curve and other dynamical observations . In these heavy spheroid models the optical depth results $`4\times 10^8`$, and is hence relevant for observations towards the LMC .
Another background for LMC searches is that of old neutron stars . These are born in the disc, but they are believed to acquire large velocities in the supernova explosions producing them and hence move to large distances where they become more efficient lenses. For the typical value $`N_{ns}2\times 10^9`$ for the number of past galactic core collapse supernovae , the resulting optical depth is $`\tau ^{NS}0.4\times 10^8`$, comparable to the one from known stars in the spheroid or thick disc.
Regarding the LMC populations, in the same way as the Milky Way has a dark halo, the rotation curve of the LMC suggests that this galaxy also has a dark halo around it. The optical depth associated to it amounts to $`8\times 10^8`$ . Hence, if all the dark matter were in MACHOs the total optical depth would be 20% higher than the prediction from the Milky Way halo alone, and this clearly worsens the discrepancies with the observations. It has to be stressed that no stellar population has been observed with the expected LMC halo distribution. The predictions for the observed distributions of LMC stars (disc and bar) are more delicate. It was actually suggested that the optical depth towards the central bar could be $`5\times 10^8`$ , and hence of the order of the observed $`\tau `$. However, outside the bar, where several of the observed events actually are, $`\tau `$ falls significantly, and the most recent estimates for the expected average depth are $`\tau 2.4\times 10^8`$, although with some spread among different models.
Summarizing, the expectations from known stellar populations amount to $`\tau 3`$–4$`\times 10^8`$, while from the Milky Way and LMC halos one would expect $`\tau 5\times 10^7`$. The observed value $`8\times 10^8`$ was clearly larger than the first one, but only a small fraction of the second. Since the main purpose of the microlensing searches was to establish whether the halo consisted of MACHOs, the results are most commonly presented in terms of the fraction $`f`$ of the Milky Way halo in the form of compact objects, which for the observed $`\tau `$ would correspond to $`f20`$%. The remaining fraction could be just in cold gas clouds or in non-baryonic dark matter. An alternative explanation was to have instead a heavy spheroid or thick disc made mainly of MACHOs or a very large LMC self-lensing contribution, and instead a completely non-baryonic halo.
## 4 The second period (1996-2000)
In 1996, the results of the analysis of the first two years of MACHO data were announced . A total of 8 candidate events were observed and the average event duration turned out to be longer. The conclusion was that now $`\tau =2.9_{0.9}^{+1.4}\times 10^7`$ and also that the typical lens masses increased to $`m0.5M_{}`$. The picture which emerged then was that 50% of the halo should be in objects with the characteristic mass of a white dwarf. This scenario was quite unexpected, and was shown to be also potentially in trouble with the chemical enrichment of the galaxy due to all the metal pollution which would result from the white dwarf progenitors . There were also attempts to explain the observed rates as due to some tidal debris of the LMC or even due to a previously undetected satellite galaxy along the l.o.s. to the LMC , but it is not clear whether these ideas are actually supported by observations .
In the period after 1996 significant improvements on the observational programs took place: the EROS and OGLE collaborations started to use better cameras in new telescopes and networks of telescopes were organized to follow ongoing alerted microlensing events (GMAN, MOA, PLANET, MPS). This allowed for instance to measure in great detail an event in the Small Magellanic Cloud caused by a binary lens . When the lens is a binary, the signatures are quite different from the single lens case. For large impact parameters there are now three images of the source, one on the exterior side of each of the lenses and the third one in between the lenses. An extra pair of images can appear however for small impact parameters when the source crosses the so called caustic, and then two images disappear again when the source leaves the region delimited by the caustic. At these caustic crossings the magnification formally diverges. This means that one has an extremely powerful magnifying glass to look at the source. Actually, effects associated to the finite size of the source become observable and in particular they limit the maximum amplification to a finite value. If one knows the radius of the source, one can also infer the relative lens-source proper motion $`\mu `$, and since this quantity is expected to differ significantly for lenses in the Magellanic Clouds ($`\mu 1`$ km/s/kpc) or in the galactic populations ($`\mu 10`$ km/s/kpc), it was possible to establish that the binary lens belonged to the SMC and not to the halo. Also one of the MACHO LMC events was a binary and its proper motion again suggested that it belonged to the LMC. This shows that indeed the contribution of the lenses in the Magellanic Clouds is significant, and actually in the SMC it is expected to be relatively larger than in the LMC due to the elongated shape of the SMC along the l.o.s..
## 5 Recent developments
The most recent microlensing results appeared a few months ago and again changed the overall picture. The analysis of 5.7 yrs of MACHO data in 30 fields (out of a total of 82) found 13–17 events (depending on the criteria used), leading to a significantly reduced depth $`\tau =1.1_{0.3}^{+0.4}\times 10^7`$ (for the 13 event sample). Another relevant observation was that no particular increase in $`\tau `$ towards the center of the Cloud was observed, contrary to the expectations from a rate dominated by LMC stars. Also the EROS group presented the reanalysis of their old data together with two new years of EROS II data , finding a total of only 4 events, a result inconsistent with the ’96 MACHO result. This shifted the situation somehow back to the initial stage in which the observed rates are a factor $`5`$ below the halo predictions but $`3`$–4 above the expectations from known populations. The inferred masses $`m0.5M_{}`$ continue to suggest however that the objects could be old white dwarfs.
Another important related result has been the recent activity related to the direct search of old white dwarfs. Hansen realized that the spectra of old white dwarfs having molecular hydrogen atmospheres (i.e. probably half of the total) would look much bluer and brighter than previously believed due to the strong absorption at wavelength larger than 1 $`\mu `$m by their atmospheres. This prediction was actually confirmed by the direct measurement of the spectrum of a cool ($`T3800^{}`$K) nearby white dwarf . With this new scenario the search for old white dwarfs becomes feasible, and indeed analyses of two Hubble Deep Fields taken two years apart were done searching for objects with large proper motions . The candidate objects they found, with colours consistent with being old ‘halo’ white dwarfs, led them to infer that their density could be comparable to the local halo density ($`10^2M_{}/`$pc<sup>3</sup>). However, recently two new searches in larger nearby volumes (not so deep but wider) have found results inconsistent with such large white dwarf densities. Flynn et al. found no candidates while expecting a few tens based on ref. , while Ibata et al. found two nearby white dwarfs , inferring a density of $`7\times 10^4M_{}/`$pc<sup>3</sup> for these hydrogen atmosphere white dwarfs. A remarkable thing is that this value is just in the required range to account for the missing mass of the heavy spheroid models, whose proper motions would be only $`20`$% smaller than those assumed for ‘dark halo’ objects and hence consistent with those found. To analyse the possibility that these old white dwarfs are genetically related to the old population II spheroid stars, and not to an independent halo population with no observed counterpart, seems then particularly relevant. If this were the case, the initial mass function of spheroid stars would have to be peaked at a few solar masses to account for the large number of white dwarf progenitors, and the gas released during the ejection of their envelopes would have ended up in the disc and bulge, but producing certainly less metal pollution than the halo white dwarf models due to the much smaller total spheroid mass. The future searches of white dwarfs should also be able to distinguish between thick disc and spheroid/halo populations due to their significantly different proper motions.
Regarding the future of microlensing observations, the MACHO experiment has finished taking data, while EROS II and OGLE II are still running. A significant increase in the number of events is then expected when all the data available gets analysed. Furthermore, a new analysis technique (Difference Image Analysis), devised for the study of lensing in crowded fields such as the Andromeda galaxy, has been successfully applied to several bulge fields by the MACHO group , doubling the number of observed events with respect to previous analyses, and also by EROS to their first CCD data . The use of this technique for all the LMC data should then also help to get more decent statistics and hence to discriminate among the population(s) responsible for the microlensing events. Clearly the possibility that the dark halo is completely made of non-baryonic dark matter still remains open, and hence the search for its even more elusive constituents is crucial to finally unravel the dark matter mystery.
## Acknowledgments
This work was supported by CONICET, ANPCyT and Fundación Antorchas, Argentina.
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# 1 Introduction
## 1 Introduction
The topic of Regge trajectories was an active area of research in the 1960’s. But there is a resurgence of interest in Regge trajectories because of the quantity of new data that need analysis and the new quark models need more complete experimental fits for testing. Despite these recent interests, some authors are still using old data to construct Chew-Frautschi plots. In this paper, we reconstruct all Regge trajectories with the most recent data and elucidate the principles of their construction. At the end, we explain why all current meson Regge trajectories models are ruled out by data.
1.1 Theoretical Developments
This paper is concerned with the properties of Regge trajectories which are graphs of the total quantum number $`J`$ versus mass squared $`M^2`$ over a set of particles of fixed principal quantum number $`N`$, isospin $`I`$, dimensionality of the symmetry group $`D`$, spin-parity and flavor. Variations in $`J`$ and $`L`$ $`(J=L+S)`$ are equivalent when $`S`$ is fixed.
Scattering processes are usually analyzed by the method of partial-waves . The wavefunction in the far zone has the form
$$\psi (𝐫)e^{i𝐤𝐫}+f(k,\mathrm{cos}\theta )\frac{e^{i𝐤𝐫}}{r},$$
(1)
where $`\theta `$ is the angle between the wave vector $`𝐤`$ and the position vector $`𝐫`$. In the case of bound states, the plane wave term is absent. The form factor $`f`$ is written as a sum of partial-waves as
$`f(k^2,\mathrm{cos}\theta )`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}(2l+1)a_l(k^2)P_l(\mathrm{cos}\theta ),`$ (2)
$`a_l(k^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _1^1}f(k^2,\mathrm{cos}\theta )P_l(\mathrm{cos}\theta )d\mathrm{cos}\theta ,`$ (3)
where $`l`$ is the orbital angular momentum quantum number and $`P_l`$ is the Legendre polynomial of order $`l`$. In 1959, Regge generalized the solution of $`f`$ by complexifying angular momenta. He interpreted the simple poles of $`a_l(k^2)`$ on the complex $`l`$-plane to be either resonances or bound states.
Chew and Frautschi applied the Regge poles theory to investigate the analyticity of $`a_l(k^2)`$ in the case of strong interactions. They postulated that all strongly interacting particles are self-generating (the bootstrap hypothesis) and that they must lie on Regge trajectories (Chew-Frautschi conjecture) .
At first, linearity was just a convenient guide in constructing the Chew-Frautschi plots because data were scarce and there were few a priori rules to direct the mesons and baryons into the same trajectories . Once linearity was found to be a good working hypothesis, justification was given through certain assumptions in the Regge poles theory as follows: For $`\mathrm{Re}l1/2`$, the partial-wave components of the scattering amplitude $`f`$ have only simple poles and are functions of $`k^2`$,
$$a_l(k^2)\frac{\beta (k^2)}{l\alpha (k^2)},$$
(4)
where $`\beta `$ is the residue (Regge residue) and $`\alpha `$ the position (Regge trajectory) of the simple poles. We can use Watson transformation to rewrite Eq. as the Sommerfeld-Watson formula to include the poles.
A Regge trajectory $`J=\alpha (k^2)`$ is also sometimes expressed as $`J=\alpha (E)`$, or more commonly in terms of the Mandelstam variable $`t`$ as $`J=\alpha (t)`$. $`t`$ is the center-of-mass energy of the quark-antiquark pair defined as $`t(p_q+p_{\overline{q}})^2`$. It is used instead of $`s`$ or $`u`$ because Regge poles generally arise in this channel. For the purpose of plotting, we use $`J=\alpha (M^2)`$. $`\alpha (t)`$ represents a set of leading Regge poles on the complex $`l`$-plane and is called the Reggeon. The condition $`\alpha (t)<0`$ does not correspond to any physical particles because $`J`$ cannot be negative .
Many authors are careful to disclaim linearity as being only approximate. For others, linearity is simply stated . By and large, it is believed that Regge trajectories for relativistic scattering are straight lines over a considerable range of energy without any sign of deviation . Attempts have been made to validate this belief on computational grounds. Kahana, Maung and Norbury calculated the numerical solutions of the relativistic Thompson equation which yield linear, non-intersecting and parallel Regge trajectories. Their calculation did not include the effects of spin which can be a factor in predicting the shape of the Regge trajectories. But earlier in 1985, Godfrey and Isgur solved a relativized Schrödinger equation which did include the spin-spin and spin-orbit interactions. Their calculation concurred with the results of linear Regge trajectories obtained from spinless particles. These two works together seem to suggest that the effect of spin on the shape of the Regge trajectories is negligible. On the other hand, if the coupling constants are not negligible, one expects the spin-orbit contributions to be significant for high $`J`$ values. Salvo *et al.* published solutions for non-linear Regge trajectories by including spin dependent terms in a 3-dimension reduction of the Bethe-Salpeter equation. Salvo’s conclusion differs from those of Kahana and Godfrey concerning the effect of spin on Regge trajectories.
The linearity of Regge trajectories has been the object of investigation once again recently. On the theoretical front, Tang used perturbative QCD to show that Regge trajectories are non-linear by studying high-energy elastic scattering with mesonic exchange in the case of both fixed and running coupling constants. On the experimental side, Brandt *et al.* affirmed the existence of non-linear Pomeron trajectories from the data analysis of the recent UA8 and ISR experiments at CERN. They published a parameterization of Pomeron trajectories containing a quadratic term,
$$\alpha (t)=1.10+0.25t+\alpha ^{\prime \prime }t^2.$$
(5)
where $`\alpha ^{\prime \prime }`$ is a constant. Recently, Burakovsky presented a phenomenological string model for logarithmic and square root Regge trajectories.
In this paper, we check the claims of non-linear Regge trajectories by plotting the most recent experimental data published in the 1998 Review of Particle Physics (RPP) . Our plots confirm the existence of non-linear trajectories. Early Chew-Frautschi plots also show that Regge trajectories fan out. We refer to this non-intersecting property as “divergence.” We also show that many trajectories intersect. Kahana et al. numerically constructed a set of hypothetical Regge trajectories by using a fully relativistic Thompson equation. They discovered that there are differences in the properties of the trajectories obtained by NRSE versus those by the Thompson equation. We summarize the conclusions of Kahana et al. in Table 1 to illustrate these differences.
Table 1: Comparisons of the predictions made by NRSE and the relativistic
Thompson equation according to Kahana et al. . “Yes” refers to a property
predicted by the theory and “No” is the prediction of the opposite property. Non-relativistic Relativistic Schrödinger Equation Thompson Equation Linearity No Yes Divergence No Yes Parallelism Yes Yes
When trajectories of different principal quantum numbers $`N`$ but all other quantum numbers fixed are plotted together, they appear parallel. We call this property “parallelism.”
1.2 Construction of Regge Trajectories
The starting point for constructing a meson Regge trajectory is the meson assignment table in RPP (Table 13.2 on p. 110 of Ref. ). We fix $`I`$ and flavor by selecting particles from a single column. From this column, we isolate different trajectories by fixing $`N`$ and spin-parity when we select particles with consecutive values of $`J`$. For example, the $`1^1S_0`$, $`1^1P_1`$ and $`1^1D_2`$ states constitute an $`N=1`$ singlet trajectory; $`1^3P_0`$ and $`1^3D_1`$ the $`N=1`$ first triplet; $`1^3P_1`$ and $`1^3D_2`$ the $`N=1`$ second triplet; $`1^3S_1`$, $`1^3P_2`$, $`1^3D_3`$, and $`1^3F_4`$ the $`N=1`$ third triplet; $`2^3S_1`$ and $`2^3P_2`$ the $`N=2`$ third triplet and so on. We use the experimental error instead of the width to measure the accuracy of the mass of a meson. The width measures the imaginary part of the complex energy while the experimental error indicates the accuracy of the measurement of the mass at the resonance peak. In case the mass of both the neutral and charged mesons are reported, the mass is taken to be the average of the three. For example, the mass of $`\pi (138)`$ is taken to be the average mass of $`\pi ^0(135)`$ and $`\pi ^\pm (140)`$. Similarly the error of mass is also taken to be the average of the errors of the two corresponding masses. This scheme does not pose any serious ambiguity because the masses and the errors of the neutral and charged mesons are usually quite close (the difference in mass is usually $`<1\%`$ and is $`3.5\%`$ in the worst case) and hence do not change our conclusions. The error of mass square, $`dM^2`$, is calculated from the mass $`M`$ and its error $`dM`$ by the relation $`dM^2=2MdM`$. The end results are 13 trajectories containing 2 particles each, 4 containing 3 particles each and 4 containing 4 particles each. Single particle trajectories are omitted from the plots. None of the $`N=1`$ second triplet trajectories are plotted because most of them are single particle trajectories except the one containing $`K_{1B}`$ and $`K_2(1820)`$, where $`K_{1B}`$ is a nearly equal $`(45^{})`$ mixture of $`K_1(1270)`$ and $`K_1(1400)`$. Since some of these trajectories contained unconfirmed mesons, not all of them are used in this paper.
The bold face entries in the assignment table refer to the mesons which are confirmed by experiments. The regular typeface entries correspond to those which are omitted from the summary table because of work in progress. For example, one of the regular typeface entries in the assignment table, $`f_4(2220)`$, is listed as $`f_J(2220)`$ in the summary table because $`J`$ may assume a value of either 2 or 4 depending on the final confirmation by experiments. There are other similar undetermined quantities in the meson data. This paper takes the conservative approach by using only the bold face (confirmed) data contained in the RPP meson assignment table.
The baryon Regge trajectories are constructed from the RPP baryon assignment table (Table 13.4 on p. 112 of Ref. ). Baryons are categorized into 4 different confidence levels according to their likelihood of existence. Confidence level 1 is assigned to the baryons which are deemed the least likely to exist and level 4 the most likely to exist. The baryon assignment table contains only the level 3 and 4 particles. These are the baryons we will analyze in this paper.
The baryon assignment table uses a set of slightly different quantum numbers, such as $`J^P`$, $`(D,L_N^P)`$ and S. As before, $`J`$ is the total angular momentum, $`P`$ the parity, $`L`$ the orbital angular momentum and $`S`$ the spin. The new quantum number $`D`$ is the dimensionality of the symmetry group and has the value of either 56 or 70. These numbers come from the dimensionalities of the irreducible representations of flavor-spin $`SU(6)`$, i.e. $`\mathrm{𝟔}\mathrm{𝟔}\mathrm{𝟔}=\mathrm{𝟓𝟔}_S\mathrm{𝟕𝟎}_M\mathrm{𝟕𝟎}_M\mathrm{𝟐𝟎}_A`$ where the subscript $`S`$ stands for “symmetric”, $`A`$ for “asymmetric” and $`M`$ for “mixed symmetry.” $`N`$ is the “band” which gives the number of excitation quanta. The construction of a baryon Regge trajectory is similar to that of the meson in that all quantum numbers except $`J`$ are fixed along a trajectory. In other words, $`D`$, $`S`$, flavor, strangeness and isospin are constant along a baryon Regge trajectory. Only $`L`$ is allowed to vary. $`N`$ changes with $`L`$ in the same integer steps so that a change in $`N`$ is the same as a change in $`L`$. Hence we can ignore the consideration of $`N`$.
Regge recurrences are separated by 2 units of $`J`$. In the case of mesons, we can plot two trajectories together in some cases because the cross channel forces between them vanish. It is known as the “exchange degeneracy” (EXD) which arises out of the cross channel forces which split $`a(l,k)`$ into even (+) and odd ($``$) signatures as $`a_\pm (l,k)`$. The separation of the even and odd signatures correspond to the two different Regge trajectories. If the cross channel forces vanish (as in the case of mesons), the even and odd signatures coincide and the even and odd trajectories overlap. It means $`\alpha _+=\alpha _{}`$ and $`\beta _+=\beta _{}`$. These are called the EXD conditions. When the EXD conditions apply, the even and odd parity mesons can be plotted along the same trajectories. In the case of baryons, the cross channel forces persist. Therefore the even and odd parity baryons cannot be plotted together in the same trajectories. The EXD criteria enable us to pick out 3 trajectories of 3 baryons each and 2 trajectories of 2 baryons each from the baryon assignment table. These selections are achieved by isolating a column (e.g. the $`N(939)`$$`N(2220)`$ column) and picking all the particles with the same $`D`$, $`S`$ and $`P`$ (e.g. $`N(939)`$, $`N(1680)`$ and $`N(2220)`$). Once a trajectory is picked from the first column, corresponding entries of the following columns also constitute baryon Regge trajectories. The spectroscopic notation for baryons is $`L_{2I,2J}`$. The $`N`$ and $`\mathrm{\Lambda }`$ trajectories are made up of the $`P_{11}`$, $`F_{15}`$, $`H_{19}`$ states and the $`\mathrm{\Delta }`$ trajectory is made up of the $`P_{33}`$, $`F_{37}`$, $`H_{\mathrm{3\hspace{0.17em}11}}`$ states.
## 2 Linearity
Linearity means that all the particles of a Regge trajectory must lie on the straight line $`M^2=\alpha J+\beta `$. In graphical analysis, non-linearity can be detected by simple inspection in only extreme cases. Linearity on the other hand is more difficult to judge. Therefore we devise a method called the “zone test” to facilitate this judgment.
2.1 Zone Test
We test linearity by the “zone test” on Regge trajectories with 3 or more particles. A test zone of an experimental Regge trajectory is defined to be the area bounded by the error bars of the first and the last particles and the straight lines joining them. Figs. 16 illustrate these test zones (regions enclosed by the dotted lines). A zone contains all the possible straight lines crossing the error bars of the first and the last particles. A Regge trajectory can be a straight line if the error bars of all other particles intersect the zone. In most cases, intersections are easily discernible by inspection. If ambiguity ever arises in borderline cases, an exact numerical version of the zone test is used.
Suppose we are given a sequence of $`N`$ mesons and their values of mass square with errors, $`\{M_i^2\pm dM_i^2\}`$. We calculate the equation of the straight line connecting $`M_1^2+dM_1^2`$ and $`M_N^2+dM_N^2`$ and then the equation of the line connecting $`M_1^2dM_1^2`$ and $`M_N^2dM_N^2`$. These two lines define the boundaries of the zone. For each $`J`$, we can calculate the bounds to be intersected by the error bar to qualify as a linear Regge trajectory. For a 3-particle trajectory in which the particles are labelled $`(1,2,3)`$, the lower and upper bounds at $`J=2`$ are calculated as
$`lb(3,2)`$ $`=`$ $`{\displaystyle \frac{(M_1^2dM_1^2)+(M_3^2dM_3^2)}{2}},`$
$`ub(3,2)`$ $`=`$ $`{\displaystyle \frac{(M_1^2+dM_1^2)+(M_3^2+dM_3^2)}{2}},`$ (6)
where $`lb(3,2)`$ stands for the lower bound and $`ub(3,2)`$ the upper bound of particle 2 along a 3-particle trajectory. Similarly, we can calculate the bounds of particles 2 and 3 along a 4-particle trajectory as
$`lb(4,2)`$ $`=`$ $`{\displaystyle \frac{2(M_1^2dM_1^2)+(M_4^2dM_4^2)}{3}},`$
$`ub(4,2)`$ $`=`$ $`{\displaystyle \frac{2(M_1^2+dM_1^2)+(M_4^2+dM_4^2)}{3}}.`$ (7)
$`lb(4,3)`$ $`=`$ $`{\displaystyle \frac{(M_1^2dM_1^2)+2(M_4^2dM_4^2)}{3}},`$
$`ub(4,3)`$ $`=`$ $`{\displaystyle \frac{(M_1^2+dM_1^2)+2(M_4^2+dM_4^2)}{3}}.`$ (8)
We can generalize these results for particle $`i`$ along an $`n`$-particle trajectory as
$`lb(N,i)`$ $`=`$ $`{\displaystyle \frac{(Ni)(M_1^2dM_1^2)+(i1)(M_N^2dM_N^2)}{N1,}}`$
$`ub(N,i)`$ $`=`$ $`{\displaystyle \frac{(Ni)(M_1^2+dM_1^2)+(i1)(M_N^2+dM_N^2)}{N1.}}`$ (9)
This numerical method is useful for checking linearity when simple inspection is inconclusive.
2.2 Conclusions from Zone Test
All of the data points in all of the graphs in this paper are shown with error bars. If the error bars are invisible in the plots, it simply means that the error bars are smaller than the symbols of the associated data points. We use the zone test to check linearity by simple inspection in Figs. 16. At least one of the error bars of the intermediate particles fails to intersect the test zone in all of the figures except Fig. 3.
Figs. 12 illustrate a group of meson Regge trajectories of the $`N=1`$, $`S=0`$ singlet states and varying $`J`$ corresponding to the $`1^1S_0`$, $`1^1P_1`$ and $`1^1D_2`$ states. Both trajectories fail the zone test and are non-linear. The $`\pi `$ trajectory has a decreasing slope.
In Figs. 36, trajectories of the $`N=1`$, $`S=1`$ third triplet states with varying $`J`$ corresponding to the $`1^3S_1`$, $`1^3P_2`$, $`1^3D_3`$, and $`1^3F_4`$ states are plotted. Trajectories in Figs. 4 and 6 fail the zone test by simple inspection. In Fig. 3, the error bars of both $`a_2(1320)`$ and $`\rho _3(1690)`$ appear to intersect the zone at the lower boundary. In this case, the numerical version of the zone test is used. The error bars of both particles must intersect the bounds to support linearity. The bounds in the case of $`a_2(1320)`$ are $`(lb,ub)(4,2)=(1.73,1.78)\mathrm{GeV}^2`$ which intersect the error bar, (1.7358, 1.7390)$`\mathrm{GeV}^2`$. The bounds for $`\rho _3(1690)`$ are $`(lb,ub)(4,3)=(2.87,2.96)\mathrm{GeV}^2`$ which also intersect the error bar, (2.843, 2.876)$`\mathrm{GeV}^2`$. The numerical test supports the existence of a straight line intersecting all the error bars of the particles along this trajectory. The $`\omega `$ trajectory has an increasing slope while both the $`\varphi `$ and $`K^{}`$ trajectories have decreasing slopes.
The zone test for baryon trajectories are illustrated in Figs. 79. The $`N`$ and $`\mathrm{\Delta }`$ trajectories in Figs. 7 and 9 clearly satisfy the zone test by simple inspection. The $`\mathrm{\Lambda }`$ trajectory in Fig. 8 is shown to be non-linear by the numerical zone test. In summary, 6 of 8 trajectories with 3 or more particles each are shown to be non-linear. Polynominal fits of the trajectories are included in the figure captions for reference only.
## 3 Divergence
Divergence seems to be a property of the Regge trajectories in the early Chew-Frautschi plots and is also a prediction of the numerical calculations by Kahana et al . Divergence is defined to be the conjunction of two properties: (1) non-intersection and (2) fanning out.
We check divergence by plotting families of meson Regge trajectories with the same isospin and spin-parity in Figs. 1015. It is observed that non-linear trajectories of similar masses intertwine. In general, Regge trajectories are not evenly saparated in a graph. Some trajectories can be obscured when many of them are plotted over a large mass range on the same graph. We adopt a numeration scheme which allows us to identify the obscured trajectories in separate plots. For example, the group denoted as 1–3 in Fig. 11 is magnified as trajectories 1–3 in Fig. 12. Divergence is clearly violated in Fig. 14 when trajectories intersect. Due to the large error bars, divergence in Fig. 12, the determination of the properties of these meson trajectories is inconclusive.
Although individual meson trajectories do not fan out, it can be seen in Figs. 10, 11, 13 and 15 that groups of them diverge on a global level. We also notice that these groups can be labelled according to mass difference. In general, the mass of the intersecting trajectories does not differ significantly. On the other hand, divergent trajectories have large mass difference. For example, in Fig. 10, the $`\pi `$, $`K`$ and $`\eta `$ trajectories have small mass difference and form a group of intersecting trajectories. The $`D`$ and $`D_S`$ trajectories also form a group with small mass difference. These two groups of trajectories diverge globally. In summary, trajectories of small mass difference do not diverge but those of large mass difference fan out in our plots. In the case of baryon Regge trajectories, there are insufficient data to test divergence. Divergence is shown to be plausible in Figs. 16 and 17.
## 4 Parallelism
Parallelism refers to the property that Regge trajectories of different values of $`N`$ (which are otherwise identical) are parallel. Two trajectories are parallel if the dynamics are similar. There is no a priori reason why parallelism must hold. There are only two $`\varphi `$ trajectories with $`N=1`$ and $`N=2`$ which qualify for this test. Fig. 18 shows that the two trajectories appear to be parallel. However these trajectories consist of only 2 or 3 mesons each. It is not clear how they will behave at $`J>2`$. The error of $`f_2(2010)`$ is also quite large compared to the separation of the two trajectories. In conclusion, the status of parallelism as a candidate for a property of Regge trajectories is still uncertain.
## 5 Conclusion
The linearity of Regge trajectories is clearly violated in Figs. 1, 2, 4 and 6 by simple inspection but is supported by the numerical zone test in Fig. 3. Divergence is not observed on an individual basis. On the other hand, divergence of groups of trajectories of small mass difference is observed on a global level. Due to insufficient data, parallelism is inconclusive.
Currently there are a variety of models predicting both linear and non-linear Regge trajectories. In general, almost all theories agree that meson Regge trajectories are linear in the small $`J`$ limit. Our plots contradict these predictions. Secondly, all non-linear Regge trajectories models predict trajectories with either increasing or decreasing slopes exclusively, but not both . Our plots show that meson Regge trajectories of both kinds exist. Therefore data rule out all the models that predict non-linear meson Regge trajectories with strictly increasing or decreasing slopes. In the end, data rule out all current meson Regge trajectories models because they are faced with at least one of the problems mentioned above.
## 6 Acknowledgment
We thank Dr. Sudha Swaminathan and Prof. Dale Snider for their comments. This work was supported in part by NASA Research Grant Numbers NCC-1-354 and NCC-1-260.
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# Pulsar optical observation with the Very Large Telescope
## 1. Optical observation of PSR1706-44 & pulsar radiation models
Optical detection of pulsars aid in the development and constraining of theoretical models of pulsar electromagnetic radiation. PSR1706-44 belongs to the set of seven $`\gamma `$-ray pulsars detected by EGRET (Thompson et al. 1996). In the optical it has not been detected yet. Detection of the optical counterparts of radio and gamma-ray pulsars is often complicated by field crowding by other stars in the optical bands. The field of this pulsar was observed with the VLT-UT1 on August 19, 1998 (SV phase). The details of observation, analysis and theoretical implications are given in Lundqvist et al (1999). Chakrabarty and Kaspi (1998) (CK98) using the radio position of the pulsar as summarised in their paper estimate that the pulsar should lie $`2.^{\prime \prime }7`$ away from the star 1. The combined error in position of the optical counterpart of the radio pulsar from various sources is $`1.^{\prime \prime }0`$. Measuring the background at a distance of $`2.^{\prime \prime }7`$ from nearby bright Star 1 in the VLT data shows that an artificial star with $`V=25.5`$ can be detected at more than 3$`\sigma `$ level. To estimate how bright a star one could hide in the PSF of Star 1 we subtracted artificial stars from this position until a hole appeared in the background. It is possible to hide a point source ($`V=25.0`$) at a distance of $`2^{\prime \prime }`$ from the Star 1. As an upper limit for a pulsar this close to Star 1 we claim $`V=24.5`$.
The optical fluxes may be correlated with the gamma-ray photon fluxes in outer gap models (see e.g. Usov (1994) and Cheng, Ho & Ruderman (1986a, 1986b)) for Vela-like pulsars. Assuming $`F_\nu `$ to be the same in both the $`V`$ and $`R`$ bands, the magnitudes predicted by the outer gap models are $`R<20.0`$ and $`V<19.8`$. Our faint $`V`$ limit, in comparison to the predictions of the standard outer-gap model, scaled from gamma-ray flux ($`V<19.8`$), requires a low frequency cutoff in its synchrotron emission spectrum. If the magnetic axis inclination with respect to spin axis $`\chi >\pi /4`$, the synchrotron cutoff frequency is $`10^{13}`$ Hz and in that case the flux of optical radiation may be very small.
## References
Chakrabarty, D., Kaspi, V.M. 1998, ApJ, 498, L37 (CK98)
Cheng, K.S., Ho, C., Ruderman, M.A. 1986, ApJ, 300, 500
Lundqvist, P., Sollerman, J., Ray, A., Leibundgut, B. & Sutaria, F., 1999, A & A, 343, L15
Ray, A., Harding, A.K., Strickman, M.S. 1999, ApJ, 513, 919
Thompson, D.J. et al. 1996, ApJ, 385, 465
Thompson, D.J. et al. 1999, ApJ, in press (astro-ph/9811219)
Usov, V.V. 1994, ApJ 427, 394
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# GEOMETRY OF MULTISYMPLECTIC HAMILTONIAN FIRST-ORDER FIELD THEORIES
## 1 Introduction
The application of techniques of differential geometry to the study of physical theories has been revealed as a very suitable method for better understanding many features of these theories. In particular, the geometric description of classical Field Theories is an area of increasing interest.
The standard geometrical techniques used for the covariant Lagrangian description of first-order Field Theories, involve first order jet bundles $`J^1E\stackrel{\pi ^1}{}E\stackrel{\pi }{}M`$ and their canonical structures (see, for instance, , and references quoted therein). Nevertheless, for the covariant Hamiltonian formalism of these theories the situation is rather different, and there are different kinds of geometrical descriptions for this formalism. For instance, we can find models such as those described in , and , which use $`k`$-symplectic forms, or in , , and , where the essential geometric structure are the $`k`$-cosymplectic forms, or also as in , and , where use is made of polysymplectic forms (in fact, $`k`$-symplectic, $`k`$-cosymplectic and polysymplectic structures are essentially equivalent objects). In this work, we consider only the multisymplectic models , , , , , , , and depending on the choice of the multimomentum phase space there are different ones. In fact:
1. There are some models where the multimomentum phase space is taken to be $`\pi \mathrm{\Lambda }_1^m\mathrm{T}^{}E`$, the bundle of $`m`$-forms on $`E`$ ($`m`$ being the dimension of $`M`$) vanishing by the action of two $`\pi `$-vertical vector fields. This choice is made in works such as , and , as a refinement of the techniques previously given in , and (see also and ).
2. The multimomentum phase space $`J^1\pi ^{}\mathrm{\Lambda }_1^m\mathrm{T}^{}E/\mathrm{\Lambda }_0^m\mathrm{T}^{}E`$ (where $`\mathrm{\Lambda }_0^m\mathrm{T}^{}E`$ is the bundle of $`\pi `$-semibasic $`m`$-forms in $`E`$) has been studied in and used, later on, in , and for the analysis of different aspects of Hamiltonian Field Theories.
3. Finally, in , , , , , and the basic choice is the bundle $`\mathrm{\Pi }\pi ^{}\mathrm{T}M\mathrm{V}^{}(\pi )\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M`$ (here $`\mathrm{V}^{}(\pi )`$ denotes the dual bundle of the $`\pi `$-vertical subbundle $`\mathrm{V}(\pi )`$ of $`\mathrm{T}E`$) which, in turn, is canonically related to $`J^1E^{}\pi ^{}\mathrm{T}M\mathrm{T}^{}E\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M`$.
Although in (and later papers by these authors), a covariant Hamiltonian formalism is constructed in $`\pi `$, in most of the works, this multimomentum bundle is not really used in order to establish a Hamiltonian formalism on $`\pi `$, but just for defining canonical differential structures which, translated to $`J^1E`$ and $`J^1\pi ^{}`$, are used for setting the Lagrangian and Hamiltonian formalisms, respectively. The choice of $`J^1\pi ^{}`$ or $`\mathrm{\Pi }`$ as multimomentum phase space allows us to state covariant Hamiltonian formalisms for Field Theories. Nevertheless, none of them have canonical structures, so the Hamiltonian forms of the Hamiltonian formalism must be obtained from the canonical forms of the multicotangent bundle $`\mathrm{\Lambda }^m\mathrm{T}^{}E`$. This is done by using sections of the projection $`\pi J^1\pi ^{}`$, (or $`J^1E^{}\mathrm{\Pi }`$) which are called Hamiltonian sections , or the so-called Hamiltonian densities , . To our knowledge, a rigorous analysis comparing these formulations and their equivalence has not been done. The aim of this work is to carry out a comparative study of some of these Hamiltonian formulations, establishing the equivalence between them. In every case, the geometrical structures needed for setting the field equations in the Hamiltonian formalism are introduced, as well as the corresponding Legendre maps when the multimomentum bundles are related to a Lagrangian system.
The question of whether the use of connections in the bundle $`\pi :EM`$ is needed for the construction of the covariant formalisms in Field theories is studied. It was analized for the first time in , where a connection was used to define Hamiltonian densities in the Hamiltonian formalism, and in for the case of the density of Lagrangian energy in the Lagrangian formalism. In this work we make a deeper analysis on the role played by connections in the construction of Hamiltonian systems.
An obvious subject of interest is the statement of the Hamiltonian field equations. In all the multisymplectic models field equations are obtained by characterizing the critical sections which are solutions of the problem by means of the multisymplectic form , , , . This characterization can be derived from a variational principle: the so-called Hamilton principle in the Lagrangian formalism and Hamilton-Jacobi principle in the Hamiltonian one. Nevertheless, this aspect of the theory is overlooked in many papers. We give an accurate derivation of the Hamiltonian equations starting from the Hamilton-Jacobi principle, and the role played by connections in the statement of covariant Hamiltonian equations is discussed.
An important kind of Hamiltonian systems are those which are the Hamiltonian counterpart of Lagrangian systems. The construction of such systems starting from the Lagrangian formalism is carried out by using a Legendre map associated with the Lagrangian density and the corresponding multimomentum bundle. This problem has been studied by different authors in the (hyper) regular case (see, for instance, , ), and in the singular (almost-regular) case , , . In this work we review some of these constructions, developing new methods, and giving a unified perspective of all of them.
The structure of the work is as follows:
Section 2 is devoted to a review of the main features of the Lagrangian formalism of Field theories, and afterwards the definition of the different multimomentum bundles for the Hamiltonian formalism, as well as the construction and characterization of the canonical forms with which some of them are endowed. Furthermore, when these multimomentum bundles are related with a Lagrangian system, the corresponding Legendre maps are introduced for both the (hyper)-regular and the almost-regular cases.
In section 3, we undertake the construction of Hamiltonian systems in the multimomentum bundle $`\mathrm{\Pi }`$. As a first step, we will define the Hamiltonian forms which allow us to set the field equations in an intrinsic way. Since $`\mathrm{\Pi }`$ has no canonical geometric form, we must use the canonical forms with which $`\pi `$ and $`J^1E^{}`$ are endowed. Ways of constructing Hamiltonian systems are studied and compared, and in this multimomentum bundle we make a careful deduction of the Hamiltonian equations from the variational principle. In addition, the Hamiltonian formalism associated to a Lagrangian system is developed, both for the hyper-regular and almost-regular cases. Finally, the equivalence between the Lagrangian and Hamiltonian formalisms is proved (for the hyper-regular case).
The construction of Hamiltonian systems in the multimomentum bundle $`J^1\pi ^{}`$ is stated and analyzed in Section 4, following the same pattern as in the above section, and proving the equivalence between the formalisms developed for both multimomentum bundles.
As typical examples, time-dependent mechanics and the electromagnetic field are analyzed (in this context) in Section 5.
An appendix describing the basic geometrical structures in first-order jet bundles is included.
All manifolds are real, paracompact, connected and $`C^{\mathrm{}}`$. All maps are $`C^{\mathrm{}}`$. Sum over crossed repeated indices is understood. Throughout this paper $`\pi :EM`$ will be a fiber bundle ($`dimM=m`$, $`dimE=N+m`$), where $`M`$ is an oriented manifold with volume form $`\omega \mathrm{\Omega }^m(M)`$, and $`\pi ^1:J^1EE`$ will be the jet bundle of local sections of $`\pi `$. The map $`\overline{\pi }^1=\pi \pi ^1:J^1EM`$ defines another structure of differentiable bundle. We denote by $`\mathrm{V}(\overline{\pi }^1)`$ the vertical bundle associated with $`\overline{\pi }^1`$, that is, $`\mathrm{V}(\overline{\pi }^1)=Ker\mathrm{T}\overline{\pi }^1`$, and by $`\text{X}^{\mathrm{V}(\overline{\pi }^1)}(J^1E)`$ the corresponding sections or vertical vector fields. Finally, $`(x^\nu ,y^A,v_\nu ^A)`$ (with $`\nu =1,\mathrm{},m`$; $`A=1,\mathrm{},N`$) will be natural local systems of coordinates in $`J^1E`$ adapted to the bundle $`\pi :EM`$, and such that $`\omega =\mathrm{d}x^1\mathrm{}\mathrm{d}x^m\mathrm{d}^mx`$.
## 2 Geometrical background of the Lagrangian and Hamiltonian formalisms
### 2.1 Lagrangian systems
From the Lagrangian point of view, a first-order classical Field Theory is described by its configuration bundle $`\pi :EM`$, and a Lagrangian density which is a $`\overline{\pi }^1`$-semibasic $`m`$-form on $`J^1E`$ (see the appendix for notation and terminology). A Lagrangian density is usually written as $`=\mathrm{\pounds }(\overline{\pi }^1^{}\omega )`$, where $`\mathrm{\pounds }\mathrm{C}^{\mathrm{}}(J^1E)`$ is the Lagrangian function associated with $``$ and $`\omega `$. The Poincaré-Cartan $`m`$ and $`(m+1)`$-forms associated with the Lagrangian density $``$ are defined using the vertical endomorphism $`𝒱`$ of the bundle $`J^1E`$:
$$\mathrm{\Theta }_{}:=𝑖(𝒱)+\theta _{}+\mathrm{\Omega }^m(J^1E);\mathrm{\Omega }_{}:=\mathrm{d}\mathrm{\Theta }_{}\mathrm{\Omega }^{m+1}(J^1E)$$
In a natural chart in $`J^1E`$ we have
$`\mathrm{\Theta }_{}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\pounds }}{v_\nu ^A}}\mathrm{d}y^A\mathrm{d}^{m1}x_\nu \left({\displaystyle \frac{\mathrm{\pounds }}{v_\nu ^A}}v_\nu ^A\mathrm{\pounds }\right)\mathrm{d}^mx`$
$`\mathrm{\Omega }_{}`$ $`=`$ $`{\displaystyle \frac{^2\mathrm{\pounds }}{v_\eta ^Bv_\nu ^A}}\mathrm{d}v_\eta ^B\mathrm{d}y^A\mathrm{d}^{m1}x_\nu {\displaystyle \frac{^2\mathrm{\pounds }}{y^Bv_\nu ^A}}\mathrm{d}y^B\mathrm{d}y^A\mathrm{d}^{m1}x_\nu +`$
$`{\displaystyle \frac{^2\mathrm{\pounds }}{v_\eta ^Bv_\nu ^A}}v_\nu ^A\mathrm{d}v_\eta ^B\mathrm{d}^mx+\left({\displaystyle \frac{^2\mathrm{\pounds }}{y^Bv_\nu ^A}}v_\nu ^A{\displaystyle \frac{\mathrm{\pounds }}{y^B}}+{\displaystyle \frac{^2\mathrm{\pounds }}{x^\nu v_\nu ^B}}\right)\mathrm{d}y^B\mathrm{d}^mx`$
(See, for instance, , , , , and , for details). Then a Lagrangian system is a couple $`(J^1E,\mathrm{\Omega }_{})`$.
As we can see, the factor $`\mathrm{E}_\mathrm{\pounds }{\displaystyle \frac{\mathrm{\pounds }}{v_\nu ^A}}v_\nu ^A\mathrm{\pounds }`$ appears in the local expression of the Poincaré-Cartan $`(m+1)`$-form, and it is recognized as the classical expression of the Lagrangian energy associated with the Lagrangian function $`\mathrm{\pounds }`$. In fact, the existence of such a function as a global object, and by extension a density of Lagrangian energy, is closely related to the existence of a connection in the bundle $`\pi :EM`$, in the same way that happens in non-autonomous mechanics . As shown in , we can define the density of Lagrangian energy using the vertical endomorphisms in $`J^1E`$. In fact, given a connection $``$ in $`\pi :EM`$, we can identify $`\mathrm{V}^{}(\pi )`$ as a subbundle of $`\mathrm{T}^{}E`$. Then the operation $`𝒮^{}𝒱`$ makes sense, where $`𝒮`$ and $`𝒱`$ are the vertical endomorphisms of the bundle $`J^1E`$, and $`𝒮^{}`$ denotes the action of $`𝒮`$ followed by the injection of $`\mathrm{V}^{}(\pi )`$ in $`\mathrm{T}^{}E`$ induced by $``$ (see the appendix). Therefore:
###### Definition 1
Let $`(J^1E,\mathrm{\Omega }_{})`$ be a Lagrangian system and $``$ a connection in the bundle $`\pi :EM`$. The density of Lagrangian energy associated with the Lagrangian density $``$ and the connection $``$ is given by
$$_{}^{}=𝑖(𝒮^{}𝒱)\mathrm{d}=𝑖(𝒮^{})\mathrm{d}\mathrm{\Theta }_{}\mathrm{\Theta }_{}^{}\mathrm{\Theta }_{}$$
It is a $`\overline{\pi }^1`$-vertical $`m`$-form in $`J^1E`$. Hence, we can write $`_{}^{}=\mathrm{E}_{}^{}(\overline{\pi }^1\omega )`$, where $`\mathrm{E}_{}^{}\mathrm{C}^{\mathrm{}}(J^1E)`$ is the Lagrangian energy function associated with $``$, $``$ and $`\omega `$.
Remark:
* Note that every connection $``$ in $`\pi :EM`$ allows us to split the Poincaré-Cartan forms as
$$\mathrm{\Theta }_{}=\mathrm{\Theta }_{}^{}_{}^{},\mathrm{\Omega }_{}=\mathrm{d}\mathrm{\Theta }_{}^{}+\mathrm{d}_{}^{}\mathrm{\Omega }_{}^{}+\mathrm{d}_{}^{}$$
Using natural systems of coordinates, and $`\mathrm{\Gamma }_\nu ^A`$ being the component functions of the connection, we have the following local expressions
$$_{}^{}=(\frac{\mathrm{\pounds }}{v_\nu ^A}(v_\nu ^A\mathrm{\Gamma }_\nu ^A)\mathrm{\pounds })\mathrm{d}^mx;\mathrm{E}_{}^{}=\frac{\mathrm{\pounds }}{v_\nu ^A}(v_\nu ^A\mathrm{\Gamma }_\nu ^A)\mathrm{\pounds }$$
Observe also that if we take a local connection with $`\mathrm{\Gamma }_\nu ^A=0`$, then the Lagrangian energy associated with this natural connection has the classical local expression given above.
A variational problem can be posed from the Lagrangian density $``$, which is called the Hamilton principle of the Lagrangian formalism: the states of the field are the sections of $`\pi `$ which are critical for the functional $`𝐋:\mathrm{\Gamma }_c(M,E)\text{}`$ defined by
$$𝐋(\varphi ):=_M(j^1\varphi )^{},\text{, for every }\varphi \mathrm{\Gamma }_c(M,E)$$
where $`\mathrm{\Gamma }_c(M,E)`$ is the set of compact supported sections of $`\pi `$. These (compact-supported) critical sections can be characterized in several equivalent ways. In fact (see , , and ):
###### Theorem 1
The critical sections of the Hamilton’s principle are sections $`\varphi :ME`$ whose canonical liftings $`j^1\varphi :MJ^1E`$ satisfy the following equivalent conditions:
1. $`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}|_{t=0}{\displaystyle _M}(j^1\varphi _t)^{}=0`$ , being $`\varphi _t=\sigma _t\varphi `$, where $`\{\sigma _t\}`$ denotes a local one-parameter group of any $`\pi `$-vertical vector field $`Z\text{X}(E)`$.
2. $`{\displaystyle _M}(j^1\varphi )^{}L(j^1Z)=0`$ , for every $`Z\text{X}^{\mathrm{V}(\pi )}(E)`$.
3. $`{\displaystyle _M}(j^1\varphi )^{}L(j^1Z)\mathrm{\Theta }_{}=0`$ , for every $`Z\text{X}^{\mathrm{V}(\pi )}(E)`$.
4. $`(j^1\varphi )^{}𝑖(j^1Z)\mathrm{\Omega }_{}=0`$, for every $`Z\text{X}^{\mathrm{V}(\pi )}(E)`$.
5. $`(j^1\varphi )^{}𝑖(X)\mathrm{\Omega }_{}=0`$ , for every $`X\text{X}(J^1E)`$.
6. The coordinates of $`\varphi `$ satisfy the Euler-Lagrange equations: $`{\displaystyle \frac{\mathrm{\pounds }}{y^A}}|_{j^1\varphi }{\displaystyle \frac{}{x^\nu }}{\displaystyle \frac{\mathrm{\pounds }}{v_\nu ^A}}|_{j^1\varphi }=0`$ (for $`A=1,\mathrm{},N`$).
### 2.2 Multimomentum bundles and Legendre maps
(See for a more detailed study of all these constructions).
Let $`\overline{y}J^1E`$, with $`\overline{y}\stackrel{\pi ^1}{}y\stackrel{\pi }{}x`$. We have that $`\mathrm{T}_{\overline{y}}J_y^1E=\mathrm{V}_{\overline{y}}\pi ^1`$ is canonically isomorphic to $`\mathrm{T}_x^{}M\mathrm{V}_y\pi `$, by means of the directional derivatives; therefore $`\mathrm{V}(\pi ^1)=\overline{\pi }^1\mathrm{T}^{}M_{J^1E}\pi ^1\mathrm{V}(\pi )`$. Moreover, if $`𝒟\mathrm{T}J^1E`$ denotes the subbundle of total derivatives (which in a system of natural coordinates in $`J^1E`$, is generated by $`\left\{{\displaystyle \frac{}{x^\nu }}+v_\nu ^A{\displaystyle \frac{}{y^A}}\right\}`$ ), we have that $`\pi ^1\mathrm{T}E=\pi ^1\mathrm{V}(\pi )𝒟`$ with $`\mathrm{T}_yE|_{\overline{y}}=\mathrm{V}_y\pi |_{\overline{y}}𝒟_{\overline{y}}`$ (see for details). Hence there is a natural projection $`\sigma :\pi ^1\mathrm{T}E\pi ^1\mathrm{V}(\pi )`$ and its dual injection $`\sigma ^{}:\pi ^1\mathrm{V}^{}(\pi )\pi ^1\mathrm{T}^{}E`$ and so we can consider the projection
$$\mathrm{Id}\sigma :\overline{\pi }^1\mathrm{T}^{}M\pi ^1\mathrm{T}E\overline{\pi }^1\mathrm{T}^{}M\pi ^1\mathrm{V}(\pi )=\mathrm{V}(\pi ^1)$$
In a natural chart $`(x^\nu ,y^A)`$ adapted to the bundle $`\pi :EM`$, the local expression of this mapping is
$$\sigma \left(\left(\frac{}{x^\nu }\right)_y|_{\overline{y}}\right)=v_\nu ^A(\overline{y})\left(\frac{}{y^A}\right)_y|_{\overline{y}};\sigma \left(\left(\frac{}{y^A}\right)_y|_{\overline{y}}\right)=\left(\frac{}{y^A}\right)_y|_{\overline{y}}$$
and, if $`\{\zeta ^A\}`$ is the dual basis of $`\left\{{\displaystyle \frac{}{y^A}}\right\}`$ in $`\mathrm{V}^{}(\pi )`$, we have that $`\sigma ^{}(\zeta ^A)=\mathrm{d}y^Av_\nu ^A\mathrm{d}x^\nu `$.
Now let $`(J^1E,\mathrm{\Omega }_{})`$ be a Lagrangian system, and consider the restriction $`_y:J_y^1E\mathrm{\Lambda }^m\mathrm{T}_x^{}M`$. Its differential map at $`\overline{y}J_y^1E`$ is
$$D_{\overline{y}}_y:\mathrm{T}_{\overline{y}}J_y^1E\mathrm{T}_{_y(\overline{y})}\mathrm{\Lambda }^m\mathrm{T}_x^{}M$$
(which, bearing in mind that $`\mathrm{\Lambda }^m\mathrm{T}_x^{}M`$ is a vector space, it is just the vertical differential of $``$). Thus, using the defined projection $`\sigma `$, we have
$$\begin{array}{cccc}\text{}& \text{}& \text{}& \text{}\end{array}$$
###### Definition 2
1. The bundle (over $`E`$)
$$J^1E^{}:=\pi ^{}\mathrm{T}M_E\mathrm{T}^{}E_E\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M$$
is called the generalized multimomentum bundle associated with the bundle $`\pi :EM`$. We denote the natural projections by $`\widehat{\rho }^1:J^1E^{}E`$ and $`\overline{\widehat{\rho }}^1:=\pi \widehat{\rho }^1:J^1E^{}M`$.
2. The generalized Legendre map associated with a Lagrangian density $``$ is the $`\mathrm{C}^{\mathrm{}}`$-map
$$\begin{array}{ccccc}\widehat{\mathrm{F}}& :& J^1E& & J^1E^{}\\ & & \overline{y}& & D_{\overline{y}}_y(\mathrm{Id}\sigma )_{\overline{y}}\end{array}$$
(We have departed a little from the notation in this definition, because $`\sigma `$ acts on $`\pi ^1\mathrm{T}E`$, and not on $`\mathrm{T}E`$. Then, given $`\overline{y}J^1E`$, the right way consists in taking $`\mathrm{T}_yE`$ and lifting it to $`\overline{y}`$).
Natural coordinates in $`J^1E^{}`$ will be denoted by $`(x^\nu ,y^A,\mathrm{p}_\nu ^\eta ,\mathrm{p}_A^\nu )`$, and for every $`𝐲J^1E^{}`$, with $`𝐲\stackrel{\widehat{\rho }^1}{}y\stackrel{\pi }{}x`$, we have
$$𝐲=\frac{}{x^\nu }|_y(\mathrm{p}_\eta ^\nu (𝐲)\mathrm{d}x^\eta +\mathrm{p}_A^\nu (𝐲)\mathrm{d}y^A)_y\mathrm{d}^mx|_y$$
and the local expression of the generalized Legendre map is
$$\widehat{\mathrm{F}}^{}x^\nu =x^\nu ,\widehat{\mathrm{F}}^{}y^A=y^A,\widehat{\mathrm{F}}^{}\mathrm{p}_\eta ^\nu =v_\eta ^A\frac{\mathrm{\pounds }}{v_\nu ^A},\widehat{\mathrm{F}}^{}\mathrm{p}_A^\nu =\frac{\mathrm{\pounds }}{v_\nu ^A}$$
(1)
Now, let $`\overline{y}J^1E`$ with $`\overline{y}\stackrel{\pi ^1}{}y\stackrel{\pi }{}x`$ . We define the map $`_y:J_y^1E\mathrm{\Lambda }^m\mathrm{T}_x^{}M`$ as $`_y:=|_{J_y^1E}`$ . It is a $`\mathrm{C}^{\mathrm{}}`$-map of the affine space $`J_y^1E`$, modeled on $`\mathrm{T}_x^{}M\mathrm{V}_y(\pi )`$, with values on $`\mathrm{\Lambda }^m\mathrm{T}_x^{}M`$. Then, the tangent map $`\mathrm{T}_{\overline{y}}_y`$ allows us to construct the following diagram (where the vertical arrows are canonical isomorphisms given by the directional derivatives)
$$\begin{array}{ccc}\mathrm{T}_{\overline{y}}J_y^1E& & \mathrm{T}_{_y(\overline{y})}\mathrm{\Lambda }^m\mathrm{T}_x^{}M\\ & \mathrm{T}_{\overline{y}}_y& \\ & & \\ & \stackrel{~}{\mathrm{T}}_{\overline{y}}_y& \\ \mathrm{T}_x^{}M\mathrm{V}_y\pi & & \mathrm{\Lambda }^m\mathrm{T}_x^{}M\end{array}$$
Hence, taking into account these identifications, we have that $`\stackrel{~}{\mathrm{T}}_{\overline{y}}_y`$ is an element of $`\mathrm{T}_xM\mathrm{V}_y^{}(\pi )\mathrm{\Lambda }^m\mathrm{T}_x^{}M`$, and so, bearing in mind the analogy with classical mechanics, we define:
###### Definition 3
1. The bundle (over $`E`$)
$$\mathrm{\Pi }:=\pi ^{}\mathrm{T}M_E\mathrm{V}^{}(\pi )_E\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M$$
is called the reduced multimomentum bundle associated with the bundle $`\pi :EM`$. We denote the natural projections by $`\rho ^1:\mathrm{\Pi }E`$ and $`\overline{\rho }^1:=\pi \rho ^1:\mathrm{\Pi }M`$.
2. The reduced Legendre map associated with a Lagrangian density $``$ is the $`\mathrm{C}^{\mathrm{}}`$-map
$$\begin{array}{ccccc}\mathrm{F}& :& J^1E& & \mathrm{\Pi }\\ & & \overline{y}& & \stackrel{~}{\mathrm{T}}_{\overline{y}}_y\end{array}$$
Natural coordinates in $`\mathrm{\Pi }`$ are denoted by $`(x^\nu ,y^A,\mathrm{p}_A^\nu )`$, and for every $`\stackrel{~}{y}\mathrm{\Pi }`$ with $`\stackrel{~}{y}\stackrel{\rho ^1}{}y\stackrel{\pi }{}x`$,
$$\stackrel{~}{y}=\mathrm{p}_A^\nu (\stackrel{~}{y})\frac{}{x^\nu }\zeta ^A\mathrm{d}^mx|_y$$
(We have departed from the notation by denoting the momentum coordinates in $`\mathrm{\Pi }`$ and $`J^1E^{}`$ with the same symbol, $`\mathrm{p}_A^\nu `$,. This departure will be repeated frequently throughout the work).
The local expression of the reduced Legendre map is
$$\mathrm{F}^{}x^\nu =x^\nu ,\mathrm{F}^{}y^A=y^A,\mathrm{F}^{}\mathrm{p}_A^\nu =\frac{\mathrm{\pounds }}{v_\nu ^A}$$
(2)
If we recall that $`J^1E^{}:=\pi ^{}\mathrm{T}M\mathrm{T}^{}E\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M`$, then the natural projection $`\mathrm{T}^{}E\mathrm{V}^{}(\pi )`$ (which is the transpose of the natural injection $`\mathrm{V}(\pi )\mathrm{T}E`$) induces another one
$$\delta :J^1E^{}\mathrm{\Pi }$$
###### Proposition 1
The natural map $`\delta `$ is onto, and $`\mathrm{F}=\delta \widehat{\mathrm{F}}`$.
Furthermore, we can introduce the following map:
###### Definition 4
The canonical contraction in $`J^1E^{}`$ is the map
$$\iota :J^1E^{}\pi ^{}\mathrm{T}M\mathrm{T}^{}E\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M\mathrm{\Lambda }^m\mathrm{T}^{}E$$
defined as follows: $`\iota (𝐲):=\alpha ^k\pi ^{}𝑖(u_k)\beta `$, for every $`𝐲=u_k\alpha ^k\beta J^1E^{}`$.
In a chart of natural coordinates in $`J^1E^{}`$, we have that
$$\iota (𝐲)=(\mathrm{p}_\eta ^\nu (𝐲)\mathrm{d}x^\eta +\mathrm{p}_A^\nu (𝐲)\mathrm{d}y^A)_y𝑖\left(\frac{}{x^\nu }\right)\mathrm{d}^mx|_y=(\mathrm{p}_\nu ^\nu (𝐲)\mathrm{d}^mx+\mathrm{p}_A^\nu (𝐲)\mathrm{d}y^A\mathrm{d}^{m1}x_\nu )_y$$
(3)
(let us recall that $`\mathrm{p}_\nu ^\nu `$ denotes $`{\displaystyle \underset{\nu =1}{\overset{m}{}}}\mathrm{p}_\nu ^\nu `$ ).
For every $`yE`$, we have that
$$\iota (J^1E^{})_y=\{\gamma \mathrm{\Lambda }^m\mathrm{T}_y^{}E;𝑖(u_1)𝑖(u_2)\gamma =0,u_1,u_2\mathrm{V}_y(\pi )\}\mathrm{\Lambda }_1^m\mathrm{T}_y^{}E$$
We will denote $`\iota _0:J^1E^{}\iota (J^1E^{})=\mathrm{\Lambda }_1^m\mathrm{T}^{}E=_{yE}\{(y,\alpha );\alpha \mathrm{\Lambda }_1^m\mathrm{T}_y^{}E\}`$ the restriction of $`\iota `$ onto its image.
###### Definition 5
1. The bundle (over $`E`$)
$$\pi :=\mathrm{\Lambda }_1^m\mathrm{T}^{}E$$
will be called the extended multimomentum bundle associated with the bundle $`\pi :EM`$. We denote the natural projections by $`\widehat{\tau }^1:\pi E`$ and $`\overline{\widehat{\tau }}^1:\pi M`$.
2. The (first) extended Legendre map associated with a Lagrangian density $``$ is the $`\mathrm{C}^{\mathrm{}}`$-map
$$\widehat{}:=\iota _0\widehat{\mathrm{F}}$$
The (second) extended Legendre map is the $`\mathrm{C}^{\mathrm{}}`$-map $`:J^1E\pi `$ given by
$$\stackrel{~}{}=\widehat{}+\pi ^{}$$
Natural coordinates in $`\pi `$ are denoted by $`(x^\nu ,y^A,p,p_A^\nu )`$, and for every $`𝐲J^1E^{}`$ we have
$$\iota :𝐲(x^\nu ,y^A,\mathrm{p}_A^\nu ,\mathrm{p}_\nu ^\eta )\widehat{y}(x^\nu ,y^A,\mathrm{p}_A^\nu ,\mathrm{p}=p_\nu ^\nu )$$
The local expressions of the extended Legendre maps are
$$\begin{array}{ccccccc}\widehat{}^{}x^\nu =x^\nu & ,& \widehat{}^{}y^A=y^A& ,& \widehat{}^{}p_A^\nu =\frac{\mathrm{\pounds }}{v_\nu ^A}& ,& \widehat{}^{}p=v_\nu ^A\frac{\mathrm{\pounds }}{v_\nu ^A}\\ \stackrel{~}{}^{}x^\nu =x^\nu & ,& \stackrel{~}{}^{}y^A=y^A& ,& \stackrel{~}{}^{}p_A^\nu =\frac{\mathrm{\pounds }}{v_\nu ^A}& ,& \stackrel{~}{}^{}p=\mathrm{\pounds }v_\nu ^A\frac{\mathrm{\pounds }}{v_\nu ^A}\end{array}$$
(4)
Remarks:
* It can be proved , that $`\pi \mathrm{\Lambda }_1^m\mathrm{T}^{}E`$ is canonically isomorphic to $`\mathrm{Aff}(J^1E,\mathrm{\Lambda }^m\mathrm{T}^{}M)`$.
* It is interesting to point out that, as $`\mathrm{\Theta }_{}`$ and $`\theta _{}`$ can be thought of as $`m`$-forms on $`J^1E`$ along the projection $`\pi ^1:J^1EE`$, the extended Legendre maps can be defined as
$`(\widehat{}(\overline{y}))(Z_1,\mathrm{},Z_m)`$ $`=`$ $`(\theta _{})_{\overline{y}}(\overline{Z}_1,\mathrm{},\overline{Z}_m)`$
$`(\stackrel{~}{}(\overline{y}))(Z_1,\mathrm{},Z_m)`$ $`=`$ $`(\mathrm{\Theta }_{})_{\overline{y}}(\overline{Z}_1,\mathrm{},\overline{Z}_m)`$
where $`\overline{y}J^1E`$, $`Z_1,\mathrm{},Z_m\mathrm{T}_{\pi ^1(\overline{y})}E`$, and $`\overline{Z}_1,\mathrm{},\overline{Z}_m\mathrm{T}_{\overline{y}}J^1E`$ are such that $`\mathrm{T}_{\overline{y}}\pi ^1\overline{Z}_\nu =Z_\nu `$.
In addition, the (second) extended Legendre map can also be defined as the “first order vertical Taylor approximation to $`\mathrm{\pounds }`$, .
For the construction of the last multimomentum bundle, observe that the sections of the bundle $`\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}ME`$ are the $`\pi `$-semibasic $`m`$-forms on $`E`$; therefore we introduce the notation $`\mathrm{\Lambda }_0^m\mathrm{T}^{}E\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M`$, and then:
###### Definition 6
1. The bundle (over $`E`$)
$$J^1\pi ^{}:=\mathrm{\Lambda }_1^m\mathrm{T}^{}E/\mathrm{\Lambda }_0^m\mathrm{T}^{}E\pi /\mathrm{\Lambda }_0^m\mathrm{T}^{}E$$
will be called the restricted multimomentum bundle associated with the bundle $`\pi :EM`$. We denote the natural projections by $`\tau ^1:J^1\pi ^{}E`$ and $`\overline{\tau }^1:=\pi \tau ^1:J^1\pi ^{}M`$.
2. The restricted Legendre map associated with a Lagrangian density $``$ is the $`\mathrm{C}^{\mathrm{}}`$-map
$$:=\mu \widehat{}=\mu \stackrel{~}{}$$
where $`\mu :\pi J^1\pi ^{}`$ is the natural projection.
Natural coordinates in $`J^1\pi ^{}`$ will also be denoted as $`(x^\nu ,y^A,p_A^\nu )`$. As is evident in this system, the local expression of the restricted Legendre map is
$$^{}x^\nu =x^\nu ,^{}y^A=y^A,^{}p_A^\nu =\frac{\mathrm{\pounds }}{v_\nu ^A}$$
(5)
###### Theorem 2
The multimomentum bundles $`J^1\pi ^{}`$ and $`\mathrm{\Pi }`$ are canonically diffeomorphic as vector bundles over $`E`$, and denoting this diffeomorphism by $`\Psi :J^1\pi ^{}\mathrm{\Pi }`$, therefore $`\mathrm{F}=\Psi `$.
( Proof ) Consider the diagram
$$\begin{array}{ccc}J^1E^{}& \text{}& J^1\pi ^{}\\ \text{}& \text{}& \text{}\\ \mathrm{\Pi }& \text{}& E\end{array}$$
We have that the maps $`\mu \iota _0`$ and $`\delta `$ are sobrejective, linear on the fibers, and restrict to the identity on the base. On the other hand, for every $`yE`$, we have that $`\mathrm{ker}\delta _y=\mathrm{ker}(\mu \iota _0)_y`$ (as can be shown from the corresponding expressions in coordinates). Hence we conclude that $`J^1\pi ^{}`$ and $`\mathrm{\Pi }`$ are canonically isomorphic as vector bundles over $`E`$.
(See for another version of this proof, and an explicit construction of $`\mathrm{\Psi }`$).
$`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$ are fiber bundles over $`E`$, then $`\mathrm{\Psi }`$ is a fiber-diffeomorphism (it is the identity on the base), whose local expression in natural coordinates in $`J^1\pi ^{}`$ and $`\mathrm{\Pi }`$ is
$$\mathrm{\Psi }^{}x^\nu =x^\nu ,\mathrm{\Psi }^{}y_A=y_A,\mathrm{\Psi }^{}\mathrm{p}_A^\nu =p_A^\nu ;(\nu ,A)$$
### 2.3 Canonical forms
As is known , the multicotangent bundle $`\mathrm{\Lambda }^m\mathrm{T}^{}E`$ is endowed with canonical forms: $`𝚯\mathrm{\Omega }^m(\mathrm{\Lambda }^m\mathrm{T}^{}E)`$ and the multisymplectic form $`𝛀:=\mathrm{d}𝚯\mathrm{\Omega }^{m+1}(\mathrm{\Lambda }^m\mathrm{T}^{}E)`$. Then:
###### Definition 7
The canonical $`m`$ and $`(m+1)`$ forms of $`J^1E^{}`$ are
$$\widehat{\mathrm{\Theta }}=\iota ^{}𝚯\mathrm{\Omega }^m(J^1E^{}),\widehat{\mathrm{\Omega }}:=\mathrm{d}\widehat{\mathrm{\Theta }}=\iota ^{}𝛀\mathrm{\Omega }^{m+1}(J^1E^{})$$
On the other hand, observe that $`\pi \mathrm{\Lambda }_1^m\mathrm{T}^{}E`$ is a subbundle of the multicotangent bundle $`\mathrm{\Lambda }^m\mathrm{T}^{}E`$. Let
$$\lambda :\mathrm{\Lambda }_1^m\mathrm{T}^{}E\mathrm{\Lambda }^m\mathrm{T}^{}E$$
be the natural imbedding (hence $`\lambda \iota _0=\iota `$). Then:
###### Definition 8
The canonical $`m`$ and $`(m+1)`$ forms of $`\pi `$ (multimomentum Liouville $`m`$ and $`(m+1)`$ forms of $`\pi `$) are
$$\mathrm{\Theta }:=\lambda ^{}𝚯\mathrm{\Omega }^m(\pi ),\mathrm{\Omega }=\mathrm{d}\mathrm{\Theta }=\lambda ^{}𝛀\mathrm{\Omega }^{m+1}(\pi )$$
* Of course, $`\widehat{\mathrm{\Theta }}=\iota _0^{}\mathrm{\Theta }`$ and $`\widehat{\mathrm{\Omega }}=\iota _0^{}\mathrm{\Omega }`$
* $`\mathrm{\Omega }`$ is $`1`$-nondegenerate, and hence $`(\pi ,\mathrm{\Omega })`$ is a multisymplectic manifold.
The canonical forms $`\widehat{\mathrm{\Theta }}`$ and $`\mathrm{\Theta }`$ can also be characterized as follows (see ):
* $`\widehat{\mathrm{\Theta }}`$ is the only $`m`$-form in $`J^1E^{}`$, such that if $`𝐲J^1E^{}`$, and $`X_1,\mathrm{},X_m\mathrm{T}_𝐲J^1E^{}`$, then
$$\widehat{\mathrm{\Theta }}(𝐲;X_1,\mathrm{},X_m)=\iota (𝐲)[\mathrm{T}_𝐲\widehat{\rho }^1(X_1),\mathrm{},\mathrm{T}_𝐲\widehat{\rho }^1(X_m)]$$
* In turn, considering the natural projection $`\widehat{\kappa }^1:\mathrm{\Lambda }_1^m\mathrm{T}^{}EE`$, then
$$\mathrm{\Theta }((y,\alpha );X_1,\mathrm{},X_m):=\alpha (y;\mathrm{T}_{(y,\alpha )}\widehat{\kappa }^1(X_1),\mathrm{},\mathrm{T}_{(y,\alpha )}\widehat{\kappa }^1(X_m))$$
for every $`(y,\alpha )\mathrm{\Lambda }_1^m\mathrm{T}^{}E`$ (where $`yE`$ and $`\alpha \mathrm{\Lambda }_1^m\mathrm{T}_y^{}E`$), and $`X_i\text{X}(\mathrm{\Lambda }_1^m\mathrm{T}^{}E)`$.
Bearing in mind the following diagram
$$\begin{array}{ccc}& & \pi \\ & \text{}& \text{}\\ J^1E^{}& \text{}& E\end{array}$$
we observe that the map $`\iota _0:J^1E^{}\pi `$ is a form along the projection $`\widehat{\rho }^1:J^1E^{}E`$. Then:
###### Lemma 1
$`\widehat{\mathrm{\Theta }}=\iota _0^{}\mathrm{\Theta }=\widehat{\rho }^1\iota _0`$.
( Proof ) Let $`𝐲J^1E^{}`$, and $`X_1,\mathrm{},X_m\mathrm{T}_𝐲J^1E^{}`$. We have
$`\iota _0^{}\mathrm{\Theta }(𝐲;X_1,\mathrm{},X_m)`$ $`=`$ $`\mathrm{\Theta }(\iota _0(𝐲);\mathrm{T}_𝐲\iota _0(X_1),\mathrm{},\mathrm{T}_𝐲\iota _0(X_m))`$
$`=`$ $`(\iota _0(𝐲)[\mathrm{T}_𝐲(\widehat{\tau }^1\iota _0)(X_1),\mathrm{},\mathrm{T}_𝐲(\widehat{\tau }^1\iota _0)(X_m)]`$
$`=`$ $`(\iota _0(𝐲)(\mathrm{T}_𝐲\widehat{\rho }^1(X_1),\mathrm{},\mathrm{T}_𝐲\rho ^1(X_m))=(\widehat{\rho }^1\iota )(𝐲;X_1,\mathrm{},X_m)`$
In natural coordinates in $`J^1E^{}`$ and $`\pi `$, the local expressions of these forms are
$`\widehat{\mathrm{\Theta }}=\mathrm{p}_\eta ^\eta \mathrm{d}^mx+\mathrm{p}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu `$ $`,`$ $`\widehat{\mathrm{\Omega }}=\mathrm{dp}_\eta ^\eta \mathrm{d}^mx\mathrm{dp}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu `$
$`\mathrm{\Theta }=p\mathrm{d}^mx+p_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu `$ $`,`$ $`\mathrm{\Omega }=\mathrm{d}p\mathrm{d}^mx\mathrm{d}p_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu `$
###### Proposition 2
Let $`(J^1E,\mathrm{\Omega }_{})`$ be a Lagrangian system. Let $`\widehat{\mathrm{F}}`$ be the generalized Legendre map, and $`\widehat{}`$ and $`\stackrel{~}{}`$ the extended Legendre maps. Then
$`\widehat{\mathrm{F}}^{}\widehat{\mathrm{\Theta }}=\mathrm{\Theta }_{}=\theta _{}`$ $`;`$ $`\widehat{\mathrm{F}}^{}\widehat{\mathrm{\Omega }}=\mathrm{\Omega }_{}\mathrm{d}=\mathrm{d}\theta _{}`$
$`\widehat{}^{}\mathrm{\Theta }=\mathrm{\Theta }_{}=\theta _{}`$ $`;`$ $`\widehat{}^{}\mathrm{\Omega }=\mathrm{\Omega }_{}\mathrm{d}=\mathrm{d}\theta _{}`$
$`\stackrel{~}{}^{}\mathrm{\Theta }=\mathrm{\Theta }_{}`$ $`;`$ $`\stackrel{~}{}^{}\mathrm{\Omega }=\mathrm{\Omega }_{}`$
### 2.4 Regular and singular systems
###### Definition 9
Let $`(J^1E,\mathrm{\Omega }_{})`$ be a Lagrangian system.
1. $`(J^1E,\mathrm{\Omega }_{})`$ is said to be a regular or non-degenerate Lagrangian system if $``$, and hence, $`\mathrm{F}`$ are local diffeomorphisms.
As a particular case, $`(J^1E,\mathrm{\Omega }_{})`$ is said to be a hyper-regular Lagrangian system if $``$, and hence $`\mathrm{F}`$, are global diffeomorphisms.
2. Elsewhere $`(J^1E,\mathrm{\Omega }_{})`$ is said to be a singular or degenerate Lagrangian system.
The matrix of the tangent maps $`_{}`$ and $`\mathrm{F}_{}`$ in a natural coordinate system is
$$\left(\begin{array}{ccc}\mathrm{Id}& 0& 0\\ 0& \mathrm{Id}& 0\\ \frac{^2\mathrm{\pounds }}{x^\nu v_\mu ^A}& \frac{^2\mathrm{\pounds }}{y^Bv_\mu ^A}& \frac{^2\mathrm{\pounds }}{v_\nu ^Bv_\mu ^A}\end{array}\right)$$
(6)
where the sub-matrix $`\left({\displaystyle \frac{^2\mathrm{\pounds }}{v_\nu ^Bv_\mu ^A}}\right)`$ is the partial Hessian matrix of $`\mathrm{\pounds }`$. Then, the regularity condition is equivalent to demanding that this matrix is regular everywhere in $`J^1E`$. This fact establishes the relation to the concept of regularity given in an equivalent way by saying that a Lagrangian system $`(J^1E,\mathrm{\Omega }_{})`$ is regular if $`\mathrm{\Omega }_{}`$ is $`1`$-nondegenerate. (See also for a different definition of this concept).
###### Proposition 3
(See and ). Let $`(J^1E,\mathrm{\Omega }_{})`$ be a hyper-regular Lagrangian system. Then
1. $`\widehat{\mathrm{F}}(J^1E)`$ is a $`m^2`$-codimensional imbedded submanifold of $`J^1E^{}`$, which is transverse to the projection $`\delta `$.
2. $`\widehat{}(J^1E)`$ and $`\stackrel{~}{}(J^1E)`$ are 1-codimensional imbedded submanifolds of $`\pi `$, which are transverse to the projection $`\mu `$.
3. The manifolds $`J^1\pi ^{}`$, $`\widehat{}(J^1E)`$, $`\stackrel{~}{}(J^1E)`$, $`\widehat{\mathrm{F}}(J^1E)`$ and $`\mathrm{\Pi }`$ are diffeomorphic.
Hence, $`\widehat{\mathrm{F}}`$, $`\widehat{}`$ and $`\stackrel{~}{}`$ are diffeomorphisms on their images; and the maps $`\mu `$, restricted to $`\widehat{}(J^1E)`$ or to $`\stackrel{~}{}(J^1E)`$, and $`\iota _0`$ and $`\delta `$, restricted to $`\widehat{\mathrm{F}}(J^1E)`$, are also diffeomorphisms.
In this way we have the following diagram
$$\begin{array}{ccccccccc}\text{}& \text{}& \text{}& \text{}& \text{}& \text{}& & & \end{array}$$
(7)
where the map $`\mu ^{}:\pi \pi `$ is defined by the relation
$$\mu ^{}:=\stackrel{~}{}^1\mu $$
and it satisfies that $`\mu \mu ^{}=\mu `$. Observe also that the restriction $`\mu ^{}:\widehat{}(J^1E)\pi \stackrel{~}{}(J^1E)\pi `$, is a diffeomorphism, which is also defined by the relation $`\stackrel{~}{}=\mu ^{}\widehat{}`$.
For dealing with singular Lagrangians, we must assume minimal “regularity” conditions. Hence we introduce the following terminology:
###### Definition 10
A singular Lagrangian system $`(J^1E,\mathrm{\Omega }_{})`$ is said to be almost-regular if:
1. $`𝒫:=(J^1E)`$ and $`P:=\mathrm{F}(J^1E)`$ are closed submanifolds of $`J^1\pi ^{}`$ and $`\mathrm{\Pi }`$, respectively.
(We will denote the corresponding imbeddings by $`ȷ_0:𝒫J^1\pi ^{}`$ and $`ȷ_0:P\mathrm{\Pi }`$).
2. $``$, and hence $`\mathrm{F}`$, are submersions onto their images (with connected fibers).
3. For every $`\overline{y}J^1E`$, the fibers $`^1((\overline{y}))`$ and hence $`\mathrm{F}^1(\mathrm{F}(\overline{y}))`$ are connected submanifolds of $`J^1E`$.
(This definition is equivalent to that in reference , but slightly different from that in references and ).
Let $`(J^1E,\mathrm{\Omega }_{})`$ be an almost-regular Lagrangian system. Denote
$$\widehat{𝒫}:=\widehat{}(J^1E),\stackrel{~}{𝒫}:=\stackrel{~}{}(J^1E),\widehat{P}:=\widehat{\mathrm{F}}(J^1E)$$
Let $`\widehat{ȷ}_0:\widehat{𝒫}\pi `$, $`\stackrel{~}{ȷ}_0:\stackrel{~}{𝒫}\pi `$, $`\widehat{ȷ}_0:\widehat{P}J^1E^{}`$ be the canonical inclusions, and
$$\widehat{\mu }:\widehat{𝒫}𝒫,\stackrel{~}{\mu }:\stackrel{~}{𝒫}𝒫,\widehat{\iota }_0:\widehat{P}\widehat{𝒫},\widehat{\delta }:\widehat{P}P,\mathrm{\Psi }_0:\widehat{𝒫}P$$
the restrictions of the maps $`\mu `$, $`\iota _0`$, $`\delta `$ and the diffeomorphism $`\mathrm{\Psi }`$, respectively. Finally, define the restriction mappings
$$_0:J^1E𝒫,\stackrel{~}{}_0:J^1E\stackrel{~}{𝒫},\widehat{}_0:J^1E\widehat{𝒫},\widehat{\mathrm{F}}_0:J^1E\widehat{P},\mathrm{F}_0:J^1EP$$
###### Proposition 4
(See ,,). Let $`(J^1E,\mathrm{\Omega }_{})`$ be an almost-regular Lagrangian system. Then:
1. The maps $`\mathrm{\Psi }_0`$ and $`\stackrel{~}{\mu }`$ are diffeomorphisms.
2. For every $`\overline{y}J^1E`$,
$$\stackrel{~}{_0}^1(\stackrel{~}{_0}(\overline{y}))=_0^1(_0(\overline{y}))=\mathrm{F}_0^1(\mathrm{F}_0(\overline{y}))$$
(8)
3. $`\stackrel{~}{𝒫}`$ and $`\widehat{𝒫}`$ are submanifolds of $`\pi `$, $`\widehat{P}`$ is a submanifold of $`J^1E^{}`$, and $`\stackrel{~}{ȷ}_0:\stackrel{~}{𝒫}\pi `$, $`\widehat{ȷ}_0:\widehat{𝒫}\pi `$, $`\widehat{ȷ}_0:\widehat{P}J^1E^{}`$ are imbeddings.
4. The restriction mappings $`\stackrel{~}{}_0`$, $`\widehat{}_0`$ and $`\widehat{\mathrm{F}}_0`$ are submersions with connected fibers.
Thus we have the diagram
$$\begin{array}{ccccccc}\text{}& \text{}& \text{}& \text{}& \text{}& \text{}& \text{}\end{array}$$
(9)
where $`\widehat{\mu }^{}:\widehat{𝒫}\stackrel{~}{𝒫}`$ is defined by the relation $`\widehat{\mu }^{}:=\stackrel{~}{\mu }^1\widehat{\mu }`$.
The maps $`\widehat{\mu }`$ and $`\widehat{\mu }^{}`$ are not diffeomorphisms in general, since $`\mathrm{rank}\widehat{}_0\mathrm{rank}\stackrel{~}{}_0=\mathrm{rank}_0`$, as is evident from the analysis of the corresponding Jacobian matrices.
###### Proposition 5
Let $`(J^1E,\mathrm{\Omega }_{})`$ be an almost-regular Lagrangian system. Then:
$$\mathrm{ker}\mathrm{F}_{}=\mathrm{ker}_{}=\stackrel{~}{\mathrm{ker}}_{}=\mathrm{ker}\mathrm{\Omega }_{}\text{X}^{\mathrm{V}(\pi ^1)}(J^1E)$$
( Proof ) The first two equalities are immediate, since $`P:=\mathrm{F}(J^1E)`$, $`𝒫:=(J^1E)`$ and $`\stackrel{~}{P}:=\stackrel{~}{}(J^1E)`$ are diffeomorphic.
For the last equality, as $`\mathrm{F}`$, $``$ and $`\stackrel{~}{}`$ are the identity on the basis $`E`$ of the bundle $`\pi ^1:J^1EE`$, first we have that
$$\mathrm{ker}\mathrm{F}_{}=\mathrm{ker}_{}=\mathrm{ker}\stackrel{~}{}\text{X}^{\mathrm{V}(\pi ^1)}(J^1E)$$
(and this relation holds also for the other Legendre maps). Then, for every $`X\mathrm{ker}\stackrel{~}{}_{}`$ we have
$$𝑖(X)\mathrm{\Omega }_{}=𝑖(X)(\stackrel{~}{}^{}\mathrm{\Omega })=\stackrel{~}{}^{}[𝑖(\stackrel{~}{}_{}X)\mathrm{\Omega }]=0$$
and hence
$$\mathrm{ker}\stackrel{~}{}_{}=\mathrm{ker}_{}=\mathrm{ker}\mathrm{F}_{}\mathrm{ker}\mathrm{\Omega }_{}\text{X}^{\mathrm{V}(\pi ^1)}(J^1E)$$
Conversely, if $`X\text{X}^{\mathrm{V}(\pi ^1)}(J^1E)`$, in a natural system of coordinates in $`J^1E`$ we have $`X=f_\nu ^B{\displaystyle \frac{}{v_\nu ^B}}`$ , and if, in addition, $`X\mathrm{ker}\mathrm{\Omega }_{}`$, we obtain
$$0=𝑖(X)\mathrm{\Omega }_{}=\frac{^2\mathrm{\pounds }}{v_\eta ^Bv_\nu ^A}f_\eta ^B\mathrm{d}y^A\mathrm{d}^{m1}x_\nu +\frac{^2\mathrm{\pounds }}{v_\eta ^Bv_\nu ^A}v_\nu ^Af_\eta ^B\mathrm{d}^mx$$
this is equivalent to demanding that $`f_\eta ^B{\displaystyle \frac{^2\mathrm{\pounds }}{v_\eta ^Bv_\nu ^A}}=0`$ , and this is the condition which characterizes locally the vector fields belonging to $`\mathrm{ker}\mathrm{F}_{}=\mathrm{ker}_{}=\mathrm{ker}\stackrel{~}{}_{}`$ (see (6)).
## 3 Hamiltonian formalism in the reduced multimomentum bundle
### 3.1 Hamiltonian systems
(Compare this presentation with references , and ).
The more standard way for constructing Hamiltonian systems in $`\mathrm{\Pi }`$ consists in using sections of the projection $`\delta `$, and it is similar to that developed in for the Hamiltonian formalism in $`J^1\pi ^{}`$ (which we will review later). Thus:
###### Definition 11
Consider the bundle $`\overline{\rho }^1:\mathrm{\Pi }M`$.
1. A section $`h_\delta :\mathrm{\Pi }J^1E^{}`$ of the projection $`\delta `$ is called a Hamiltonian section of $`\delta `$.
2. The differentiable forms
$$\mathrm{\Theta }_{h_\delta }:=h_\delta ^{}\widehat{\mathrm{\Theta }}=(\iota _0h_\delta )^{}\mathrm{\Theta },\mathrm{\Omega }_{h_\delta }:=\mathrm{d}\mathrm{\Theta }_{h_\delta }=h_\delta ^{}\widehat{\mathrm{\Omega }}=(\iota _0h_\delta )^{}\mathrm{\Omega }$$
are called the Hamilton-Cartan $`m`$ and $`(m+1)`$ forms of $`\mathrm{\Pi }`$ associated with the Hamiltonian section $`h_\delta `$.
3. The couple $`(\mathrm{\Pi },\mathrm{\Omega }_{h_\delta })`$ is said to be a Hamiltonian system.
Using charts of natural coordinates in $`\mathrm{\Pi }`$ and $`J^1E^{}`$, a Hamiltonian section is specified by a set of local functions $`H_\nu ^\eta \mathrm{C}^{\mathrm{}}(U)`$, $`U\mathrm{\Pi }`$, such that
$$h_\delta (x^\nu ,y^A,\mathrm{p}_A^\nu )(x^\nu ,y^A,\mathrm{p}_\nu ^\eta =H_\nu ^\eta (x^\gamma ,y^B,\mathrm{p}_B^\gamma ),\mathrm{p}_A^\nu )$$
(10)
Then, the local expressions of these Hamilton-Cartan forms are
$`\mathrm{\Theta }_{h_\delta }`$ $`=`$ $`\mathrm{p}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu H_{h_\delta }\mathrm{d}^mx`$
$`\mathrm{\Omega }_{h_\delta }`$ $`=`$ $`\mathrm{dp}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu +\mathrm{d}H_{h_\delta }\mathrm{d}^mx`$ (11)
where $`H_{h_\delta }H_\nu ^\nu `$ is a local Hamiltonian function associated with the Hamiltonian section $`h_\delta `$.
As $`\iota _0`$ is a submersion, we can state the following:
###### Definition 12
There is a natural equivalence relation in the set of Hamiltonian sections of $`\delta `$, which is defined as follows: two Hamiltonian sections $`h_\delta ^1,h_\delta ^2`$ are equivalent if $`\iota _0h_\delta ^1=\iota _0h_\delta ^2`$. We denote by $`\{h_\delta \}`$ the equivalence classes of this relation.
Remarks:
* Of course, Hamiltonian sections belonging to the same equivalence class give the same Hamilton-Cartan forms, and hence the same Hamiltonian system.
* Observe that all the Hamiltonian sections of the same equivalence class have the same local Hamiltonian function $`H_{h_\delta }H_\nu ^\nu `$ (in the same open set $`U\mathrm{\Pi }`$).
There is a relation between sections of $`\mu `$ and of $`\delta `$. In fact:
###### Proposition 6
There is a bijective correspondence between the set of sections of the projection $`\mu :\pi J^1\pi ^{}`$ and the set of equivalence classes of sections of the projection $`\delta :J^1E^{}\mathrm{\Pi }`$.
( Proof ) In fact, this correspondence is established by the commutativity of the following diagram
that is, a section $`h_\mu :J^1\pi ^{}\pi `$ and a class $`\{h_\delta \}:\mathrm{\Pi }J^1E^{}`$ are in correspondence if, and only if,
$$h_\mu =\iota _0h_\delta \mathrm{\Psi },\text{for every }h_\delta \{h_\delta \}$$
and this correspondence is one-to-one.
Now we can study the structure of the set of Hamilton-Cartan forms, and hence of Hamiltonian systems. So, for every Hamiltonian section $`h_\delta `$ of $`\delta `$, consider the diagram
$$\begin{array}{ccc}& & \pi \\ & \text{}& \text{}\\ \mathrm{\Pi }& \text{}J^1E^{}\text{}& E\\ & \text{}& \\ & \text{}& \text{}\end{array}$$
###### Lemma 2
Let $`h_\delta ^1,h_\delta ^2:\mathrm{\Pi }J^1E^{}`$ be two Hamiltonian sections of $`\delta `$, then:
1. $`h_\delta ^1\widehat{\mathrm{\Theta }}h_\delta ^2\widehat{\mathrm{\Theta }}=(\iota _0h_\delta ^1)^{}\mathrm{\Theta }(\iota _0h_\delta ^2)^{}\mathrm{\Theta }=\rho ^1(\iota _0h_\delta ^1\iota _0h_\delta ^2)`$.
2. $`\rho ^1(\iota _0h_\delta ^1\iota _0h_\delta ^2)`$ is a $`\overline{\rho }^1`$-semibasic form in $`\mathrm{\Pi }`$.
( Proof )
1. For every Hamiltonian section $`h_\delta `$, the map $`\iota _0h_\delta `$ is a form along the map $`\widehat{\rho }^1h_\delta =\rho ^1`$. Therefore, following the same pattern as in Lemma 1, we obtain that
$$(\iota _0h_\delta )^{}\mathrm{\Theta }=\rho ^1(\iota _0h_\delta )$$
and hence the result is immediate.
2. As $`\mathrm{\Psi }\mu \iota _0=\delta `$, for every section $`h_\delta `$ we have that $`\mathrm{\Psi }\mu \iota _0h_\delta =\mathrm{Id}_\mathrm{\Pi }`$, then $`\mu (\iota _0h_\delta ^1\iota _0h_\delta ^2)=0_\mathrm{\Pi }`$, and therefore $`\mathrm{Im}(\iota _0h_\delta ^1\iota _0h_\delta ^2)\mathrm{ker}\mu =\mathrm{\Lambda }_m^0\mathrm{T}^{}E`$ (that is, the $`\overline{\rho }^1`$-semibasic forms in $`\mathrm{\Pi }`$.)
From the local expressions (3) and (10), for every $`\stackrel{~}{y}U\mathrm{\Pi }`$, we have
$$(\iota _0h_\delta ^1\iota _0h_\delta ^2)(\stackrel{~}{y})=(H_{h_\delta ^1}H_{h_\delta ^2})(\stackrel{~}{y})\mathrm{d}^mx|_{\stackrel{~}{y}}$$
which is the local expression (at $`\stackrel{~}{y}`$) of a $`\overline{\rho }^1`$-semibasic form in $`\mathrm{\Pi }`$. In addition, this is also the local expression of the form
$$(\iota _0h_\delta ^1)^{}\mathrm{\Theta }(\iota _0h_\delta ^2)^{}\mathrm{\Theta }=\mathrm{\Theta }_{h_\delta ^1}\mathrm{\Theta }_{h_\delta ^2}\text{}$$
(12)
###### Definition 13
A $`\overline{\rho }^1`$-semibasic form $`\text{}\Omega ^m(\mathrm{\Pi })`$ is said to be a Hamiltonian density in $`\mathrm{\Pi }`$.
It can be written as $`\text{}=\mathrm{H}(\overline{\rho }^1^{}\omega )`$, where $`\mathrm{H}\mathrm{C}^{\mathrm{}}(\mathrm{\Pi })`$ is the global Hamiltonian function associated with and $`\omega `$.
In this way, we have proved that two Hamiltonian systems generated by two Hamiltonian sections of $`\delta `$ belonging to different equivalence classes, are related by means of a Hamiltonian density. We can state this result as follows:
###### Theorem 3
The set of Hamilton-Cartan $`m`$-forms associated with Hamiltonian sections of $`\delta `$ is an affine space modelled on the set of Hamiltonian densities in $`\mathrm{\Pi }`$.
Remark:
* If $`(\mathrm{\Pi },\mathrm{\Omega }_{h_\delta })`$ is a Hamiltonian system, taking into account (12) we have that every Hamiltonian section $`h_\delta ^{}`$ (such that $`h_\delta ^{}\{h_\delta \}`$) allows us to split globally the Hamilton-Cartan forms as
$$\mathrm{\Theta }_{h_\delta }=\mathrm{\Theta }_{h_\delta ^{}}\text{};\mathrm{\Omega }_{h_\delta }=\mathrm{\Omega }_{h_\delta ^{}}+\mathrm{d}\text{}$$
(13)
If $`(x^\nu ,y^A,\mathrm{p}_A^\nu )`$ is a natural system of coordinates in $`\mathrm{\Pi }`$, such that $`\overline{\rho }^1\omega =\mathrm{d}^mx`$, and $`H_{h_\delta ^{}}(x^\nu ,y^A,\mathrm{p}_A^\nu )`$ is the local Hamiltonian function associated with the Hamiltonian section $`h_\delta ^{}`$, and $`\text{}=\mathrm{H}(x^\nu ,y^A,\mathrm{p}_A^\nu )\mathrm{d}^mx`$, then
$`\mathrm{\Theta }_{h_\delta }`$ $`=`$ $`\mathrm{p}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu (\mathrm{H}+H_{h_\delta ^{}})\mathrm{d}^mx`$
$`\mathrm{\Omega }_{h_\delta }`$ $`=`$ $`\mathrm{dp}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu +\mathrm{d}(\mathrm{H}+H_{h_\delta ^{}})\mathrm{d}^mx`$ (14)
If $`H_{h_\delta }`$ is the local Hamiltonian function associated with the Hamiltonian section $`h_\delta `$, we have the relation $`\mathrm{H}=H_{h_\delta }H_{h_\delta ^{}}`$ (in an open set $`U`$). Hence, taking this into account, the local expressions (11) and (14) are really the same thing.
### 3.2 Hamiltonian sections, Hamiltonian densities and connections
In order to obtain a Hamiltonian density using two Hamiltonian sections, it is usual for one of them to be a linear section induced by a connection. This is a natural procedure for different reasons. For instance, when we construct the Hamiltonian formalism associated with a Lagrangian system, the Hamiltonian density must be related with the density of Lagrangian energy and, as this last is defined by using a connection, this same connection must be used for constructing the related Hamiltonian density (see sections 3.4 and 3.6).
Next, we are going to show how to define the linear Hamiltonian section induced by a connection. Hence, suppose that a connection $``$ has been chosen in $`\pi :EM`$. It allows us to identify $`\mathrm{V}^{}(\pi )`$ as a subbundle of $`\mathrm{T}^{}E`$. So, if $`v_{}^{}:\mathrm{V}^{}(\pi )\mathrm{T}^{}E`$ is the dual injection of the vertical projection $`v_{}`$ induced by $``$. Then:
###### Definition 14
The linear Hamiltonian section of $`\delta `$ induced by the connection $``$ is the map
$$\begin{array}{ccccc}h_\delta ^{}& :& \mathrm{\Pi }& & \pi ^{}\mathrm{T}M\mathrm{T}^{}E\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M:=J^1E^{}\\ & & u_k\alpha ^k\beta & & u_kv_{}^{}(\alpha ^k)\beta \end{array}$$
that is, $`h_\delta ^{}:=\mathrm{Id}_{\pi ^{}\mathrm{T}M}v_{}^{}\mathrm{Id}_{\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M}`$.
Remark:
* Two linear sections $`h_\delta ^_1`$ and $`h_\delta ^_2`$ induced by two different connections $`_1`$ and $`_2`$ cannot belong to the same equivalence class of Hamiltonian sections, as can be proved comparing their coordinate expressions.
If $`\widehat{\mathrm{\Theta }}`$ is the canonical $`m`$-form in $`J^1E^{}`$, the forms
$$\mathrm{\Theta }_{h_\delta ^{}}:=(\iota _0h_\delta ^{})^{}𝚯=h_\delta ^{}\mathrm{\Theta }\mathrm{\Omega }^m(\mathrm{\Pi }),\mathrm{\Omega }_{h_\delta ^{}}:=\mathrm{d}\mathrm{\Theta }_{h_\delta ^{}}=h_\delta ^{}\mathrm{\Omega }\mathrm{\Omega }^{m+1}(\mathrm{\Pi })$$
(15)
are the Hamilton-Cartan $`m`$ and $`m+1`$ forms of $`\mathrm{\Pi }`$ associated with $``$.
Remark:
* It can be proved that the Hamilton-Cartan $`m`$-form associated with a connection $``$ is the unique form $`\mathrm{\Theta }_{h_\delta ^{}}\mathrm{\Omega }^m(\mathrm{\Pi })`$ such that, if $`\stackrel{~}{y}\mathrm{\Pi }`$ and $`w_1,\mathrm{},w_m\mathrm{T}_{\stackrel{~}{y}}\mathrm{\Pi }`$, then
$`\mathrm{\Theta }_{h_\delta ^{}}(\stackrel{~}{y};w_1,\mathrm{},w_m)`$ $`:=`$ $`(\iota _0h_\delta ^{})(\stackrel{~}{y})(\rho ^1(\stackrel{~}{y});\mathrm{T}_{\stackrel{~}{y}}\rho ^1(w_1),\mathrm{},\mathrm{T}_{\stackrel{~}{y}}\rho ^1(w_m))`$ (16)
$`=`$ $`[\rho ^1(\iota _0h_\delta ^{})](\stackrel{~}{y};w_1,\mathrm{},w_m)`$
that is, $`\mathrm{\Theta }_{h_\delta ^{}}=\rho ^1(\iota _0h_\delta ^{})`$
If $`(x^\nu ,y^A,\mathrm{p}_A^\nu )`$ is a system of natural coordinates in $`\mathrm{\Pi }`$, and $`\stackrel{~}{y}=\mathrm{p}_A^\nu (\stackrel{~}{y}){\displaystyle \frac{}{x^\nu }}\zeta ^A\mathrm{d}^mx\mathrm{\Pi }`$, taking into account the local expression of $`v_{}^{}`$ which, for a connection $`=\mathrm{d}x^\nu \left({\displaystyle \frac{}{x^\nu }}+\mathrm{\Gamma }_\nu ^A{\displaystyle \frac{}{y^A}}\right)`$, is $`v_{}^{}(\zeta ^A)=\mathrm{d}y^A\mathrm{\Gamma }_\nu ^A\mathrm{d}x^\nu `$, we have that
$`h_\delta ^{}(\stackrel{~}{y})`$ $`=`$ $`h_\delta ^{}\left(\mathrm{p}_A^\nu (\stackrel{~}{y}){\displaystyle \frac{}{x^\nu }}\zeta ^A\mathrm{d}^mx|_{\rho ^1(\stackrel{~}{y})}\right)=\mathrm{p}_A^\nu (\stackrel{~}{y}){\displaystyle \frac{}{x^\nu }}(\mathrm{d}y^A\mathrm{\Gamma }_\eta ^A\mathrm{d}x^\eta )\mathrm{d}^mx|_{\rho ^1(\stackrel{~}{y})}`$
$`(\iota _0h_\delta ^{})(\stackrel{~}{y})`$ $`=`$ $`\mathrm{p}_A^\nu (\stackrel{~}{y})(\mathrm{d}y^A\mathrm{d}^{m1}x_\nu \mathrm{\Gamma }_\nu ^A\mathrm{d}^mx)|_{\rho ^1(\stackrel{~}{y})}`$
Observe that $`\iota _0`$ restricted to the image of $`h_\delta ^{}`$ is injective. Therefore
$`\mathrm{\Theta }_{h_\delta ^{}}`$ $`=`$ $`\mathrm{p}_A^\nu (\mathrm{d}y^A\mathrm{\Gamma }_\eta ^A\mathrm{d}x^\eta )\mathrm{d}^{m1}x_\nu =\mathrm{p}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu \mathrm{p}_A^\nu \mathrm{\Gamma }_\nu ^A\mathrm{d}^mx`$
$`\mathrm{\Omega }_{h_\delta ^{}}`$ $`=`$ $`\mathrm{dp}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu +\mathrm{\Gamma }_\nu ^A\mathrm{dp}_A^\nu \mathrm{d}^mx+\mathrm{p}_A^\nu \mathrm{d}\mathrm{\Gamma }_\nu ^A\mathrm{d}^mx`$ (17)
Now, given a connection $``$ and a Hamiltonian section $`h_\delta `$, from Lemma 2 we have that
$$\rho ^1(\iota _0h_\delta ^{}\iota _0h_\delta )=(\iota _0h_\delta ^{})^{}\mathrm{\Theta }(\iota _0h_\delta )^{}\mathrm{\Theta }=\mathrm{\Theta }_{h_\delta ^{}}\mathrm{\Theta }_{h_\delta }=h_\delta ^1\widehat{\mathrm{\Theta }}h_\delta ^2\widehat{\mathrm{\Theta }}:=\text{}_{h_\delta }^{}$$
is a $`\overline{\rho }^1`$-semibasic $`m`$-form in $`\mathrm{\Pi }`$. It is usually written as $`\text{}_{h_\delta }^{}=\mathrm{H}_{h_\delta }^{}(\overline{\rho }^1^{}\omega )`$, where $`\mathrm{H}_{h_\delta }^{}\mathrm{C}^{\mathrm{}}(\mathrm{\Pi })`$ is the global Hamiltonian function associated with $`\text{}_{h_\delta }^{}`$ and $`\omega `$.
Therefore, given a Hamiltonian system $`(\mathrm{\Pi },\mathrm{\Omega }_{h_\delta })`$, taking into account (13), we have that every connection $``$ in $`\pi :EM`$ allows us to split globally the Hamilton-Cartan forms as
$$\mathrm{\Theta }_{h_\delta }=\mathrm{\Theta }_{h_\delta ^{}}\text{}_{h_\delta }^{},\mathrm{\Omega }_{h_\delta }=\mathrm{d}\mathrm{\Theta }_{h_\delta }=\mathrm{\Omega }_{h_\delta ^{}}+\mathrm{d}\text{}_{h_\delta }^{}$$
(18)
In a natural system of coordinates in $`\mathrm{\Pi }`$, such that $`\overline{\rho }^1\omega =\mathrm{d}^mx`$, we write $`\text{}_{h_\delta }^{}=\mathrm{H}_{h_\delta }^{}(x^\nu ,y^A,\mathrm{p}_A^\nu )\mathrm{d}^mx`$, and
$`\mathrm{\Theta }_{h_\delta }`$ $`=`$ $`\mathrm{p}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu (\mathrm{H}_{h_\delta }^{}+\mathrm{p}_A^\nu \mathrm{\Gamma }_\nu ^A)\mathrm{d}^mx`$
$`\mathrm{\Omega }_{h_\delta }`$ $`=`$ $`\mathrm{dp}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu +\mathrm{d}(\mathrm{H}_{h_\delta }^{}+\mathrm{p}_A^\nu \mathrm{\Gamma }_\nu ^A)\mathrm{d}^mx`$ (19)
###### Proposition 7
A couple $`(\{h_\delta \},)`$ in $`\mathrm{\Pi }`$ is equivalent to a couple $`(\text{},)`$ (that is, given a connection $``$, classes of Hamiltonian sections of $`\delta `$ and Hamiltonian densities in $`\mathrm{\Pi }`$ are in one-to-one correspondence).
( Proof ) Given a connection in $`\pi :EM`$, we have just seen that all the Hamiltonian sections belonging to the same equivalence class $`\{h_\delta \}`$ define a unique Hamiltonian density $`\text{}_{h_\delta }^{}`$ and, hence, the same Hamilton-Cartan forms.
Conversely, given a Hamiltonian density and a connection $``$, we can construct an equivalence class of Hamiltonian sections $`\{h_\delta \}`$ (which leads to the same Hamilton-Cartan forms), since, as $`\text{}:\mathrm{\Pi }\pi `$ takes values in $`\overline{\widehat{\rho }}^1\mathrm{\Lambda }^m\mathrm{T}^{}M`$, we have a map $`\iota _0h_\delta ^{}\text{}:\mathrm{\Pi }\pi `$. From the local expression of this map, it is easy to prove that there exists a local section $`h_\delta `$ of $`\delta `$, such that $`\iota _0h_\delta =\iota _0h_\delta ^{}\text{}`$. Then, using a partition of unity we can construct a global section fulfilling this condition, and hence a family of sections $`\{h_\delta \}`$ defined by the relation $`\iota _0h_\delta =\iota _0h_\delta ^{}\text{}`$.
As a direct consequence of this proposition, we have another way of obtaining a Hamiltonian system, which consists in giving a couple $`(\text{},)`$. In fact:
###### Proposition 8
Let $``$ be a connection in $`\pi :EM`$, and a Hamiltonian density. There exist a unique class $`\{h_\delta \}`$ of Hamiltonian sections of $`\delta `$ such that
$$\mathrm{\Theta }_{h_\delta }=\mathrm{\Theta }_{h_\delta ^{}}\text{},\mathrm{\Omega }_{h_\delta }=\mathrm{d}\mathrm{\Theta }_{h_\delta }=\mathrm{\Omega }_{h_\delta ^{}}+\mathrm{d}\text{}$$
Remark:
* If $`\pi :EoM`$ is a trivial bundle, then there is a natural connection (the trivial one). So, in this case, there is a bijective correspondence between Hamiltonian systems and Hamiltonian densities.
This is the situation in classical non-autonomous mechanics , , .
### 3.3 Variational principle and field equations
Now we can establish the field equations for Hamiltonian systems. First we need to introduce the notion of prolongation of diffeomorphisms and vector fields from $`E`$ to $`\mathrm{\Pi }`$.
###### Definition 15
Let $`\mathrm{\Phi }:EE`$ be a diffeomorphism of $`\pi `$-fiber bundles and $`\mathrm{\Phi }_M:MM`$ the induced diffeomorphism in $`M`$. The prolongation of $`\mathrm{\Phi }`$ to $`\mathrm{\Pi }`$ is the diffeomorphism $`j^1\mathrm{\Phi }:\mathrm{\Pi }\mathrm{\Pi }`$ defined by $`j^1\mathrm{\Phi }:=(\mathrm{\Phi }_M)_{}\mathrm{\Phi }^1\stackrel{m}{}\mathrm{\Phi }_M^1`$.
###### Proposition 9
Let $`\mathrm{\Phi }:EE`$ be a diffeomorphism of fiber bundles, $`\mathrm{\Phi }_M:MM`$ its restriction to $`M`$ and $`j^1\mathrm{\Phi }`$ its prolongation to $`\mathrm{\Pi }`$. Then
1. $`\rho ^1j^1\mathrm{\Phi }=\mathrm{\Phi }\rho ^1`$, $`\overline{\rho }^1j^1\mathrm{\Phi }=\mathrm{\Phi }_M\overline{\rho }^1`$.
2. If $`\mathrm{\Psi }:EE`$ is another fiber bundle diffeomorphism, then
$$j^1(\mathrm{\Psi }\mathrm{\Phi })=j^1\mathrm{\Psi }j^1\mathrm{\Phi }$$
3. $`j^1\mathrm{\Phi }`$ is a diffeomorphism of $`\rho ^1`$-bundles and $`\overline{\rho }^1`$-bundles, and $`(j^1\mathrm{\Phi })^1=j^1\mathrm{\Phi }^1`$.
###### Definition 16
Let $`Z\text{X}(E)`$ be a $`\pi `$-projectable vector field. The prolongation of $`Z`$ to $`\mathrm{\Pi }`$ is the vector field $`j^1Z`$ whose local one-parameter group of diffeomorphisms are the extensions $`\{j^1\sigma _t\}`$ of the local one-parameter group of diffeomorphisms $`\{\sigma _t\}`$ of $`Z`$.
Now we can state:
###### Definition 17
Let $`(\mathrm{\Pi },\mathrm{\Omega }_{h_\delta })`$ be a Hamiltonian system. Let $`\mathrm{\Gamma }_c(M,\mathrm{\Pi })`$ be the set of compact-supported sections of $`\overline{\rho }^1`$, and $`\psi \mathrm{\Gamma }_c(M,\mathrm{\Pi })`$. Consider the map
$$\begin{array}{ccccc}𝐇& :& \mathrm{\Gamma }_c(M,\mathrm{\Pi })& & \text{}\\ & & \psi & & _M\psi ^{}\mathrm{\Theta }_{h_\delta }\end{array}$$
The variational problem for this Hamiltonian system is the search of the critical (or stationary) sections of the functional $`𝐇`$, with respect to the variations of $`\psi `$ given by $`\psi _t=j^1\sigma _t\psi `$, where $`\{\sigma _t\}`$ is a local one-parameter group of every $`Z\text{X}^{\mathrm{V}(\pi )}(E)`$ (the module of $`\pi `$-vertical vector fields in $`E`$).
$$\frac{\mathrm{d}}{\mathrm{d}t}|_{t=0}_M\psi _t^{}\mathrm{\Theta }_{h_\delta }=0$$
This is the so-called Hamilton-Jacobi principle of the Hamiltonian formalism.
###### Theorem 4
Let $`(\mathrm{\Pi },\mathrm{\Omega }_{h_\delta })`$ be a Hamiltonian system. The following assertions on a section $`\psi \mathrm{\Gamma }_c(M,\mathrm{\Pi })`$ are equivalent:
1. $`\psi `$ is a critical section for the variational problem posed by $`\mathrm{\Theta }_{h_\delta }`$.
2. $`{\displaystyle _M}\psi ^{}L(j^1Z)\mathrm{\Theta }_{h_\delta }=0`$ , for every $`Z\text{X}^{\mathrm{V}(\pi )}(E)`$.
3. $`\psi ^{}𝑖(j^1Z)\mathrm{\Omega }_{h_\delta }=0`$, for every $`Z\text{X}^{\mathrm{V}(\pi )}(E)`$.
4. $`\psi ^{}𝑖(X)\mathrm{\Omega }_{h_\delta }=0`$, for every $`X\text{X}(\mathrm{\Pi })`$.
5. If $`(U;x^\nu ,y^A,\mathrm{p}_A^\nu )`$ is a natural system of coordinates in $`\mathrm{\Pi }`$, then $`\psi =(x^\nu ,y^A(x^\eta ),\mathrm{p}_A^\nu (x^\eta ))`$ in $`U`$ satisfies the system of equations
$$\frac{y^A}{x^\nu }|_\psi =\frac{H_{h_\delta }}{\mathrm{p}_A^\nu }|_\psi ;\frac{\mathrm{p}_A^\nu }{x^\nu }|_\psi =\frac{H_{h_\delta }}{y^A}|_\psi $$
(20)
which are known as the Hamilton-De Donder-Weyl equations of the Hamiltonian formalism.
( Proof ) ($`12`$) If $`Z\text{X}^{\mathrm{V}(\pi )}(E)`$ and $`\sigma _t`$ is a one-parameter local group of $`Z`$, we have
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}|_{t=0}{\displaystyle _M}(j^1\sigma _t\psi )^{}\mathrm{\Theta }_{h_\delta }`$ $`=`$ $`\underset{t0}{lim}{\displaystyle \frac{1}{t}}({\displaystyle _M}((j^1\sigma _t\psi )^{}\mathrm{\Theta }_{h_\delta }{\displaystyle _M}\psi ^{}\mathrm{\Theta }_{h_\delta })`$
$`=`$ $`\underset{t0}{lim}{\displaystyle \frac{1}{t}}\left({\displaystyle _M}\psi ^{}(j^1\sigma _t)^{}\mathrm{\Theta }_{h_\delta }{\displaystyle _M}\psi ^{}\mathrm{\Theta }_{h_\delta }\right)`$
$`=`$ $`\underset{t0}{lim}{\displaystyle \frac{1}{t}}\left({\displaystyle _M}\psi ^{}[(j^1\sigma _t)^{}\mathrm{\Theta }_{h_\delta }\mathrm{\Theta }_{h_\delta }]\right)={\displaystyle _M}\psi ^{}L(j^1Z)\mathrm{\Theta }_{h_\delta }`$
and the results follows immediately.
($`23`$) Taking into account that
$$L(j^1Z)\mathrm{\Theta }_{h_\delta }=\mathrm{d}𝑖(j^1Z)\mathrm{\Theta }_{h_\delta }+𝑖(j^1Z)\mathrm{d}\mathrm{\Theta }_{h_\delta }=\mathrm{d}𝑖(j^1Z)\mathrm{\Theta }_{h_\delta }𝑖(j^1Z)\mathrm{\Omega }_{h_\delta }$$
we obtain that
$$_M\psi ^{}L(j^1Z)\mathrm{\Theta }_{h_\delta }=_M\psi ^{}\mathrm{d}𝑖(j^1Z)\mathrm{\Theta }_{h_\delta }_M\psi ^{}𝑖(j^1Z)\mathrm{\Omega }_{h_\delta }$$
and, as $`\psi `$ has compact support, using Stoke’s theorem we have
$$_M\psi ^{}\mathrm{d}𝑖(j^1Z)\mathrm{\Theta }_{h_\delta }=_Md\psi ^{}𝑖(j^1Z)\mathrm{\Theta }_{h_\delta }=0$$
hence $`{\displaystyle _M}\psi ^{}L(j^1Z)\mathrm{\Theta }_{h_\delta }=0`$ (for every $`Z\text{X}^{\mathrm{V}(\pi )}(E)`$) if, and only if, $`{\displaystyle _M}\psi ^{}𝑖(j^1Z)\mathrm{\Omega }_{h_\delta }=0`$ , and, according to the fundamental theorem of variational calculus, this is equivalent to
$$\psi ^{}𝑖(j^1Z)\mathrm{\Omega }_{h_\delta }=0$$
($`35`$) Suppose that $`\psi `$ is a section verifying that $`\psi ^{}𝑖(j^1Z)\mathrm{\Omega }_{h_\delta }=0`$, for every $`Z\text{X}^{\mathrm{V}(\pi )}(E)`$. In a natural chart in $`\mathrm{\Pi }`$, if $`Z=\beta ^A{\displaystyle \frac{}{y^A}}`$ , then $`j^1Z=\beta ^A{\displaystyle \frac{}{y^A}}\mathrm{p}_B^\nu {\displaystyle \frac{\beta ^B}{y^A}}{\displaystyle \frac{}{\mathrm{p}_A^\nu }}`$ . Taking into account the local expression of $`\mathrm{\Omega }_{h_\delta }`$ given in (11), we have
$$𝑖(j^1Z)\mathrm{\Omega }_{h_\delta }=\beta ^A\left(\mathrm{dp}_A^\nu \mathrm{d}^{m1}x_\nu +\frac{H_{h_\delta }}{y^A}\mathrm{d}^mx\right)+\mathrm{p}_B^\nu \frac{\beta ^B}{y^A}\left(\mathrm{d}y^A\mathrm{d}^{m1}x_\nu \frac{H_{h_\delta }}{\mathrm{p}_A^\nu }\mathrm{dp}_A^\nu \mathrm{d}^mx\right)$$
As $`\psi =(x^\nu ,f^A(x^\eta ),g_A^\nu (x^\eta ))`$ is a section of $`\overline{\rho }^1`$, then on the points of the image of $`\psi `$ we have $`y^A=f^A(x^\eta )`$, $`\mathrm{p}_A^\nu =g_A^\nu (x^\eta )`$, and we obtain
$`0=\psi ^{}𝑖(j^1Z)\mathrm{\Omega }_{h_\delta }`$ $`=`$ $`\beta ^A\left({\displaystyle \frac{g_A^\nu }{x_\nu }}+{\displaystyle \frac{H_{h_\delta }}{y^A}}\right)\mathrm{d}^mx+g_C^\nu \left({\displaystyle \frac{\beta ^C}{y^A}}{\displaystyle \frac{f^A}{x_\nu }}{\displaystyle \frac{\beta ^C}{y^A}}{\displaystyle \frac{H_{h_\delta }}{\mathrm{p}_A^\nu }}\right)\mathrm{d}^mx`$
$`=`$ $`\left[\beta ^A\left({\displaystyle \frac{g_A^\nu }{x^\nu }}+{\displaystyle \frac{H_{h_\delta }}{y^A}}\right)+g_C^\nu {\displaystyle \frac{\beta ^C}{y^A}}\left({\displaystyle \frac{f^A}{x^\nu }}{\displaystyle \frac{H_{h_\delta }}{\mathrm{p}_A^\nu }}\right)\right]\mathrm{d}^mx`$
and, as this holds for every $`Z\text{X}^{\mathrm{V}(\pi )}(E)`$, this is equivalent to demanding that
$$\beta ^A\left(\frac{g_A^\nu }{x^\nu }+\frac{H_{h_\delta }}{y^A}\right)+g_C^\nu \frac{\beta ^C}{y^A}\left(\frac{f^A}{x^\nu }\frac{H_{h_\delta }}{\mathrm{p}_A^\nu }\right)=0$$
for every $`\beta ^A(x^\nu ,y^A)`$. Therefore
$$\beta ^A(\frac{g_A^\nu }{x^\nu }+\frac{H_{h_\delta }}{y^A})=0;g_C^\nu \frac{\beta ^C}{y^A}(\frac{f^A}{x^\nu }\frac{H_{h_\delta }}{\mathrm{p}_A^\nu })=0$$
From the first equalities we obtain the first group of the Hamiltonian equations. For the second ones, let $`(W;x^\nu ,y^A,\mathrm{p}_A^\nu )`$ a natural chart, $`U=\overline{\rho }^1(W)`$, and $`\psi `$ a critical section. Then, for every $`xU`$ we have
$$g_C^\nu (x)\frac{\beta ^C}{y^A}|_{(x,f^A(x))}\left(\frac{f^A}{x^\nu }\frac{H_{h_\delta }}{\mathrm{p}_A^\nu }\right)|_{\psi (x)}=0$$
but, as there are critical sections passing through every point in $`W`$, we obtain that
$$\frac{\beta ^C}{y^A}|_{(x,f^A(x))}\left(\frac{f^A}{x^\nu }\frac{H_{h_\delta }}{\mathrm{p}_A^\nu }\right)|_{\psi (x)}=0\text{(for every }B,\nu \text{)}$$
Now we can choose $`\beta ^B`$ such that $`{\displaystyle \frac{\beta ^B}{y^A}}`$ take arbitrary values, and then
$$\frac{f^A}{x^\nu }\frac{H_{h_\delta }}{\mathrm{p}_A^\nu }=0$$
which is the second group of the Hamiltonian equations. The converse is trivial.
($`45`$) Suppose that $`\psi `$ is a section verifying that $`\psi ^{}𝑖(X)\mathrm{\Omega }_{h_\delta }=0`$, for every $`X\text{X}(\mathrm{\Pi })`$. If $`X=\alpha ^\nu {\displaystyle \frac{}{x^\nu }}+\beta ^A{\displaystyle \frac{}{y^A}}+\gamma _A^\nu {\displaystyle \frac{}{\mathrm{p}_A^\nu }}`$ , taking into account (11), we have
$`𝑖(X)\mathrm{\Omega }_{h_\delta }`$ $`=`$ $`(1)^\eta \alpha ^\eta \left(\mathrm{dp}_A^\nu \mathrm{d}y^A\mathrm{d}^{m2}x_{\eta \nu }{\displaystyle \frac{H_{h_\delta }}{\mathrm{p}_A^\nu }}\mathrm{dp}_A^\nu \mathrm{d}^{m1}x_\eta \right)`$
$`+`$ $`\beta ^A\left(\mathrm{dp}_A^\nu \mathrm{d}^{m1}x_\nu +{\displaystyle \frac{H_{h_\delta }}{y^A}}\mathrm{d}^mx\right)+\gamma _A^\nu \left(\mathrm{d}y^A\mathrm{d}^{m1}x_\nu +{\displaystyle \frac{H_{h_\delta }}{\mathrm{p}_A^\nu }}\mathrm{dp}_A^\nu \mathrm{d}^mx\right)`$
but as $`\psi =(x^\nu ,f^A(x^\eta ),g_A^\nu (x^\eta ))`$ is a section of $`\overline{\rho }^1`$, then on the points of the image of $`\psi `$ we have $`y^A=f^A(x^\eta )`$, $`\mathrm{p}_A^\nu =g_A^\nu (x^\eta )`$, and
$`0=\psi ^{}𝑖(X)\mathrm{\Omega }_{h_\delta }`$ $`=`$ $`(1)^{\eta +\nu }\alpha ^\eta \left({\displaystyle \frac{f^A}{x^\nu }}{\displaystyle \frac{H_{h_\delta }}{\mathrm{p}_A^\nu }}\right){\displaystyle \frac{g_A^\nu }{x^\eta }}\mathrm{d}^mx+`$
$`\beta ^A\left({\displaystyle \frac{g_A^\nu }{x_\nu }}+{\displaystyle \frac{H_{h_\delta }}{y^A}}\right)\mathrm{d}^mx+\gamma _A^\nu \left({\displaystyle \frac{f^A}{x_\nu }}+{\displaystyle \frac{H_{h_\delta }}{\mathrm{p}_A^\nu }}\right)\mathrm{d}^mx`$
and, as this holds for every $`X\text{X}(\mathrm{\Pi })`$, we obtain the Hamilton-De Donder-Weyl equations. The converse is trivial.
Remark:
* In relation to the equations (20), it is important to point out that they are not covariant, since the Hamiltonian function $`H_{h_\delta }`$ is defined only locally, and hence it is not intrinsically defined.
In order to write a set of covariant Hamiltonian equations we must use a global Hamiltonian function, which can be obtained by introducing another Hamiltonian section $`h_\delta ^{}`$, with $`h_\delta ^{}\{h_\delta \}`$ (as we have seen in section 3.1). It is usual to take the section induced by a connection $``$ in $`\pi :EM`$, and hence we have the splitting given in (18) for the form $`\mathrm{\Omega }_{h_\delta }`$. Then, if $`\mathrm{\Gamma }_\eta ^B`$ are the local component functions of $``$ in $`U\mathrm{\Pi }`$, starting from the local expression (19), and following the same pattern as in the proof of the last item, we obtain for a critical section $`\psi =(x^\nu ,y^A(x^\eta ),\mathrm{p}_A^\nu (x^\eta ))`$ in $`U`$ the covariant Hamiltonian equations:
$$\frac{y^A}{x^\nu }|_\psi =(\frac{\mathrm{H}_{h_\delta }^{}}{\mathrm{p}_A^\nu }+\mathrm{\Gamma }_\nu ^A)|_\psi ;\frac{\mathrm{p}_A^\nu }{x^\nu }|_\psi =(\frac{\mathrm{H}_{h_\delta }^{}}{y^A}+\mathrm{p}_\eta ^B\frac{\mathrm{\Gamma }_\eta ^B}{y^A})|_\psi $$
Observe that, as $`\mathrm{H}_{h_\delta }^{}=H_{h_\delta }\mathrm{p}_A^\nu \mathrm{\Gamma }_\nu ^A`$ (on each open set $`U\mathrm{\Pi }`$, where $`H_{h_\delta }`$ is the corresponding local Hamiltonian function), then from these last equations we recover the Hamilton-De Donder-Weyl equations. (See for comments on this subject).
### 3.4 Hamiltonian system associated with a hyper-regular Lagrangian system
It is evident that different choices of equivalence classes of Hamiltonian sections of $`\delta `$ lead to different Hamiltonian systems in $`\mathrm{\Pi }`$. The question now is how to associate (if possible) a Hamiltonian system with a Lagrangian system. The answer to this question is closely related to the regularity of the Lagrangian system.
First, let $`(J^1E,\mathrm{\Omega }_{})`$ be a hyper-regular Lagrangian system. Then:
###### Lemma 3
For every section $`\mathrm{h}_\delta :\mathrm{\Pi }J^1E^{}`$ of $`\delta `$, the relation
$$\mathrm{h}_\mu :=\mu ^{}\iota _0\mathrm{h}_\delta \mathrm{\Psi }$$
defines a unique section of $`\mu `$, which is just $`\mathrm{h}_\mu =\stackrel{~}{}^1`$.
( Proof ) We have the diagram
$$\begin{array}{ccccc}\text{}& \text{}& \text{}& \text{}& \text{}\end{array}$$
(21)
Then, taking into account the commutativity of this diagram, we have
$`\mathrm{h}_\mu `$ $`=`$ $`\mu ^{}\iota _0\mathrm{h}_\delta \mathrm{\Psi }=\stackrel{~}{}^1\mu \iota _0\mathrm{h}_\delta \mathrm{\Psi }`$
$`=`$ $`\stackrel{~}{}^1\mathrm{\Psi }^1\mathrm{\Psi }=\stackrel{~}{}^1`$
So $`\mathrm{h}_\mu `$ is independent of $`\mathrm{h}_\delta `$.
Remarks:
* This result is to be expected, since $`\stackrel{~}{}(J^1E)`$ is a 1-codimensional submanifold of $`\pi `$, transverse to the projection $`\mu `$, and hence it defines a section $`\mathrm{h}_\mu `$ of $`\mu `$. This is just the natural section used in and for associating a Hamiltonian system to a hyper-regular Lagrangian one (see section 4.2).
* Observe that a natural section $`\mathrm{h}_\delta `$ of $`\delta `$ can be selected by making
$$\mathrm{h}_\delta :=\widehat{\mathrm{F}}\mathrm{F}^1$$
or, what is equivalent, its associated class can be defined by
$$\iota _0h_\delta :=\widehat{}\mathrm{F}^1,\text{for every }h_\delta \{\mathrm{h}_\delta \}$$
Observe that this section $`\mathrm{h}_\delta `$ is just the inverse of $`\delta `$ restricted to $`\widehat{\mathrm{F}}(J^1E)`$, and that $`\iota _0h_\delta `$ is a diffeomorphism.
###### Definition 18
Given a section $`\mathrm{h}_\delta :\mathrm{\Pi }J^1E^{}`$ of $`\delta `$, we define the Hamilton-Cartan forms
$$\mathrm{\Theta }_{\mathrm{h}_\delta }:=(\mu ^{}\iota _0\mathrm{h}_\delta )^{}\mathrm{\Theta };\mathrm{\Omega }_{\mathrm{h}_\delta }:=(\mu ^{}\iota _0\mathrm{h}_\delta )^{}\mathrm{\Omega }$$
###### Proposition 10
The Hamilton-Cartan forms are independent of the section $`\mathrm{h}_\delta `$, and
$$\mathrm{F}^{}\mathrm{\Theta }_{\mathrm{h}_\delta }=\mathrm{\Theta }_{},\mathrm{F}^{}\mathrm{\Omega }_{\mathrm{h}_\delta }=\mathrm{\Omega }_{}$$
(22)
Then, $`(\mathrm{\Pi },\mathrm{\Omega }_{\mathrm{h}_\delta })`$ is the (unique) Hamiltonian system which is associated with the hyper-regular Lagrangian system $`(J^1E,\mathrm{\Omega }_{})`$.
( Proof ) The independence of the section $`\mathrm{h}_\delta `$ is a consequence of lemma 3. Then, taking into account the commutativity of diagram (21), and proposition 2, for every section $`\mathrm{h}_\delta `$, we have
$$\mathrm{F}^{}\mathrm{\Theta }_{\mathrm{h}_\delta }=\mathrm{F}^{}(\mu ^{}\iota _0\mathrm{h}_\delta )^{}\mathrm{\Theta }=(\mu ^{}\iota _0\mathrm{h}_\delta \mathrm{F})^{}\mathrm{\Theta }=\stackrel{~}{}^{}\mathrm{\Theta }=\mathrm{\Theta }_{}$$
and the same result follows for $`\mathrm{\Omega }_{\mathrm{h}_\delta }`$.
Using charts of natural coordinates in $`\mathrm{\Pi }`$ and $`J^1E^{}`$, and the expressions (1) and (2) of the Legendre maps, we have that the natural Hamiltonian section $`\mathrm{h}_\delta =\widehat{\mathrm{F}}\mathrm{F}^1`$ has associated the local Hamiltonian function
$$H_{\mathrm{h}_\delta }(x^\nu ,y^A,\mathrm{p}_A^\nu )=\mathrm{F}^1\left(v_\nu ^A\frac{\mathrm{\pounds }}{v_\nu ^A}\mathrm{\pounds }\right)=\mathrm{p}_A^\nu \mathrm{F}^1^{}v_\nu ^A\mathrm{F}^1^{}\mathrm{\pounds }$$
(23)
and for the Hamilton-Cartan forms:
$`\mathrm{\Theta }_{\mathrm{h}_\delta }`$ $`=`$ $`\mathrm{p}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu (\mathrm{p}_A^\nu \mathrm{F}^1^{}v_\nu ^A\mathrm{F}^1^{}\mathrm{\pounds })\mathrm{d}^mx`$
$`\mathrm{\Omega }_{\mathrm{h}_\delta }`$ $`=`$ $`\mathrm{dp}_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu +\mathrm{d}(\mathrm{p}_A^\nu \mathrm{F}^1^{}v_\nu ^A\mathrm{F}^1^{}\mathrm{\pounds })\mathrm{d}^mx`$
There is another way of obtaining this Hamiltonian system. In fact, suppose that a connection $``$ is given in $`\pi :EM`$, and let $`h_\delta ^{}:\mathrm{\Pi }J^1E^{}`$ be the induced linear section of $`\delta `$. If we have used $``$ for constructing the associated density of Lagrangian energy $`_{}^{}\mathrm{\Omega }^m(J^1E)`$ (see definition 1), the key is to define a Hamiltonian density $`\text{}^{}\mathrm{\Omega }^m(\mathrm{\Pi })`$ which is $`\mathrm{F}`$-related with $`_{}^{}`$. We can make this construction in two ways:
###### Proposition 11
1. The $`m`$-form $`\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}}\mathrm{\Theta }_{}`$ is $`\overline{\pi }^1`$-semibasic and
$$\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}}\mathrm{\Theta }_{}=_{}^{}$$
(24)
2. There exists a unique Hamiltonian density $`\text{}^{}\mathrm{\Omega }^m(\mathrm{\Pi })`$ such that
$$\mathrm{F}^{}\text{}^{}=\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}}\mathrm{\Theta }_{}=_{}^{}$$
(25)
Let $`\text{}^{}=\mathrm{H}^{}(\overline{\rho }^1\omega )`$, with $`\mathrm{H}^{}\mathrm{C}^{\mathrm{}}(\mathrm{\Pi })`$. Then, $`\text{}^{}`$ and $`\mathrm{H}^{}`$ are called the Hamiltonian density and the Hamiltonian function associated with the Lagrangian system, the connection $``$ and $`\omega `$.
3. The Hamilton-Cartan forms of definition 18 split as
$$\mathrm{\Theta }_{\mathrm{h}_\delta }=\mathrm{\Theta }_{h_\delta ^{}}\text{}^{},\mathrm{\Omega }_{\mathrm{h}_\delta }=\mathrm{d}\mathrm{\Theta }_{\mathrm{h}_\delta }=\mathrm{\Omega }_{h_\delta ^{}}+\mathrm{d}\text{}^{}$$
(26)
( Proof )
1. Once again, it suffices to see it in a natural local system $`(x^\nu ,y^A,v_\nu ^A)`$. Then, if $`=\mathrm{\pounds }\mathrm{d}^mx`$, taking into account the corresponding local expressions we have that
$$\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}}\mathrm{\Theta }_{}=\left(\frac{\mathrm{\pounds }}{v_\nu ^A}(v_\nu ^A\mathrm{\Gamma }_\nu ^A)\mathrm{\pounds }\right)\mathrm{d}^mx$$
and the result holds. The last part follows, recalling the local expression of the density of Lagrangian energy. Thus this form is $`\overline{\pi }^1`$-semibasic.
2. It is immediate, as $`\mathrm{F}`$ is a diffeomorphism.
3. From (24) and (25) we obtain that
$$\mathrm{F}^{}(\mathrm{\Theta }_{h_\delta ^{}}\text{}^{})=\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}}\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}}+\mathrm{\Theta }_{}=\mathrm{\Theta }_{}=\mathrm{F}^{}\mathrm{\Theta }_{\mathrm{h}_\delta }$$
and therefore $`\mathrm{F}^{}\mathrm{\Omega }_{\mathrm{h}_\delta }=\mathrm{\Omega }_{}`$ too. Then, the result follows because $`\mathrm{F}`$ is a diffeomorphism.
Remark:
* Notice that the item 1 holds even if $`\mathrm{F}`$ is not a diffeomorphism.
In a system of natural coordinates we have
$$\mathrm{H}^{}(x^\nu ,y^A,\mathrm{p}_A^\nu )=\mathrm{p}_A^\nu (\mathrm{F}^1^{}v_\nu ^A\mathrm{\Gamma }_\nu ^A)\mathrm{F}^1^{}\mathrm{\pounds }$$
and thus $`\mathrm{F}^{}\mathrm{H}^{}=\mathrm{E}_{}^{}`$.
An alternative way is to obtain this Hamiltonian density using only Hamiltonian sections.
###### Proposition 12
Consider the Hamiltonian section $`\mathrm{h}_\delta =\widehat{\mathrm{F}}\mathrm{F}^1`$, and a connection $``$. Then we have that
$$(\iota _0h_\delta ^{})^{}\mathrm{\Theta }(\mu ^{}\iota _0\mathrm{h}_\delta )^{}\mathrm{\Theta }=\text{}^{}$$
and hence the splitting (26) holds
( Proof ) We have the following diagram
$$\begin{array}{ccccccc}\text{}& \text{}& \text{}& \text{}& \text{}& \text{}& \text{}\end{array}$$
Observe that $`\mu ^{}\iota _0\mathrm{h}_\delta =\stackrel{~}{}\mathrm{F}^1`$. Therefore, taking into account definition 7, Proposition 2, (15) and (25), we have
$$(\iota _0h_\delta ^{})^{}\mathrm{\Theta }(\mu ^{}\iota _0\mathrm{h}_\delta )^{}\mathrm{\Theta }=\mathrm{\Theta }_{h_\delta ^{}}(\mathrm{F}^1)^{}\stackrel{~}{}^{}\mathrm{\Theta }=\mathrm{\Theta }_{h_\delta ^{}}(\mathrm{F}^1)^{}\mathrm{\Theta }_{}=\text{}^{}$$
Then the result for the splittings of the Hamilton-Cartan forms follows straightforwardly.
Remark:
* Note that the use of both extended Legendre maps is necessary for obtaining the Hamiltonian density in this way.
As a final remark, all the results stated in section 3.3 in relation to the variational principle and the characterization of critical sections are true. In particular, field equations are the Hamilton-De Donder-Weyl equations (20), where the local Hamiltonian function $`H_{\mathrm{h}_\delta }`$ is given by (23).
### 3.5 Hamiltonian system associated with an almost-regular Lagrangian system
Now, let $`(J^1E,\mathrm{\Omega }_{})`$ be an almost-regular Lagrangian system. Bearing in mind diagram (9), first observe that the submanifold $`ȷ_0:P\mathrm{\Pi }`$, is a fiber bundle over $`E`$ (and $`M`$), and the corresponding projections will be denoted $`\kappa _0^1:PE`$ and $`\overline{\kappa }_0^1:PM`$, satisfying that $`\kappa ^1ȷ_0=\kappa _0^1`$ and $`\overline{\kappa }^1ȷ_0=\overline{\kappa }_0^1`$.
###### Proposition 13
The Lagrangian forms $`\mathrm{\Theta }_{}`$ and $`\mathrm{\Omega }_{}`$, are $`\mathrm{F}`$-projectable.
( Proof ) By Proposition 5, we have that $`\mathrm{ker}\mathrm{F}_{}=\mathrm{ker}\mathrm{\Omega }_{}\text{X}^{\mathrm{V}(\pi ^1)}(J^1E)`$. Then, for every $`X\mathrm{ker}\mathrm{F}_{}`$ we have that $`𝑖(X)\mathrm{\Theta }_{}=0`$, since $`\mathrm{\Theta }_{}`$ is a $`\pi ^1`$-semibasic $`m`$-form, and in the same way $`L(X)\mathrm{\Theta }_{}=0`$. Therefore $`\mathrm{\Theta }_{}`$ is $`\mathrm{F}`$-projectable.
As a trivial consequence of this fact, $`𝑖(X)\mathrm{\Omega }_{}=0`$, and $`L(X)\mathrm{\Omega }_{}=0`$, and therefore $`\mathrm{\Omega }_{}`$ is also $`\mathrm{F}`$-projectable.
###### Definition 19
Given a section $`\widehat{\mathrm{h}}_\delta :P\widehat{P}`$ of $`\widehat{\delta }`$, we define the Hamilton-Cartan forms
$$\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0:=(\stackrel{~}{ȷ}_0\widehat{\mu }^{}\widehat{\iota }_0\widehat{\mathrm{h}}_\delta )^{}\mathrm{\Theta };\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0:=(\stackrel{~}{ȷ}_0\widehat{\mu }^{}\widehat{\iota }_0\widehat{\mathrm{h}}_\delta )^{}\mathrm{\Omega }$$
###### Proposition 14
The Hamilton-Cartan forms $`\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0`$ and $`\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0`$ are independent of the section $`\widehat{\mathrm{h}}_\delta `$ of $`\widehat{\delta }`$, and
$$\mathrm{F}_0^{}\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0=\mathrm{\Theta }_{},\mathrm{F}_0^{}\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0=\mathrm{\Omega }_{}$$
(27)
Then $`(\mathrm{\Pi },P,\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0)`$ is the unique Hamiltonian system associated with the almost-regular Lagrangian system $`(J^1E,\mathrm{\Omega }_{})`$.
( Proof ) We have the following diagram
$$\begin{array}{ccccc}\text{}& \text{}& \text{}& \text{}& \text{}\end{array}$$
Then, taking into account the commutativity of this diagram, and proposition 2, for every section $`\widehat{\mathrm{h}}_\delta `$ of $`\widehat{\delta }`$ we have that
$`\mathrm{F}_0^{}\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0`$ $`=`$ $`\mathrm{F}_0^{}(\stackrel{~}{ȷ}_0\widehat{\mu }^{}\widehat{\iota }_0\widehat{\mathrm{h}}_\delta )^{}\mathrm{\Theta }=(\stackrel{~}{ȷ}_0\widehat{\mu }^{}\widehat{\iota }_0\widehat{\mathrm{h}}_\delta \mathrm{F}_0)^{}\mathrm{\Theta }`$
$`=`$ $`(\stackrel{~}{ȷ}_0\stackrel{~}{}_0)^{}\mathrm{\Theta }=\stackrel{~}{}^{}\mathrm{\Theta }=\mathrm{\Theta }_{}`$
and the same result follows for $`\mathrm{\Omega }_{\mathrm{h}_\delta }^0`$.
Remarks:
* Following the terminology of sections above, we have that all the sections $`\widehat{\mathrm{h}}_\delta `$ belong to the same equivalence class.
* In the particular situation that $`\mathrm{rank}\widehat{\mathrm{F}}_0=\mathrm{rank}\widehat{}_0`$, we have that $`\widehat{\iota }_0`$ is a diffeomorphism and, as the fibers of $`\widehat{\delta }`$ are also the fibers of $`\widehat{\iota }_0`$, then so is $`\widehat{\delta }`$. In this case there is only one map $`\widehat{\mathrm{h}}_\delta `$, which is just $`\widehat{\delta }^1`$.
As in the hyper-regular case, we can construct this Hamiltonian system using a connection. Thus, let $``$ be connection in $`\pi :EM`$, and $`h_\delta ^{}:\mathrm{\Pi }J^1E^{}`$ the induced linear section of $`\delta `$. Let $`_{}^{}\mathrm{\Omega }^m(J^1E)`$ be the density of Lagrangian energy associated with $``$ (see definition 1). Then:
###### Proposition 15
1. The density of Lagrangian energy $`_{}^{}`$ is $`\mathrm{F}`$-projectable.
2. The $`\overline{\pi }^1`$-semibasic $`m`$-form $`\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}}\mathrm{\Theta }_{}`$ is $`\mathrm{F}`$-projectable and denoting $`\mathrm{\Theta }_{h_\delta ^{}}^0=ȷ_0^{}\mathrm{\Theta }_{h_\delta ^{}}`$, we have
$$_{}^{}=\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}}\mathrm{\Theta }_{}=\mathrm{F}_0^{}\mathrm{\Theta }_{h_\delta ^{}}^0\mathrm{\Theta }_{}$$
3. There exists a unique $`\overline{\rho }_0^1`$-semibasic form $`\text{}_0^{}\mathrm{\Omega }^m(P)`$, such that
$$\mathrm{F}_0^{}\text{}_0^{}=_{}^{}$$
Let $`\text{}_0^{}=\mathrm{H}_0^{}(\overline{\rho }_0^1\omega )`$, with $`\mathrm{H}_0^{}\mathrm{C}^{\mathrm{}}(P)`$. Then $`\text{}_0^{}`$ and $`\mathrm{H}_0^{}`$ are called the Hamiltonian density and the Hamiltonian function associated with the Lagrangian system, the connection $``$ and $`\omega `$. Obviously we have that $`\mathrm{F}_0^{}\mathrm{H}_0^{}=\mathrm{E}_{}^{}`$.
4. The Hamilton-Cartan forms of definition 19 split as
$`\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0`$ $`=`$ $`ȷ_0^{}\mathrm{\Theta }_{h_\delta ^{}}\text{}_0^{}=\mathrm{\Theta }_{h_\delta ^{}}^0\text{}_0^{}`$
$`\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0`$ $`=`$ $`\mathrm{d}\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0=ȷ_0^{}\mathrm{\Omega }_{h_\delta ^{}}+\mathrm{d}\text{}_0^{}=\mathrm{\Omega }_{h_\delta ^{}}^0+\mathrm{d}\text{}_0^{}`$ (28)
( Proof )
1. As $`_{}^{}=\mathrm{E}_{}^{}(\overline{\pi }^1\omega )`$, it suffices to prove that the Lagrangian energy $`\mathrm{E}_{}^{}`$ is $`\mathrm{F}`$-projectable. Then, for every $`X\mathrm{ker}\mathrm{F}_{}`$, using natural coordinates we have
$`L(X)\mathrm{E}_{}^{}`$ $`=`$ $`L\left(f_\eta ^B{\displaystyle \frac{}{v_\eta ^B}}\right)\left({\displaystyle \frac{\mathrm{\pounds }}{v_\nu ^A}}(v_\nu ^A\mathrm{\Gamma }_\nu ^A)\mathrm{\pounds }\right)`$
$`=`$ $`f_\eta ^B{\displaystyle \frac{^2\mathrm{\pounds }}{v_\eta ^Bv_\nu ^A}}(v_\nu ^A\mathrm{\Gamma }_\nu ^A)+f_\eta ^B{\displaystyle \frac{\mathrm{\pounds }}{v_\nu ^A}}\delta _B^A\delta _\nu ^\eta f_\eta ^B{\displaystyle \frac{\mathrm{\pounds }}{v_\eta ^B}}=0`$
therefore $`\mathrm{E}_{}^{}`$ is $`\mathrm{F}`$-projectable, and so is $`_{}^{}`$.
2. It is immediate, taking into account that $`\mathrm{\Theta }_{}`$ is $`\mathrm{F}`$-projectable, and the first item of Proposition 11.
3. The existence is assured, since $`_{}^{}`$ is $`\mathrm{F}`$-projectable and the uniqueness because $`\mathrm{F}_0`$ is a submersion.
Next we prove that $`\text{}_0^{}`$ is $`\overline{\rho }_0^1`$-semibasic. As $`\mathrm{F}_0`$ is a submersion, for every $`yJ^1E`$ and $`\stackrel{~}{u}\mathrm{V}_{\mathrm{F}_0(y)}(\overline{\rho }_0^1)`$, there exist $`u\mathrm{T}_y(J^1E)`$ such that $`\stackrel{~}{u}=\mathrm{T}_y\mathrm{F}_0(u)`$ and, in addition, $`u\mathrm{V}_y(\overline{\pi }^1)`$ because
$$\mathrm{T}_y\overline{\pi }^1(u)=(\mathrm{T}_{\mathrm{F}_0(\overline{y})}\overline{\rho }_0^1\mathrm{T}_y\mathrm{F}_0)(u)=\mathrm{T}_{\mathrm{F}_0(\overline{y})}\overline{\rho }_0^1(\stackrel{~}{u})=0$$
Furthermore, $`_{}^{}`$ is $`\overline{\pi }^1`$-semibasic, and hence
$$0=𝑖(u)[_{}^{}(\overline{y})]=𝑖(u)[(\mathrm{F}_0^{}\text{}_0^{})(y)]=(\mathrm{F}_0)_{\mathrm{F}_0(\overline{y})}^{}[𝑖(\stackrel{~}{u})(\text{}_0^{}(\mathrm{F}_0(\overline{y}))]$$
then, for every $`yJ^1E`$ and $`\stackrel{~}{u}\mathrm{V}_{\mathrm{F}_0(\overline{y})}(\overline{\rho }_0^1)`$ we have $`𝑖(\stackrel{~}{u})(\text{}_0^{}(\mathrm{F}_0(\overline{y}))\mathrm{ker}(\mathrm{F}_0)_{\mathrm{F}_0(\overline{y})}^{}=\{0\}`$, since $`\mathrm{F}_0`$ is a submersion. So $`\text{}_0^{}`$ is $`\overline{\rho }_0^1`$-semibasic.
4. Taking into account items 3 and 2, we obtain
$$\mathrm{F}_0^{}(\mathrm{\Theta }_{h_\delta ^{}}^0\text{}_0^{})=\mathrm{F}_0^{}\mathrm{\Theta }_{h_\delta ^{}}^0\mathrm{F}_0^{}\text{}_0^{}=\mathrm{F}_0^{}\mathrm{\Theta }_{h_\delta ^{}}^0_{}^{}=\mathrm{\Theta }_{}=\mathrm{F}_0^{}\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0$$
and therefore $`\mathrm{\Omega }_{}=\mathrm{F}_0^{}\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0`$ too. Then the result follows because $`\mathrm{F}`$ is a submersion.
We can construct the above Hamiltonian density in an alternative way, as follows:
###### Proposition 16
Let $`\widehat{\mathrm{h}}_\delta :P\widehat{P}`$ be a section of $`\widehat{\delta }`$, and $``$ a connection. Then, $`h_\delta ^{}`$ induces a map $`\widehat{h}_\delta ^{}:P\widehat{P}`$ defined by the relation $`\widehat{ȷ}_0\widehat{h}_\delta ^{}=h_\delta ^{}ȷ_0`$. Therefore
$$(\widehat{ȷ}_0\widehat{\iota }_0\widehat{h}_\delta ^{})^{}\mathrm{\Theta }(\stackrel{~}{ȷ}_0\widehat{\mu }^{}\widehat{\iota }_0\widehat{\mathrm{h}}_\delta )^{}\mathrm{\Theta }=\text{}_0^{}$$
and hence the splitting (28) holds.
( Proof ) We have the following diagram
$$\begin{array}{cccccc}\text{}& \text{}& \text{}& \text{}& \text{}& \text{}\end{array}$$
Taking into account the commutativity of this diagram, we have that every section $`\widehat{\mathrm{h}}_\delta `$ of $`\widehat{\delta }`$ satisfies that
$$\widehat{\mu }^{}\widehat{\iota }_0\widehat{\mathrm{h}}_\delta \mathrm{F}_0=\stackrel{~}{}_0$$
Then, bearing in mind the second item of Proposition 15 and (27), we have that
$`\mathrm{F}_0^{}[(\widehat{ȷ}_0\widehat{\iota }_0\widehat{h}_\delta ^{})^{}\mathrm{\Theta }(\stackrel{~}{ȷ}_0\widehat{\mu }^{}\widehat{\iota }_0\widehat{\mathrm{h}}_\delta )^{}\mathrm{\Theta }]`$ $`=`$ $`(\widehat{ȷ}_0\widehat{\iota }_0\widehat{h}_\delta ^{}\mathrm{F}_0)^{}\mathrm{\Theta }(\stackrel{~}{ȷ}_0\widehat{\mu }^{}\widehat{\iota }_0\widehat{\mathrm{h}}_\delta \mathrm{F}_0)^{}\mathrm{\Theta }`$
$`=`$ $`(\iota _0\widehat{ȷ}_0\widehat{h}_\delta ^{}\mathrm{F}_0)^{}\mathrm{\Theta }(\stackrel{~}{ȷ}_0\widehat{\mu }^{}\widehat{\iota }_0\widehat{\mathrm{h}}_\delta \mathrm{F}_0)^{}\mathrm{\Theta }`$
$`=`$ $`(\iota _0h_\delta ^{}ȷ_0\mathrm{F}_0)^{}\mathrm{\Theta }(\stackrel{~}{ȷ}_0\stackrel{~}{}_0)^{}\mathrm{\Theta }`$
$`=`$ $`(ȷ_0\mathrm{F}_0)^{}\mathrm{\Theta }_{h_\delta ^{}}(\stackrel{~}{ȷ}_0\stackrel{~}{}_0)^{}\mathrm{\Theta }`$
$`=`$ $`\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}}^{}\mathrm{\Theta }=_{}^{}+\mathrm{\Theta }_{}\mathrm{\Theta }_{}=_{}^{}`$
and the result follows as a consequence of the third item of Proposition 15.
The statement for the splittings of $`\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0`$ and $`\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0`$ is immediate.
Note that, once again, the use of both extended Legendre maps is necessary to obtain the Hamiltonian density in this way.
Finally, in the almost-regular case, the Hamilton-Jacobi variational principle of definition 17 is stated in the same way, now using sections of $`\overline{\rho }_0^1:PM`$, and the form $`\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0`$. So we look for sections $`\psi _0\mathrm{\Gamma }_c(M,P)`$ which are stationary with respect to the variations given by $`\psi _{0t}=\sigma _t\psi _0`$, where $`\{\sigma _t\}`$ is a local one-parameter group of any $`\overline{\rho }_0^1`$-vertical vector field $`Z\text{X}(P)`$, such that
$$\frac{\mathrm{d}}{\mathrm{d}t}|_{t=0}_M\psi _{0t}^{}\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0=0$$
Then these critical sections will be characterized by the condition (analogous to Theorem 4)
$$\psi _0^{}𝑖(\stackrel{~}{X}^0)\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0=0;\text{for every }\stackrel{~}{X}^0\text{X}(P)$$
(29)
### 3.6 Equivalence between the Lagrangian and Hamiltonian formalisms
One expects that both the Lagrangian and Hamiltonian formalism must be equivalent. As in mechanics, this equivalence can be proved by using the (reduced) Legendre map.
First, using the Legendre map, we can lift sections of $`\pi `$ from $`E`$ to $`\mathrm{\Pi }`$ as follows:
###### Definition 20
Let $`(J^1E,\mathrm{\Omega }_{})`$ be a hyper-regular Lagrangian system, $`\mathrm{F}`$ the induced Legendre transformation, $`\varphi :ME`$ a section of $`\pi `$ and $`j^1\varphi :MJ^1E`$ its canonical prolongation to $`J^1E`$. The Lagrangian prolongation of $`\varphi `$ to $`\mathrm{\Pi }`$ is the section
$$j^1\varphi :=\mathrm{F}j^1\varphi :M\mathrm{\Pi }$$
If $`(J^1E,\mathrm{\Omega }_{})`$ is an almost-regular Lagrangian system, the Lagrangian prolongation of a section $`\varphi :ME`$ to $`\mathrm{P}`$ is
$$j_0^1\varphi :=\mathrm{F}_0j^1\varphi :M\mathrm{P}$$
###### Theorem 5
(Equivalence theorem for sections) Let $`(J^1E,\mathrm{\Omega }_{})`$ and $`(\mathrm{\Pi },\mathrm{\Omega }_{\mathrm{h}_\delta })`$ be the Lagrangian and Hamiltonian descriptions of a hyper-regular system.
If a section $`\varphi \mathrm{\Gamma }_c(M,E)`$ is a solution of the Lagrangian variational problem (Hamilton principle) then the section $`\psi j^1\varphi :=\mathrm{F}j^1\varphi \mathrm{\Gamma }_c(M,\mathrm{\Pi })`$ is a solution of the Hamiltonian variational problem (Hamilton-Jacobi principle).
Conversely, if a section $`\psi \mathrm{\Gamma }_c(M,\mathrm{\Pi })`$ is a solution of the Hamiltonian variational problem, then the section $`\varphi \rho ^1\psi \mathrm{\Gamma }_c(M,E)`$ is a solution of the Lagrangian variational problem.
( Proof ) Bearing in mind the diagram
$$\begin{array}{ccc}J^1E& \text{}& \mathrm{\Pi }\\ & \text{}& \end{array}$$
(30)
If $`\varphi `$ is a solution of the Lagrangian variational problem then $`(j^1\varphi )^{}𝑖(X)\mathrm{\Omega }_{}=0`$, for every $`X\text{X}(J^1E)`$ (Theorem 1); therefore, as $`\mathrm{F}`$ is a local diffeomorphism,
$`0`$ $`=`$ $`(j^1\varphi )^{}𝑖(X)\mathrm{\Omega }_{}=(j^1\varphi )^{}𝑖(X)(\mathrm{F}^{}\mathrm{\Omega }_{\mathrm{h}_\delta })`$
$`=`$ $`(j^1\varphi )^{}\mathrm{F}^{}(𝑖(\mathrm{F}_{}^1X)\mathrm{\Omega }_{\mathrm{h}_\delta })=(\mathrm{F}j^1\varphi )^{}(𝑖(X^{})\mathrm{\Omega }_{\mathrm{h}_\delta })`$
which holds for every $`X^{}\text{X}(\mathrm{\Pi })`$ and thus, by Theorem 4, $`\psi \mathrm{F}j^1\varphi `$ is a solution of the Hamiltonian variational problem. (This proof holds also for the almost-regular case).
Conversely, let $`\psi \mathrm{\Gamma }_c(M,\mathrm{\Pi })`$ be a solution of the Hamiltonian variational problem. Reversing the above reasoning we obtain that $`(\mathrm{F}^1\psi )^{}𝑖(X)\mathrm{\Omega }_{}=0`$, for every $`X\text{X}(J^1E)`$, and hence $`\sigma \mathrm{F}^1\psi \mathrm{\Gamma }_c(M,J^1E)`$ is a critical section for the Lagrangian variational problem. Then, as we are in the hyper-regular case, $`\sigma `$ must be a holonomic section , , , $`\sigma =j^1\varphi `$, and since (30) is commutative, $`\varphi =\rho ^1\psi \mathrm{\Gamma }_c(M,E)`$, necessarily.
Remarks:
* Observe that every section $`\psi :M\mathrm{\Pi }`$ which is solution of the Hamilton-Jacobi variational principle is necessarily a Lagrangian prolongation of a section $`\varphi :ME`$.
* In the almost-regular case, if $`\varphi `$ is a critical section of the Lagrangian problem, then $`\psi =\mathrm{F}j^1\varphi `$ is a critical section of the Hamiltonian problem. Furthermore, $`\mathrm{F}:j^1\varphi (M)\mathrm{F}(j^1\varphi (M))`$ is a diffeomorfism because sections of $`\overline{\pi }^1`$ are transversal to the fibres of $`\mathrm{F}`$.
As critical sections are integral manifolds of multivector fields , or, what is equivalent, Ehresmann connections , , , then critical sections through different points in the same fiber of $`\mathrm{F}`$ have the same image by $`\mathrm{F}`$.
On the other hand, we can prove the equivalence between the Lagrangian and Hamiltonian formalisms from the variational point of view. First, we need the following lemma:
###### Lemma 4
Let $`\beta \Omega ^m(J^1E)`$ and $`f\mathrm{C}^{\mathrm{}}(J^1E)`$. For every differentiable section $`\varphi :UME`$, the following conditions are equivalent:
1. $`(j^1\varphi )^{}[f(\overline{\pi }^1\omega )]=(j^1\varphi )^{}\beta `$.
2. $`{\displaystyle _{j^1\varphi }}f(\overline{\pi }^1\omega )={\displaystyle _{j^1\varphi }}\beta `$ .
( Proof ) Trivially $`12`$.
Conversely, if we suppose $`1`$ is not true, then there exists one section $`\varphi :UME`$ with $`(j^1\varphi )^{}[f(\overline{\pi }^1\omega )\beta ]0`$ and hence there is $`xU`$ and a closed neighbourhood $`V`$ of $`x`$ in $`U`$ such that, taking $`\gamma :VE`$ with $`\gamma =\varphi |_V`$, then $`{\displaystyle _{j^1\gamma }}[f(\overline{\pi }^1\omega )\beta ]0`$ , so $`2`$ is false.
Now, let $``$ be a connection in $`\pi :EM`$, and let $`_{}^{}=\mathrm{E}_{}^{}(\overline{\pi }^1\omega )`$ be the density of Lagrangian energy associated with $``$ and $``$. Then:
###### Theorem 6
The Lagrangian energy function is the unique function in $`J^1E`$ verifying the following condition: for every section $`\varphi :UME`$,
$$(j^1\varphi )^{}[\mathrm{E}_{}^{}(\overline{\pi }^1\omega )]=(j^1\varphi )^{}(\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}})$$
( Proof ) (Uniqueness): Let $`f`$ and $`g`$ be two functions verifying this condition. Obviously $`(j^1\varphi )^{}[(fg)(\overline{\pi }^1\omega )]=0`$, but $`0=(j^1\varphi )^{}[(fg)(\overline{\pi }^1\omega )]=(fg)(j^1\varphi (x))(\overline{\pi }^1\omega )`$, for every $`xU`$. Hence, $`(fg)(j^1\varphi (x))=0`$, and this implies $`fg=0`$, because every point in $`J^1E`$ is in the image of some section $`j^1\varphi `$.
(Existence): From (24) we obtain
$`(j^1\varphi )^{}(\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}})`$ $`=`$ $`(j^1\varphi )^{}[\mathrm{\Theta }_{}+\mathrm{E}_{}^{}(\overline{\pi }^1\omega )]`$
$`=`$ $`(j^1\varphi )^{}[𝑖(𝒱)\mathrm{d}++\mathrm{E}_{}^{}(\overline{\pi }^1\omega )]=(j^1\varphi )^{}[\mathrm{E}_{}^{}(\overline{\pi }^1\omega )]`$
since $`(j^1\varphi )^{}(𝑖(𝒱)\mathrm{d})=0`$ (as can be proved by using expressions in coordinates). So, the energy function introduced in definition 1 satisfies this condition.
And this result leads to the following consequence:
###### Theorem 7
Let $`(J^1E,\mathrm{\Omega }_{})`$ and $`(\mathrm{\Pi },\mathrm{\Omega }_{\mathrm{h}_\delta })`$ be the Lagrangian and Hamiltonian descriptions of a hyper-regular system. Then, the Hamilton variational principle of the Lagrangian formalism and the Hamilton-Jacobi variational principle of the Hamiltonian formalisms are equivalent. That is, for every section $`\varphi \mathrm{\Gamma }_c(M,E)`$ we have that
$$_{j^1\varphi }=_{j^1\varphi }\mathrm{\Theta }_{\mathrm{h}_\delta }$$
( Proof ) The standpoint is the relation stated in Theorem 6 which, by Lemma 4, is equivalent to
$$_{j^1\varphi }\mathrm{E}_{}^{}(\overline{\pi }^1\omega )=_{j^1\varphi }(\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}})$$
therefore, from this equality, and using (24) and (22), we obtain
$$_{j^1\varphi }=_{j^1\varphi }[\mathrm{F}^{}\mathrm{\Theta }_{h_\delta ^{}}\mathrm{E}_{}^{}(\overline{\pi }^1\omega )]=_{j^1\varphi }\mathrm{\Theta }_{}=_{j^1\varphi }\mathrm{F}^{}\mathrm{\Theta }_{\mathrm{h}_\delta }=_{\mathrm{F}j^1\varphi }\mathrm{\Theta }_{\mathrm{h}_\delta }=_{j^1\varphi }\mathrm{\Theta }_{\mathrm{h}_\delta }$$
It is important to remark the essential role played by the Lagrangian energy function in the proof of this equivalence.
(The above results are generalizations of others in non-autonomous mechanics ).
## 4 Hamiltonian formalism in the restricted multimomentum bundle $`J^1\pi ^{}`$. Relation with the formalism in $`\mathrm{\Pi }`$
### 4.1 Hamiltonian systems
The construction of the Hamiltonian formalism in $`J^1\pi ^{}`$ is posed, for the first time, in , and the particular case of Hamiltonian systems associated with hyper-regular and almost-regular systems is stated in . The procedure is essentially similar to that developed in section 3.1 for $`\mathrm{\Pi }`$. Next we sketch this construction, relating it with the above one in $`\mathrm{\Pi }`$. As we have proved the existence of the canonical diffeomorphism $`\mathrm{\Psi }:\mathrm{\Pi }J^1\pi ^{}`$, we can use it to prove the equivalence between the Hamiltonian formalisms in $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$.
###### Definition 21
Consider the bundle $`\overline{\tau }^1:J^1\pi ^{}M`$.
1. A section $`h_\mu :J^1\pi ^{}\pi `$ of the projection $`\mu `$ is called a Hamiltonian section of $`\mu `$.
2. The differentiable forms
$$\mathrm{\Theta }_{h_\mu }:=h_\mu ^{}\mathrm{\Theta },\mathrm{\Omega }_{h_\mu }:=\mathrm{d}\mathrm{\Theta }_{h_\mu }=h_\mu ^{}\mathrm{\Omega }$$
are called the Hamilton-Cartan $`m`$ and $`(m+1)`$ forms of $`J^1\pi ^{}`$ associated with the Hamiltonian section $`h_\mu `$.
3. The couple $`(J^1\pi ^{},\mathrm{\Omega }_{h_\mu })`$ is said to be a Hamiltonian system.
In a local chart of natural coordinates, a Hamiltonian section is specified by a local Hamiltonian function $`H_{h_\mu }\mathrm{C}^{\mathrm{}}(U)`$, $`UJ^1\pi ^{}`$, such that $`H_{h_\mu }(x^\nu ,y^A,p_A^\nu )(x^\nu ,y^A,p=H_{h_\mu }(x^\gamma ,y^B,p_B^\eta ),p_A^\nu )`$. The local expressions of the Hamilton-Cartan forms associated with $`h_\mu `$ are similar to (11), but changing $`H_{h_\delta }`$ by $`H_{h_\mu }`$, and $`\mathrm{p}_A^\nu `$ by $`p_A^\nu `$.
For Hamiltonian sections of $`\mu `$, we have a similar result to that in Lemma 2, and so, if $`h_\mu ^1,h_\mu ^2`$ are two sections of $`\mu `$, then
$$\mathrm{\Theta }_{h_\mu ^1}\mathrm{\Theta }_{h_\mu ^2}=h_\mu ^1\mathrm{\Theta }h_\mu ^2\mathrm{\Theta }=\tau ^1(h_\mu ^1h_\mu ^2)$$
(31)
is a $`\overline{\tau }^1`$-semibasic $`m`$-form in $`J^1\pi ^{}`$.
###### Definition 22
A $`\overline{\tau }^1`$-semibasic form $`\Omega ^m(J^1\pi ^{})`$ is said to be a Hamiltonian density in $`J^1\pi ^{}`$.
It can be written as $`=\mathrm{H}(\overline{\tau }^1^{}\omega )`$, where $`\mathrm{H}\mathrm{C}^{\mathrm{}}(J^1\pi ^{})`$ is the global Hamiltonian function associated with $``$ and $`\omega `$.
If (31) holds, then the relation between the global Hamiltonian function $`\mathrm{H}`$ associated with $``$, and the local Hamiltonian functions $`H_{h_\mu ^1}`$, $`H_{h_\mu ^2}`$ associated with $`h_\mu ^1`$ and $`h_\mu ^2`$ is $`\mathrm{H}=H_{h_\mu ^1}H_{h_\mu ^2}`$ (in an open set $`U`$).
In this way we have the analogous result as in Theorem 3:
###### Theorem 8
The set of Hamilton-Cartan $`m`$-forms associated with Hamiltonian sections of $`\mu `$ is an affine space modelled on the set of Hamiltonian densities in $`J^1\pi ^{}`$.
Hence, if $`(J^1\pi ^{},\mathrm{\Omega }_{h_\mu })`$ is a Hamiltonian system, we have that every Hamiltonian section $`h_\mu ^{}h_\mu `$ allows us to split globally the Hamilton-Cartan forms as
$$\mathrm{\Theta }_{h_\mu }=\mathrm{\Theta }_{h_\mu ^{}},\mathrm{\Omega }_{h_\mu }=\mathrm{\Omega }_{h_\mu ^{}}+\mathrm{d}$$
The local expressions of these splittings are similar to (14), but changing $`H_{h_\delta ^{}}`$ by $`H_{h_\mu ^{}}`$, and $`\mathrm{p}_A^\nu `$ by $`p_A^\nu `$.
Now, if we have a connection $``$ in $`\pi :EM`$, it induces a linear section $`h_\delta ^{}:\mathrm{\Pi }J^1E^{}`$ of $`\delta `$, and hence there exists another linear section $`h_\mu ^{}:J^1\pi ^{}\pi `$ of $`\mu `$ given by $`h_\mu ^{}\mathrm{\Psi }^1=\iota _0h_\delta ^{}`$ (see for an alternative definition). Then, if $`\mathrm{\Theta }`$ is the canonical $`m`$-form in $`\mathrm{\Omega }^m(\pi )`$, the forms
$$\mathrm{\Theta }_{h_\mu ^{}}:=h_\mu ^{}\mathrm{\Theta }\mathrm{\Omega }^m(J^1\pi ^{}),\mathrm{\Omega }_{h_\mu ^{}}:=\mathrm{d}\mathrm{\Theta }_{h_\mu ^{}}\mathrm{\Omega }^{m+1}(J^1\pi ^{})$$
are the Hamilton-Cartan $`m`$ and $`(m+1)`$ forms of $`J^1\pi ^{}`$ associated with the connection $``$. Of course, a characterization of $`\mathrm{\Theta }_{h_\mu ^{}}`$ can be stated in the same way as in (16). The local expression of these Hamilton-Cartan forms associated with $``$ is similar to (17).
Therefore, given a connection $``$ and a Hamiltonian section $`h_\mu `$, from the above results we have that
$$\tau ^1(h_\mu ^{}h_\mu )=h_\mu ^{}\mathrm{\Theta }h_\mu ^{}\mathrm{\Theta }=\mathrm{\Theta }_{h_\mu ^{}}\mathrm{\Theta }_{h_\mu }:=_{h_\mu }^{}$$
is a Hamiltonian density in $`J^1\pi ^{}`$, which is written as $`_{h_\mu }^{}=\mathrm{H}_{h_\mu }^{}(\overline{\tau }^1^{}\omega )`$, where $`\mathrm{H}_{h_\mu }^{}\mathrm{C}^{\mathrm{}}(J^1\pi ^{})`$ is the global Hamiltonian function associated with $`_{h_\mu }^{}`$ and $`\omega `$. Then, the Hamilton-Cartan forms associated with $`h_\mu `$ split as
$$\mathrm{\Theta }_{h_\mu }=\mathrm{\Theta }_{h_\mu ^{}}_{h_\mu }^{},\mathrm{\Omega }_{h_\mu }=\mathrm{\Omega }_{h_\mu ^{}}+\mathrm{d}_{h_\mu }^{}$$
The local expressions of these splittings are similar to (19), but changing $`\mathrm{H}_{h_\delta }^{}`$ by $`\mathrm{H}_{h_\mu }^{}`$, and $`\mathrm{p}_A^\nu `$ by $`p_A^\nu `$.
If, conversely, we take a connection $``$ and a Hamiltonian density $``$, then making $`h_\mu ^{}`$ we obtain a Hamiltonian section $`h_\mu `$, since $`:J^1\pi ^{}\pi `$ takes values in $`\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M`$. Hence:
###### Proposition 17
A couple $`(h_\mu ,)`$ in $`J^1\pi ^{}`$ is equivalent to a couple $`(,)`$ (that is, given a connection $``$, Hamiltonian sections of $`\mu `$ and Hamiltonian densities in $`J^1\pi ^{}`$ are in one-to-one correspondence).
Bearing in mind this last result, we have another way of obtaining a Hamiltonian system, which consists in giving a couple $`(,)`$. In fact:
###### Proposition 18
Let $``$ be a connection in $`\pi :EM`$, and $``$ a Hamiltonian density. There exists a unique Hamiltonian section $`h_\mu `$ of $`\mu `$ such that
$$\mathrm{\Theta }_{h_\mu }=\mathrm{\Theta }_{h_\mu ^{}},\mathrm{\Omega }_{h_\mu }=\mathrm{d}\mathrm{\Theta }_{h_\mu }=\mathrm{\Omega }_{h_\mu ^{}}+\mathrm{d}$$
(32)
Concerning field equations, observe that diffeomorphisms in $`E`$ (and hence vector fields in $`E`$) can be lifted to $`J^1\pi ^{}`$, for instance, lifting them to $`\mathrm{\Pi }`$ (see definitions 15 and 16), and translating them to $`J^1\pi ^{}`$ using the diffeomorphism $`\mathrm{\Psi }`$. Hence, for a Hamiltonian system $`(J^1\pi ^{},\mathrm{\Omega }_{h_\mu })`$, we can set the Hamilton-Jacobi variational principle as in definition 17 (but with the form $`\mathrm{\Omega }_{h_\mu }`$ instead of $`\mathrm{\Omega }_{h_\delta }`$), and state the same results and comments as in Theorem 4.
Hamiltonian systems in $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$ are equivalent. In fact; as a first result we have:
###### Proposition 19
Let $``$ be a connection in $`\pi :EM`$, and and $``$ Hamiltonian densities in $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$, respectively, such that $`\Psi ^{}\text{}=`$. Then
$$\mathrm{\Psi }^{}\mathrm{\Omega }_{h_\mu ^{}}=\mathrm{\Omega }_{h_\delta ^{}},\mathrm{\Psi }^{}\mathrm{\Omega }_{h_\mu }=\mathrm{\Omega }_{h_\delta }$$
( Proof ) The proof is based in the following fact:
$$\mathrm{\Theta }_{h_\mu ^{}}=h_\mu ^{}\mathrm{\Theta }=(\iota _0h_\delta ^{}\mathrm{\Psi })^{}\mathrm{\Theta }=\mathrm{\Psi }^{}(\iota _0h_\delta ^{})^{}\mathrm{\Theta }=\mathrm{\Psi }^{}\mathrm{\Theta }_{h_\delta ^{}}$$
and the result is immediate.
And therefore, as a direct consequence of Propositions 6 and 19, we can set the relation between the Hamiltonian systems in $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$:
###### Theorem 9
Every Hamiltonian system $`(\mathrm{\Pi },\mathrm{\Omega }_{h_\delta })`$ is equivalent to a Hamiltonian system $`(J^1\pi ^{},\mathrm{\Omega }_{h_\mu })`$, and conversely.
At this point, we can study the relation between the set of connections $``$ in the bundle $`\pi :EM`$, and the sets of linear (Hamiltonian) sections of the projections $`\mu :\pi J^1\pi ^{}`$ and $`\delta :J^1E^{}\mathrm{\Pi }`$:
###### Theorem 10
The map $`h_\mu ^{}`$ is a bijective affine map between the set of connections in the bundle $`\pi :EM`$ and the set of linear (Hamiltonian) sections of the projection $`\mu :\pi J^1\pi ^{}`$ or, what is equivalent, the set of linear (Hamiltonian) sections of the projection $`\delta :J^1E^{}\mathrm{\Pi }`$.
( Proof ) Let the projection $`\mu :\mathrm{\Lambda }^m\mathrm{T}^{}E\mathrm{\Lambda }_1^m\mathrm{T}^{}E/\mathrm{\Lambda }_0^m\mathrm{T}^{}E`$, and the set of linear sections of $`\mu `$
$$\mathrm{\Gamma }(\mu ):=\{\mathrm{}Ł(J^1\pi ^{},\mathrm{\Lambda }_1^m\mathrm{T}^{}E),\mu \mathrm{}=\mathrm{Id}_{J^1\pi ^{}}\}$$
which is an affine bundle modeled on the vector bundle
$`(J^1\pi ^{})^{}\mathrm{\Lambda }_0^m\mathrm{T}^{}E`$ $``$ $`(\pi ^{}\mathrm{T}M\mathrm{V}^{}(\pi )\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M)^{}\mathrm{\Lambda }_0^m\mathrm{T}^{}E`$
$`=`$ $`\pi ^{}\mathrm{T}^{}M\mathrm{V}(\pi )\pi ^{}\mathrm{\Lambda }^m\mathrm{T}M\mathrm{\Lambda }_0^m\mathrm{T}^{}E\pi ^{}\mathrm{T}^{}M\mathrm{V}(\pi )`$
But this last bundle is just the vector bundle on which the affine bundle of the connection forms in $`\pi :EM`$ is modeled. Then the result follows.
Finally, the equivalence with the set of linear (Hamiltonian) sections of the projection $`\delta :J^1E^{}\mathrm{\Pi }`$) is proved by taking into account that every linear section of $`\delta `$ is associated with a connection $``$, since this linear section defines a linear map from $`\mathrm{V}(\pi )`$ to $`\mathrm{T}^{}E`$, and hence a projection $`\mathrm{T}E\mathrm{V}(\pi )`$ (that is, a connection).
### 4.2 Hamiltonian system associated with a hyper-regular Lagrangian system
The procedure is analogous to that in Section 3.4 (see also diagram (7)). Let $`(J^1E,\mathrm{\Omega }_{})`$ be a hyper-regular Lagrangian system, then:
###### Definition 23
Let $`\mathrm{h}_\mu :J^1\pi ^{}\pi `$ be the section of $`\mu `$ given by
$$\mathrm{h}_\mu :=\stackrel{~}{}^1$$
which is a diffeomorphism connecting $`J^1\pi ^{}`$ and $`\stackrel{~}{}(J^1E)`$ (observe that it is just the inverse of $`\mu `$ restricted to $`\stackrel{~}{}(J^1E)`$). We define the Hamilton-Cartan forms
$$\mathrm{\Theta }_{\mathrm{h}_\mu }:=\mathrm{h}_\mu ^{}\mathrm{\Theta };\mathrm{\Omega }_{\mathrm{h}_\mu }:=\mathrm{h}_\mu ^{}\mathrm{\Omega }$$
###### Proposition 20
The Hamilton-Cartan forms satisfy that
$$^{}\mathrm{\Theta }_{\mathrm{h}_\mu }=\mathrm{\Theta }_{},^{}\mathrm{\Omega }_{\mathrm{h}_\mu }=\mathrm{\Omega }_{}$$
Then $`(J^1\pi ^{},\mathrm{\Omega }_{\mathrm{h}_\mu })`$ is the (unique) Hamiltonian system associated with the hyper-regular Lagrangian system $`(J^1E,\mathrm{\Omega }_{})`$.
( Proof ) We have the diagram
$$\begin{array}{ccccc}\text{}& \text{}& \text{}& \text{}& \text{}\end{array}$$
Taking into account the commutativity of this diagram, and proposition 2, we have
$$^{}\mathrm{\Theta }_{\mathrm{h}_\mu }=^{}\mathrm{h}_\mu ^{}\mathrm{\Theta }=\stackrel{~}{}^{}\mathrm{\Theta }=\mathrm{\Theta }_{}$$
and the same result follows for $`\mathrm{\Omega }_{\mathrm{h}_\delta }`$.
Using charts of natural coordinates in $`J^1\pi ^{}`$ and $`\pi `$, and the expression (5) of the Legendre map, we obtain that the local Hamiltonian function $`H_{\mathrm{h}_\mu }`$ representing this Hamiltonian section is
$$H_{\mathrm{h}_\mu }(x^\nu ,y^A,p_A^\nu )=^1\left(v_\nu ^A\frac{\mathrm{\pounds }}{v_\nu ^A}\mathrm{\pounds }\right)=p_A^\nu ^1^{}v_\nu ^A^1^{}\mathrm{\pounds }$$
(33)
and the local expressions of the corresponding Hamilton-Cartan forms are
$`\mathrm{\Theta }_{\mathrm{h}_\mu }`$ $`=`$ $`p_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu (p_A^\nu ^1^{}v_\nu ^A^1^{}\mathrm{\pounds })\mathrm{d}^mx`$
$`\mathrm{\Omega }_{\mathrm{h}_\mu }`$ $`=`$ $`\mathrm{d}p_A^\nu \mathrm{d}y^A\mathrm{d}^{m1}x_\nu +\mathrm{d}(p_A^\nu ^1^{}v_\nu ^A^1^{}\mathrm{\pounds })\mathrm{d}^mx`$
We can construct this Hamiltonian system using connections. Thus, if $``$ is a connection in $`\pi :EM`$, and $`h_\mu ^{}`$ is the linear Hamiltonian section of $`\mu `$ associated with $``$, following the same pattern as in Proposition 11, we can prove:
###### Proposition 21
1. The $`m`$-form $`^{}\mathrm{\Theta }_{h_\mu ^{}}\mathrm{\Theta }_{}`$ is $`\overline{\pi }^1`$-semibasic and
$$^{}\mathrm{\Theta }_{h_\mu ^{}}\mathrm{\Theta }_{}=_{}^{}$$
2. There exists a unique Hamiltonian density $`^{}\mathrm{\Omega }^m(J^1\pi ^{})`$ such that
$$^{}^{}=^{}\mathrm{\Theta }_{h_\mu ^{}}\mathrm{\Theta }_{}=_{}^{}$$
(34)
Then there exists a function $`H^{}\mathrm{C}^{\mathrm{}}(J^1\pi ^{})`$ such that $`^{}=H^{}(\overline{\tau }^1\omega )`$.
$`^{}`$ and $`H^{}`$ are called the Hamiltonian density and the Hamiltonian function associated with the Lagrangian system, the connection $``$ and $`\omega `$.
3. The Hamilton-Cartan forms of definition 23 split as
$`\mathrm{\Theta }_{\mathrm{h}_\mu }`$ $`=`$ $`\mathrm{\Theta }_{h_\mu ^{}}^{}=\mathrm{\Theta }_{h_\mu ^{}}^{}`$
$`\mathrm{\Omega }_{\mathrm{h}_\mu }`$ $`=`$ $`\mathrm{d}\mathrm{\Theta }_{\mathrm{h}_\mu }=\mathrm{\Omega }_{h_\mu ^{}}+\mathrm{d}^{}`$ (35)
We can obtain this Hamiltonian density using only Hamiltonian sections. In fact:
###### Proposition 22
1. Consider the Hamiltonian section $`\mathrm{h}_\mu :=\stackrel{~}{}^1`$, and a connection $``$. Then we have
$$h_\mu ^{}\mathrm{\Theta }\mathrm{h}_\mu ^{}\mathrm{\Theta }=^{}$$
and hence the splitting (35) holds.
( Proof ) We have the diagram
$$\begin{array}{ccccccc}\text{}& \text{}& \text{}& \text{}& \text{}& \text{}& \text{}\end{array}$$
The first item is a consequence of the third item of Proposition 3. For the second item, taking into account definition 8 and (32), Proposition 2, and (34), we have
$$h_\mu ^{}\mathrm{\Theta }\mathrm{h}_\mu ^{}\mathrm{\Theta }=\mathrm{\Theta }_{h_\mu ^{}}(^1)^{}\stackrel{~}{}^{}\mathrm{\Theta }=\mathrm{\Theta }_{h_\mu ^{}}(^1)^{}\mathrm{\Theta }_{}=^{}$$
Then the result for the Hamilton-Cartan forms follows immediately.
Of course, all the results stated in section 3.3 concerning to the variational principle and the characterization of critical sections are true, and the local Hamiltonian function $`H_{\mathrm{h}_\mu }`$ appearing in the Hamilton-De Donder-Weyl equations is given by (33).
Finally, the equivalence between the Hamiltonian formalisms (in $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$) associated with a hyper-regular Lagrangian system, and between the Lagrangian formalism and the Hamiltonian formalism in $`J^1\pi ^{}`$ is given by the following:
###### Theorem 11
Let $`(J^1E,\mathrm{\Omega }_{})`$ be a hyper-regular Lagrangian system. Then
$$\mathrm{\Psi }^{}\mathrm{\Theta }_{\mathrm{h}_\delta }=\mathrm{\Theta }_{\mathrm{h}_\mu },\mathrm{\Psi }^{}\mathrm{\Omega }_{\mathrm{h}_\delta }=\mathrm{\Omega }_{\mathrm{h}_\mu }$$
and hence its associated Hamiltonian systems $`(\mathrm{\Pi },\mathrm{\Omega }_{\mathrm{h}_\delta })`$ and $`(J^1\pi ^{},\mathrm{\Omega }_{\mathrm{h}_\mu })`$ are equivalent.
( Proof ) It is immediate. Observe also that the sections $`\mathrm{h}_\mu `$ and $`\mathrm{h}_\delta `$ are equivalent, by the commutativity of diagram (7).
Observe that, as $``$ is a diffeomorphism, we also have that $`\mathrm{\Psi }^{}\text{}^{}=^{}`$.
### 4.3 Hamiltonian system associated with an almost-regular Lagrangian system
The procedure is analogous to that in Section 3.5 (see also diagram (9)). Now, $`(J^1E,\mathrm{\Omega }_{})`$ is an almost-regular Lagrangian system, and the submanifold $`ȷ_0:𝒫J^1\pi ^{}`$, is a fiber bundle over $`E`$ (and $`M`$). The corresponding projections will be denoted by $`\tau _0^1:𝒫E`$ and $`\overline{\tau }_0^1:𝒫M`$, satisfying that $`\tau ^1ȷ_0=\tau _0^1`$ and $`\overline{\tau }^1ȷ_0=\overline{\tau }_0^1`$.
Taking into account Proposition 5, and following the same pattern as in Propositions 13, we can prove that the Lagrangian forms $`\mathrm{\Theta }_{}`$ and $`\mathrm{\Omega }_{}`$ are $``$-projectable, and then:
###### Definition 24
Given the diffeomorphism $`\stackrel{~}{\mathrm{h}}_\mu =\stackrel{~}{\mu }^1`$, we define the Hamilton-Cartan forms
$$\mathrm{\Theta }_{\stackrel{~}{\mathrm{h}}_\mu }^0=\stackrel{~}{\mathrm{h}}_\mu ^{}\mathrm{\Theta };\mathrm{\Omega }_{\stackrel{~}{\mathrm{h}}_\mu }^0=\stackrel{~}{\mathrm{h}}_\mu ^{}\mathrm{\Omega }$$
###### Proposition 23
The Hamilton-Cartan forms satisfy that
$$_0^{}\mathrm{\Theta }_{\stackrel{~}{\mathrm{h}}_\mu }^0=\mathrm{\Theta }_{},_0^{}\mathrm{\Omega }_{\stackrel{~}{\mathrm{h}}_\mu }^0=\mathrm{\Omega }_{}$$
Then $`(J^1\pi ^{},𝒫,\mathrm{\Omega }_{\stackrel{~}{\mathrm{h}}_\mu }^0)`$ is the (unique) Hamiltonian system associated with the almost-regular Lagrangian system $`(J^1E,\mathrm{\Omega }_{})`$.
( Proof ) In fact, taking into account the commutativity of diagram (9), and proposition 2, we have
$$_0^{}\mathrm{\Theta }_{\widehat{\mathrm{h}}_\mu }^0=_0^{}(\stackrel{~}{ȷ}_0\widehat{\mathrm{h}}_\mu )^{}\mathrm{\Theta }=(\stackrel{~}{ȷ}_0\widehat{\mathrm{h}}_\mu _0)^{}\mathrm{\Theta }=(\stackrel{~}{ȷ}_0\stackrel{~}{}_0)^{}\mathrm{\Theta }=\stackrel{~}{}^{}\mathrm{\Theta }=\mathrm{\Theta }_{}$$
and the same result follows for $`\mathrm{\Omega }_{\mathrm{h}_\mu }^0`$.
We can construct this Hamiltonian system using a connection. Thus, let $``$ be a connection in $`\pi :EM`$, and $`h_\mu ^{}:J^1\pi ^{}\pi `$ the induced linear section of $`\mu `$. Let $`_{}^{}\mathrm{\Omega }^m(J^1E)`$ be the density of Lagrangian energy associated with $``$. Then, as in Proposition 15 we can prove:
###### Proposition 24
1. The $`\overline{\pi }^1`$-semibasic $`m`$-form $`^{}\mathrm{\Theta }_{h_\mu ^{}}\mathrm{\Theta }_{}`$ is $``$-projectable and, if $`\mathrm{\Theta }_{h_\mu ^{}}^0=ȷ_0^{}\mathrm{\Theta }_{h_\mu ^{}}`$, then
$$_{}^{}=^{}\mathrm{\Theta }_{h_\mu ^{}}\mathrm{\Theta }_{}=_0^{}\mathrm{\Theta }_{h_\mu ^{}}^0\mathrm{\Theta }_{}$$
2. There exists a unique $`\overline{\tau }_0^1`$-semibasic form $`_0^{}\mathrm{\Omega }^m(𝒫)`$, such that $`_0^{}_0^{}=_{}^{}`$. Then, there is a function $`H_0^{}\mathrm{C}^{\mathrm{}}(𝒫)`$ such that $`_0^{}=H_0^{}(\overline{\tau }_0^1\omega )`$. Obviously we have that $`_0^{}H_0^{}=\mathrm{E}_{}^{}`$.
$`_0^{}`$ and $`H_0^{}`$ are called the Hamiltonian density and the Hamiltonian function associated with the Lagrangian system, the connection $``$ and $`\omega `$.
3. The Hamilton-Cartan forms of definition 24 split as
$`\mathrm{\Theta }_{\stackrel{~}{\mathrm{h}}_\mu }^0`$ $`=`$ $`ȷ_0^{}\mathrm{\Theta }_{h_\mu ^{}}_0^{}=\mathrm{\Theta }_{h_\mu ^{}}^0_0^{}`$
$`\mathrm{\Omega }_{\stackrel{~}{\mathrm{h}}_\mu }^0`$ $`=`$ $`\mathrm{d}\mathrm{\Theta }_{\stackrel{~}{\mathrm{h}}_\mu }^0=ȷ_0^{}\mathrm{\Omega }_{h_\mu ^{}}+\mathrm{d}_0^{}=\mathrm{\Omega }_{h_\mu ^{}}^0+\mathrm{d}_0^{}`$ (36)
We can obtain this Hamiltonian system in the following equivalent way:
###### Proposition 25
Consider the map $`\stackrel{~}{\mathrm{h}}_\mu `$, and a connection $``$. Then $`h_\mu ^{}`$ induces a map $`\stackrel{~}{h}_\mu ^{}:𝒫\stackrel{~}{𝒫}`$, defined by the relation $`\stackrel{~}{ȷ}_0\stackrel{~}{h}_\mu ^{}=h_\mu ^{}ȷ_0`$. Therefore
$$(\stackrel{~}{ȷ}_0\stackrel{~}{h}_0^{})^{}\mathrm{\Theta }(\stackrel{~}{ȷ}_0\stackrel{~}{\mathrm{h}}_\delta )^{}\mathrm{\Theta }=_0^{}$$
and hence the splitting (36) holds.
( Proof ) We have the diagram
$$\begin{array}{cccccc}\text{}& \text{}& \text{}& \text{}& \text{}& \text{}\end{array}$$
The first part of the statement is a consequence of the fact that $`\stackrel{~}{\mu }`$ is a diffeomorphism. For the second part, taking into account the commutativity of this diagram, and bearing in mind the first item of Proposition 24 and Proposition 2, we have
$`_0^{}_0^{}`$ $`=`$ $`_0^{}[(\stackrel{~}{ȷ}_0\stackrel{~}{h}_0^{})^{}\mathrm{\Theta }(\stackrel{~}{ȷ}_0\stackrel{~}{\mathrm{h}}_\delta )^{}\mathrm{\Theta }]=_0^{}\mathrm{\Theta }_{h_\mu ^{}}^0_0^{}\stackrel{~}{\mathrm{h}}_\mu ^{}ȷ_0^{}\mathrm{\Theta }`$
$`=`$ $`_0^{}\mathrm{\Theta }_{h_\mu ^{}}^0\stackrel{~}{}^{}\mathrm{\Theta }=_{}^{}+\mathrm{\Theta }_{}\mathrm{\Theta }_{}=_{}^{}`$
and the result follows as a consequence of the above Proposition. The last statement is immediate.
Of course, the result stated in (29) concerning to the variational principle and the characterization of critical sections holds in the same way.
Finally, the equivalence between the Hamiltonian formalisms (in $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$) associated with an almost-regular Lagrangian system, and between the Lagrangian formalism and the Hamiltonian formalism in $`J^1\pi ^{}`$ is given by the following:
###### Theorem 12
Let $`(J^1E,\mathrm{\Omega }_{})`$ be an almost-regular Lagrangian system. Then
$$\mathrm{\Psi }_0^{}\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0=\mathrm{\Theta }_{\stackrel{~}{\mathrm{h}}_\mu }^0,\mathrm{\Psi }_0^{}\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0=\mathrm{\Omega }_{\stackrel{~}{\mathrm{h}}_\mu }^0$$
and hence its associated Hamiltonian systems $`(\mathrm{\Pi },P,\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0)`$ and $`(J^1\pi ^{},𝒫,\mathrm{\Omega }_{\stackrel{~}{\mathrm{h}}_\mu }^0)`$ are equivalent.
( Proof ) First, we have the following relation
$$\mathrm{\Theta }_{h_\mu ^{}}^0=ȷ_0^{}\mathrm{\Theta }_{h_\mu ^{}}=ȷ_0^{}\mathrm{\Psi }^{}\mathrm{\Theta }_{h_\delta ^{}}=(\mathrm{\Psi }ȷ_0)^{}\mathrm{\Theta }_{h_\delta ^{}}=(ȷ_0\mathrm{\Psi }_0)^{}\mathrm{\Theta }_{h_\delta ^{}}=\mathrm{\Psi }_0^{}ȷ_0^{}\mathrm{\Theta }_{h_\delta ^{}}=\mathrm{\Psi }_0^{}\mathrm{\Theta }_{h_\delta ^{}}^0$$
Furthermore, as $`\mathrm{F}_0=\mathrm{\Psi }_0_0`$, we have that
$$_{}^{}=_0^{}_0^{}=\mathrm{F}_0^{}\text{}_0^{}=(\mathrm{\Psi }_0_0)^{}\text{}_0^{}=_0^{}\mathrm{\Psi }_0^{}\text{}_0^{}\mathrm{\Psi }_0^{}\text{}_0^{}=_0^{}$$
since $`_0`$ is a submersion. (In the same way $`\mathrm{\Psi }_0^{}\mathrm{H}_0^{}=H_0^{}`$). Therefore, $`\mathrm{\Psi }_0^{}\mathrm{\Theta }_{\widehat{\mathrm{h}}_\delta }^0=\mathrm{\Theta }_{\stackrel{~}{\mathrm{h}}_\mu }^0`$, and hence $`\mathrm{\Psi }_0^{}\mathrm{\Omega }_{\widehat{\mathrm{h}}_\delta }^0=\mathrm{\Omega }_{\stackrel{~}{\mathrm{h}}_\mu }^0`$.
Observe also that the section $`\stackrel{~}{\mathrm{h}}_\mu `$ of $`\mu `$ and the family of sections $`\widehat{\mathrm{h}}_\delta `$ of $`\delta `$ (given in Proposition 16) are equivalent, by the commutativity of diagram (9).
## 5 Examples
### 5.1 Non-autonomous Mechanics
The jet bundle description of time-dependent mechanical systems (see, for instance, and ) takes $`M=\text{}`$, and $`E=\text{}\times Q`$, where $`Q`$ is a $`N`$-dimensional manifold (and thus $`\pi :\text{}\times Q\text{}`$ is a trivial bundle). Then $`J^1E=\text{}\times \mathrm{T}Q`$. Natural adapted coordinates are denoted by $`(t,q^i,v^i)`$. Lagrangian densities are written $`=\mathrm{\pounds }\mathrm{d}t`$, where $`\mathrm{\pounds }\mathrm{C}^{\mathrm{}}(\text{}\times \mathrm{T}Q)`$ is a time-dependent Lagrangian function.
Next we will identify the different multimomentum bundles. First observe that, as $`\mathrm{T}M\text{}\times \text{}`$, we obtain
$$\mathrm{\Lambda }^m\mathrm{T}^{}M\mathrm{\Lambda }^1\mathrm{T}^{}M\text{}\times \text{}^{}$$
Therefore we have:
Generalized multimomentum bundle: Observe that
$`\pi ^{}\mathrm{T}M`$ $`=`$ $`\pi ^{}(\text{}\times \text{})(\text{}\times Q)\times _{\text{}}(\text{}\times \text{})\text{}\times Q\times \text{}`$
$`\mathrm{T}^{}E`$ $`=`$ $`\mathrm{T}^{}(\text{}\times Q)=\mathrm{T}^{}\text{}\times \mathrm{T}^{}Q`$
$`\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M`$ $``$ $`\text{}\times Q\times \text{}^{}Q\times \mathrm{T}^{}\text{}`$
and hence
$`J^1E^{}`$ $`:=`$ $`\pi ^{}\mathrm{T}M\times _E\mathrm{T}^{}E\times _E\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M=(\text{}\times Q\times \text{})\times _{(\text{}\times Q)}(\mathrm{T}^{}\text{}\times \mathrm{T}^{}Q)\times _{(\text{}\times Q)}(Q\times \mathrm{T}^{}\text{})`$
$``$ $`(\text{}\times Q\times \text{})\times _{(\text{}\times Q)}(\mathrm{T}^{}\text{}\times \mathrm{T}^{}Q)\times _{(\text{}\times Q)}((\text{}\times Q)\times \text{}^{})\mathrm{T}^{}\text{}\times \mathrm{T}^{}Q\mathrm{T}^{}(\text{}\times Q)`$
Notice that $`dimJ^1E^{}=2N+2`$.
Reduced multimomentum bundle: As $`\mathrm{V}^{}(\pi )=\text{}\times \mathrm{T}^{}Q`$, we have that
$`\mathrm{\Pi }`$ $`:=`$ $`\pi ^{}\mathrm{T}M\times _E\mathrm{V}^{}(\pi )\times _E\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M=(\text{}\times Q\times \text{})\times _{(\text{}\times Q)}(\text{}\times \mathrm{T}^{}Q)\times _{(\text{}\times Q)}(Q\times \mathrm{T}^{}\text{})`$
$``$ $`(\text{}\times Q\times \text{})\times _{(\text{}\times Q)}(\mathrm{T}^{}\text{}\times \mathrm{T}^{}Q)\times _{(\text{}\times Q)}((\text{}\times Q)\times \text{}^{})\text{}\times \mathrm{T}^{}Q`$
and $`dim\mathrm{\Pi }=2N+1`$.
Extended multimomentum bundle: Now we have
$$\pi :=\mathrm{\Lambda }_1^m\mathrm{T}^{}E\mathrm{\Lambda }_1^1\mathrm{T}^{}E=\mathrm{\Lambda }_1^1\mathrm{T}^{}E\mathrm{T}^{}E\mathrm{T}^{}(\text{}\times Q)\mathrm{T}^{}\text{}\times \mathrm{T}^{}Q$$
with $`dim\pi =2N+2`$.
Restricted multimomentum bundle: Observe that $`\mathrm{\Lambda }_0^m\mathrm{T}^{}E\mathrm{\Lambda }_0^1\mathrm{T}^{}E=(\text{}\times Q)\times \text{}^{}`$ and then
$$J^1\pi ^{}:=\pi /\mathrm{\Lambda }_0^m\mathrm{T}^{}E(\mathrm{T}^{}\text{}\times \mathrm{T}^{}Q)/(\text{}\times Q\times \text{}^{})\text{}\times \mathrm{T}^{}Q$$
with $`dimJ^1\pi ^{}=2N+1`$.
Comments:
* It is interesting to point out that in the Hamiltonian formalism of non-autonomous mechanics, $`M\text{}\times \mathrm{T}^{}Q`$ and $`J^1\pi ^{}\text{}\times \mathrm{T}^{}Q`$ make the canonical diffeomorphism between the generalized and the extended multimomentum bundle evident. They correspond to the so-called extended momentum phase space of the symplectic formulation of time-dependent systems , , . As a consequence, the generalized and the (first) extended Legendre maps are really the same.
* The case of singular (almost-regular) time-dependent mechanical systems has been extensively studied in this context in , .
### 5.2 Electromagnetic field (with fixed background)
In this case $`M`$, is space-time endowed with a semi-Riemannian metric $`g`$, $`E=\mathrm{T}^{}M`$ is a vector bundle over $`M`$ and $`\pi :\mathrm{T}^{}MM`$ denotes the natural projection. Sections of $`\pi `$ are the so-called electromagnetic potentials. Using the linear connection associated with the metric $`g`$, one assures that $`J^1E\mathrm{T}^{}M`$ is a vector bundle and, since $`\mathrm{V}E=\pi ^{}\mathrm{T}^{}M`$, we have $`J^1E=\pi ^{}\mathrm{T}^{}M_E\pi ^{}\mathrm{T}^{}M`$.
Let $`\varphi :M\mathrm{T}^{}M`$ be a section of $`\pi `$. Then $`j^1\varphi :M\pi ^{}\mathrm{T}^{}M\times \pi ^{}\mathrm{T}^{}M`$ is just $`j^1\varphi =\mathrm{T}\varphi `$ (observe that $`j^1\varphi `$ is a metric tensor on $`M`$). Now, considering $`\overline{y}J^1E`$, and $`\varphi :M\mathrm{T}^{}M`$ being a representative of $`\overline{y}`$; we have that the Lagrangian density is
$$=\frac{1}{4}\mathrm{d}\varphi ^2\mathrm{d}V_g$$
where $``$ denotes the norm induced by the metric $`g`$ on the $`2`$-forms on $`M`$, and $`\mathrm{d}V_g`$ is the volume element associated with the metric $`g`$. Observe that $`\mathrm{d}\varphi `$ is the skew-symmetric part of the matrix giving $`\mathrm{T}\varphi `$ or, in other words, the skew-symmetric part of the metric $`\mathrm{T}\varphi `$ on $`M`$.
For simplifying calculations, we take $`M=\text{}^3`$ and the metric is $`++`$. Then $`E=\text{}^3\times \text{}^3`$ and
$$J^1E=(\text{}^3\times \text{}^3)\times (\text{}^3\times \text{}^3)$$
with $`dimJ^1E=15`$. Coordinates in $`J^1E`$ are usually denoted $`(x^\nu ,A^j,v_\nu ^j)`$, with $`\nu ,j=0,1,2`$. The coordinates $`(A^1,A^2)`$ constitute the vector potential and $`A^0`$ is the scalar potential. Then, locally $`\varphi =\varphi _\eta \mathrm{d}x^\eta `$, and therefore $`j^1\varphi ={\displaystyle \frac{\varphi _\eta }{x^\nu }}\mathrm{d}x^\nu \mathrm{d}x^\eta `$ . It is usual to write $`\varphi `$ in the form $`A=\delta _{j\nu }A^j\mathrm{d}x^\nu `$ and $`\mathrm{d}A=\delta _{j\eta }v_\nu ^j\mathrm{d}x^\eta \mathrm{d}x^\nu `$ ($`\delta _{j\eta }`$ is the Kronecker’s delta). Then, in natural coordinates we have the following expression for the Lagrangian function
$$\mathrm{\pounds }=\frac{1}{4}[(v_2^1v_1^2)^2(v_0^2v_2^0)^2(v_0^1v_1^0)^2]$$
Obviously, this is a singular Lagrangian, since its Hessian matrix
$$\frac{^2\mathrm{\pounds }}{v_\nu ^iv_\eta ^j}=\frac{1}{2}\left(\begin{array}{ccccccccc}0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 1& 0& 0\\ 0& 1& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 1& 0\\ 0& 0& 1& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0\end{array}\right)$$
is singular (its rank is equal to 3).
Next we study the several Legendre maps associated with this system. As we know, all of them leave the coordinates $`(x^\nu ,A^j)`$ invariant, thus we will write the relations only for the multimomentum coordinates.
Generalized Legendre map: The generalized multimomentum bundle is
$$J^1E^{}=(\text{}^3\times \text{}^3)\times (\text{}^3(\text{}^3\text{}^3)\mathrm{\Lambda }^3\text{}^3)$$
From (1) we have
$$\begin{array}{ccccccccc}\widehat{\mathrm{F}}^{}\mathrm{p}_0^0& =& 0& \widehat{\mathrm{F}}^{}\mathrm{p}_1^0& =& \frac{1}{2}(v_0^1v_1^0)& \widehat{\mathrm{F}}^{}\mathrm{p}_2^0& =& \frac{1}{2}(v_0^2v_2^0)\\ \widehat{\mathrm{F}}^{}\mathrm{p}_0^1& =& \frac{1}{2}(v_0^1v_1^0)& \widehat{\mathrm{F}}^{}\mathrm{p}_1^1& =& 0& \widehat{\mathrm{F}}^{}\mathrm{p}_2^1& =& \frac{1}{2}(v_2^1v_1^2)\\ \widehat{\mathrm{F}}^{}\mathrm{p}_0^2& =& \frac{1}{2}(v_0^2v_2^0)& \widehat{\mathrm{F}}^{}\mathrm{p}_1^2& =& \frac{1}{2}(v_2^1v_1^2)& \widehat{\mathrm{F}}^{}\mathrm{p}_2^2& =& 0\end{array}$$
and for the additional multimomentum coordinates ($`\mathrm{p}_\eta ^\nu `$ in (1))
$$\begin{array}{ccc}\widehat{\mathrm{F}}^{}\widehat{\mathrm{p}}_0^0=\frac{1}{2}(v_0^1(v_0^1v_1^0)+v_0^2(v_0^2v_2^0))& & \widehat{\mathrm{F}}^{}\widehat{\mathrm{p}}_1^0=\frac{1}{2}(v_1^1(v_0^1v_1^0)+v_1^2(v_0^2v_2^0))\\ \widehat{\mathrm{F}}^{}\widehat{\mathrm{p}}_2^0=\frac{1}{2}(v_2^1(v_0^1v_1^0)+v_2^2(v_0^2v_2^0))& & \widehat{\mathrm{F}}^{}\widehat{\mathrm{p}}_0^1=\frac{1}{2}(v_0^0(v_0^1v_1^0)+v_0^2(v_2^1v_1^2))\\ \widehat{\mathrm{F}}^{}\widehat{\mathrm{p}}_1^1=\frac{1}{2}(v_1^0(v_0^1v_1^0)+v_1^2(v_2^1v_1^2))& & \widehat{\mathrm{F}}^{}\widehat{\mathrm{p}}_2^1=\frac{1}{2}(v_2^0(v_0^1v_1^0)+v_2^2(v_2^1v_1^2))\\ \widehat{\mathrm{F}}^{}\widehat{\mathrm{p}}_0^2=\frac{1}{2}(v_0^0(v_0^2v_2^0)v_0^1(v_2^1v_1^2))& & \widehat{\mathrm{F}}^{}\widehat{\mathrm{p}}_1^2=\frac{1}{2}(v_1^0(v_0^2v_2^0)v_1^1(v_2^1v_1^2))\\ \widehat{\mathrm{F}}^{}\widehat{\mathrm{p}}_2^2=\frac{1}{2}(v_2^0(v_0^2v_2^0)v_2^1(v_2^1v_1^2))& & \end{array}$$
We have the Hamiltonian constraints
$$\begin{array}{ccc}\widehat{\xi }^1\mathrm{p}_0^0=0& \widehat{\xi }^2\mathrm{p}_1^1=0& \widehat{\xi }^3\mathrm{p}_2^2=0\\ \widehat{\xi }^4\mathrm{p}_1^0+\mathrm{p}_0^1=0& \widehat{\xi }^5\mathrm{p}_2^0+\mathrm{p}_0^2=0& \widehat{\xi }^6\mathrm{p}_2^1+\mathrm{p}_1^2=0\end{array}$$
and the additional ones
$`\widehat{\xi }^7`$ $``$ $`\mathrm{p}_0^2\mathrm{p}_1^2(\widehat{\mathrm{p}}_0^0+\widehat{\mathrm{p}}_1^12(\widehat{\mathrm{p}}_0^1)^2))+(\mathrm{p}_0^2)^2\widehat{\mathrm{p}}^1_0+(\mathrm{p}_1^2)^2\widehat{\mathrm{p}}^0_1+\mathrm{p}^1_0(\mathrm{p}_1^2\widehat{\mathrm{p}}_1^2\mathrm{p}_0^2\widehat{\mathrm{p}}_0^2)=0`$
$`\widehat{\xi }^8`$ $``$ $`\mathrm{p}_0^1\mathrm{p}_1^2(\widehat{\mathrm{p}}_0^0+\widehat{\mathrm{p}}_2^22(\widehat{\mathrm{p}}_0^2)^2))+(\mathrm{p}_0^1)^2\widehat{\mathrm{p}}^2_0+(\mathrm{p}_1^2)^2\widehat{\mathrm{p}}^0_2+\mathrm{p}^2_0(\mathrm{p}_1^2\widehat{\mathrm{p}}_2^1\mathrm{p}_0^1\widehat{\mathrm{p}}_0^1)=0`$
$`\widehat{\xi }^9`$ $``$ $`\mathrm{p}_0^1\mathrm{p}_0^2(\widehat{\mathrm{p}}_1^1\widehat{\mathrm{p}}_2^22(\widehat{\mathrm{p}}_1^2)^2))+(\mathrm{p}_0^2)^2\widehat{\mathrm{p}}^1_2(\mathrm{p}_0^1)^2\widehat{\mathrm{p}}^2_1+\mathrm{p}^2_1(\mathrm{p}_0^2\widehat{\mathrm{p}}_2^0\mathrm{p}_0^1\widehat{\mathrm{p}}_1^0)=0`$
which define locally the submanifold $`\widehat{P}`$ in $`J^1E^{}`$.
Observe that $`dim\widehat{P}=dimJ^1E`$, as $`\mathrm{rank}\widehat{\mathrm{F}}_{}`$ is maximal. Then, taking into account the commutativity of diagram (9), we can conclude that, for this system, the degeneracy is on the fibers of the projection $`\delta `$.
Reduced Legendre map: The reduced multimomentum bundle is
$$\mathrm{\Pi }=(\text{}^3\times \text{}^3)\times (\text{}^3\text{}^3\mathrm{\Lambda }^3\text{}^3)$$
For the reduced Legendre map the results are the same as for the restricted Legendre map, but changing the multimomentum coordinates $`p_\nu ^j`$ by $`\mathrm{p}_\nu ^j`$. So, from (2) we obtain
$$\begin{array}{ccccccccc}\mathrm{F}^{}\mathrm{p}_0^0& =& 0& \mathrm{F}^{}\mathrm{p}_1^0& =& \frac{1}{2}(v_0^1v_1^0)& \mathrm{F}^{}\mathrm{p}_2^0& =& \frac{1}{2}(v_0^2v_2^0)\\ \mathrm{F}^{}\mathrm{p}_0^1& =& \frac{1}{2}(v_0^1v_1^0)& \mathrm{F}^{}\mathrm{p}_1^1& =& 0& \mathrm{F}^{}\mathrm{p}_2^1& =& \frac{1}{2}(v_2^1v_1^2)\\ \mathrm{F}^{}\mathrm{p}_0^2& =& \frac{1}{2}(v_0^2v_2^0)& \mathrm{F}^{}\mathrm{p}_2^2& =& \frac{1}{2}(v_2^1v_1^2)& \mathrm{F}^{}\mathrm{p}_2^2& =& 0\end{array}$$
($`\mathrm{F}`$ is a submersion onto its image, and hence the system is almost-regular).
Now we have the same Hamiltonian constraints
$$\begin{array}{ccc}\xi ^1\mathrm{p}_0^0=0& \xi ^2\mathrm{p}_1^1=0& \xi ^3\mathrm{p}_2^2=0\\ \xi ^4\mathrm{p}_1^0+\mathrm{p}_0^1=0& \xi ^5\mathrm{p}_2^0+\mathrm{p}_0^2=0& \xi ^6\mathrm{p}_2^1+\mathrm{p}_1^2=0\end{array}$$
which define locally the submanifold $`P`$ in $`\mathrm{\Pi }`$. Observe that, as these constraints are conditions of skew-symmetry, we have that
$$\mathrm{P}=\mathrm{F}(J^1E)=\pi ^{}(𝒜(\text{}^3)\mathrm{\Lambda }^3(\text{}^3))$$
where $`𝒜(\text{}^3)`$ denotes the bundle whose sections are the $`2`$-contravariant skew-symmetric tensor fields on $`\text{}^3`$. Note also that $`\widehat{P}`$ is not diffeomorphic to $`P`$.
First extended Legendre map: The extended multimomentum bundle is
$$\pi =(\text{}^3\times \text{}^3)\times \mathrm{\Lambda }_1^3(\text{}^3\times \text{}^3)$$
From (4) we obtain that
$$\begin{array}{ccccccccc}\widehat{}^{}p_0^0& =& 0& \widehat{}^{}p_1^0& =& \frac{1}{2}(v_0^1v_1^0)& \widehat{}^{}p_2^0& =& \frac{1}{2}(v_0^2v_2^0)\\ \widehat{}^{}p_0^1& =& \frac{1}{2}(v_0^1v_1^0)& \widehat{}^{}p_1^1& =& 0& \widehat{}^{}p_2^1& =& \frac{1}{2}(v_2^1v_1^2)\\ \widehat{}^{}p_0^2& =& \frac{1}{2}(v_0^2v_2^0)& \widehat{}^{}p_1^2& =& \frac{1}{2}(v_2^1v_1^2)& \widehat{}^{}p_2^2& =& 0\end{array}$$
and the additional relation
$$\widehat{}^{}p=\frac{1}{2}[(v_2^1v_1^2)^2(v_0^2v_2^0)^2(v_0^1v_1^0)^2]2\mathrm{\pounds }$$
The corresponding Hamiltonian constraints are
$$\begin{array}{ccc}\widehat{\chi }^1p_0^0=0& \widehat{\chi }^2p_1^1=0& \widehat{\chi }^3p_2^2=0\\ \widehat{\chi }^4p_1^0+p_0^1=0& \widehat{\chi }^5p_2^0+p_0^2=0& \widehat{\chi }^6p_2^1+p_1^2=0\end{array}$$
and the additional one
$$\widehat{\chi }^7p+2(p_1^2)^22(p_0^2)^22(p_0^1)^2=0$$
All of them define locally the submanifold $`\widehat{𝒫}`$ in $`\pi `$.
Second extended Legendre map: From (4) we obtain that
$$\begin{array}{ccccccccc}\stackrel{~}{}^{}p_0^0& =& 0& \stackrel{~}{}^{}p_1^0& =& \frac{1}{2}(v_0^1v_1^0)& \stackrel{~}{}^{}p_2^0& =& \frac{1}{2}(v_0^2v_2^0)\\ \stackrel{~}{}^{}p_0^1& =& \frac{1}{2}(v_0^1v_1^0)& \stackrel{~}{}^{}p_1^1& =& 0& \stackrel{~}{}^{}p_2^1& =& \frac{1}{2}(v_2^1v_1^2)\\ \stackrel{~}{}^{}p_0^2& =& \frac{1}{2}(v_0^2v_2^0)& \stackrel{~}{}^{}p_1^2& =& \frac{1}{2}(v_2^1v_1^2)& \stackrel{~}{}^{}p_2^2& =& 0\end{array}$$
and the additional relation
$$\stackrel{~}{}^{}p=\frac{1}{4}[(v_2^1v_1^2)^2(v_0^2v_2^0)^2(v_0^1v_1^0)^2]\mathrm{\pounds }$$
and the Hamiltonian constraints are now
$$\begin{array}{ccc}\stackrel{~}{\chi }^1p_0^0=0& \stackrel{~}{\chi }^2p_1^1=0& \stackrel{~}{\chi }^3p_2^2=0\\ \stackrel{~}{\chi }^4p_1^0+p_0^1=0& \stackrel{~}{\chi }^5p_2^0+p_0^2=0& \stackrel{~}{\chi }^6p_2^1+p_1^2=0\end{array}$$
and the additional one
$$\stackrel{~}{\chi }^7p+(p_1^2)^2(p_0^2)^2(p_0^1)^2=0$$
All of them define locally the submanifold $`\stackrel{~}{𝒫}`$ in $`\pi `$. Note that the last constraint identifies the extra coordinate $`p`$ with the Hamiltonian function which, for this system, is $`H={\displaystyle \frac{1}{4}}((p_1^2)^2(p_0^2)^2(p_0^1)^2)`$ .
Restricted Legendre map: The restricted multimomentum bundle is
$$J^1\pi ^{}=(\text{}^3\times \text{}^3)\times [\mathrm{\Lambda }_1^3(\text{}^3\times \text{}^3)/\mathrm{\Lambda }_0^3(\text{}^3\times \text{}^3)]$$
From (5) we obtain
$$\begin{array}{ccccccccc}^{}p_0^0& =& 0& ^{}p_1^0& =& \frac{1}{2}(v_0^1v_1^0)& ^{}p_2^0& =& \frac{1}{2}(v_0^2v_2^0)\\ ^{}p_0^1& =& \frac{1}{2}(v_0^1v_1^0)& ^{}p_1^1& =& 0& ^{}p_2^1& =& \frac{1}{2}(v_2^1v_1^2)\\ ^{}p_0^2& =& \frac{1}{2}(v_0^2v_2^0)& ^{}p_1^2& =& \frac{1}{2}(v_2^1v_1^2)& ^{}p_2^2& =& 0\end{array}$$
Hence we obtain the following set of Hamiltonian constraints
$$\begin{array}{ccc}\chi ^1p_0^0=0& \chi ^2p_1^1=0& \chi ^3p_2^2=0\\ \chi ^4p_1^0+p_0^1=0& \chi ^5p_2^0+p_0^2=0& \chi ^6p_2^1+p_1^2=0\end{array}$$
which define locally the submanifold $`𝒫`$ in $`J^1\pi ^{}`$.
Observe that, for this example,
$$\mathrm{rank}\mathrm{F}_{}=\mathrm{rank}\widehat{}_{}=\mathrm{rank}\stackrel{~}{}_{}=\mathrm{rank}_{}$$
and the submanifolds $`P,𝒫,\widehat{𝒫}`$ and $`\stackrel{~}{𝒫}`$ are, in fact, diffeomorphic.
## 6 Conclusions
We have studied the Hamiltonian formalism for first-order Classical Field theories in the context of multisymplectic manifolds, taking different choices of multimomentum bundles as phase spaces, in particular the bundles $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$.
* First we have reviewed the construction of these and other auxiliary multimomentum bundles ($`J^1E^{}`$ and $`\pi `$), as well as the definition of suitable Legendre maps when all these bundles are thought of, in a certain sense, as the dual bundles of a Lagrangian system $`(J^1E,\mathrm{\Omega }_{})`$. The key result is the existence of a canonical diffeomorphism between $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$. (See section 2.2).
* In order to state the Hamiltonian formalism on $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$, some additional geometric element is needed for obtaining the Hamilton-Cartan forms from the canonical forms which $`J^1E^{}`$ and $`\pi `$ are endowed with. In particular, we can take sections of the projections $`\delta :J^1E^{}\mathrm{\Pi }`$ and $`\mu :\pi J^1\pi ^{}`$, (which are called Hamiltonian sections, and are the elements carrying the “physical” information in this construction), which allows us to pull-back the canonical forms from $`J^1E^{}`$ and $`\pi `$ to $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$ respectively. These are the Hamilton-Cartan forms which define the Hamiltonian system. Hamiltonian sections are associated with local Hamiltonian functions, which appear explicitly in the local expression of the corresponding Hamilton-Cartan forms. (See sections 3.1 and 4.1).
* A relevant result is that different choices of Hamiltonian sections of $`\delta `$ may lead to the same Hamilton-Cartan forms in $`\mathrm{\Pi }`$, and this allows us to establish an equivalence relation in the set of sections of the projection $`\delta `$. Then, using the diffeomorphism $`\mathrm{\Psi }`$ between $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$, it is proved that there is a one-to-one correspondence between sections of $`\mu `$ and classes of equivalent sections of $`\delta `$. Therefore, the Hamilton-Cartan forms in $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$ are $`\mathrm{\Psi }`$-related and hence, Hamiltonian systems in $`\mathrm{\Pi }`$ and $`J^1\pi ^{}`$ are equivalent. (See sections 3.1 and 4.1). Furthermore, another one-to-one correspondence exists between the set of connections in the bundle $`EM`$, and the set of linear sections of the respective projections $`\delta `$ and $`\mu `$. (See sections 3.2 and 4.1).
* The difference between two Hamilton-Cartan $`m`$-forms defined by Hamiltonian sections is a semibasic $`m`$-form which is called a Hamiltonian density. Hence the set of Hamilton-Cartan forms can be thought of as an affine space modelled on the module of Hamiltonian densities. As a particular case, Given a connection, (classes of) Hamiltonian sections and Hamiltonian densities are in one-to-one correspondence. As a consequence of this fact, a Hamiltonian system can be also constructed starting from a Hamiltonian density and a connection. (See sections 3.2 and 4.1).
* The field equations of the Hamiltonian formalism can be derived from the so-called Hamilton-Jacobi variational principle. Different but equivalent ways of characterizing the critical sections by means of the Hamilton-Cartan forms are shown. In particular, in natural coordinates of the multimomentum bundles $`\mathrm{\Pi }`$ or $`J^1\pi ^{}`$, these sections are obtained as solutions of a local system of first-order partial differential equations, which are known as the Hamilton-De Donder-Weyl equations. Nevertheless, as the Hamiltonian function appearing in the local expression of the Hamilton-Cartan form is local, these equations are not covariant. Then, for obtaining a set of covariant equations, we must introduce a Hamiltonian density, and so a global Hamiltonian function. (See section 3.3).
* The question of associating a Hamiltonian system to a Lagrangian one is also analyzed, both in the hyper-regular and the almost-regular cases. We can define this Hamiltonian system in three equivalent ways (see sections 3.4, 3.5, 4.2 and 4.3):
+ Using a natural Hamiltonian section, which is defined using the Legendre maps, for obtaining the Hamilton-Cartan forms. These forms are related to the Poincaré-Cartan forms of the Lagrangian formalism, through the Legendre map.
+ Using a connection, the density of Lagrangian energy of the Lagrangian formalism can be defined. Then, we construct a Hamiltonian density as the only semibasic $`m`$-form which is related to it by means of the suitable Legendre map.
+ This last Hamiltonian density can be obtained from a connection and the above natural Hamiltonian section. In this case, the extended Legendre maps must also be used, and in particular, for the construction in the reduced multimomentum bundle $`\mathrm{\Pi }`$, both extended Legendre maps are needed (in the hyper-regular and in the almost-regular cases). This fact would justify the introduction of two extended Legendre maps.
* As an additional result, in the hyper-regular case, the equivalence between the Lagrangian and the Hamiltonian formalism is proved from a double point of view: showing the equivalence between the sections solution of the Lagrangian and Hamiltonian problems, and proving the equivalence of the Lagrangian and Hamiltonian variational principles. This equivalence is only partially proved in the almost-regular case. (See section 3.6).
## Appendix A Geometrical structures in first-order jet bundles
(See , and ).
Let $`\pi :EM`$ be a fiber bundle ($`dimM=m`$, $`dimE=N+m`$), and $`\pi ^1:J^1EE`$ the 1-jet bundle of local sections of $`\pi `$, which is also a differentiable bundle on $`M`$ with projection $`\overline{\pi }^1=\pi `$. If $`(x^\nu ,y^A)`$ (with $`\nu =1,\mathrm{},m`$; $`A=1,\mathrm{},N`$) is a local system of coordinates adapted to the bundle $`\pi :EM`$, then we denote by $`(x^\nu ,y^A,v_\nu ^A)`$ the local system of coordinates induced in $`J^1E`$.
Let $`\varphi :ME`$ be a section of $`\pi `$, $`xM`$ and $`y=\varphi (x)`$. The vertical differential of the section $`\varphi `$ at the point $`yE`$ is the map
$$\begin{array}{ccccc}\mathrm{d}_y^v\varphi & :& \mathrm{T}_yE& & \mathrm{V}_y(\pi )\\ & & u& & u\mathrm{T}_y(\varphi \pi )u\end{array}$$
Then, considering $`\overline{y}J^1E`$ with $`\overline{y}\stackrel{\pi ^1}{}y\stackrel{\pi }{}x`$ and $`\overline{u}\mathrm{T}_{\overline{y}}J^1E`$. The structure canonical $`1`$-form of $`J^1E`$, denoted by $`\theta `$, is defined by
$$\theta (\overline{y};\overline{u}):=(\mathrm{d}_y^v\varphi )(\mathrm{T}_{\overline{y}}\pi ^1(\overline{u}))$$
where $`\varphi `$ is a representative of $`\overline{y}`$. Its expression in a natural local system is $`\theta =(\mathrm{d}y^Av_\nu ^A\mathrm{d}x^\nu ){\displaystyle \frac{}{y^A}}`$ . $`\theta `$ is an element of $`\mathrm{\Omega }^1(J^1E)_{J^1E}\mathrm{\Gamma }(J^1E,\pi ^1\mathrm{V}(\pi ))`$; then it can be thought of as a $`\mathrm{C}^{\mathrm{}}(J^1E)`$-linear map $`\theta :\mathrm{\Gamma }(J^1E,\pi ^1\mathrm{V}(\pi ))^{}\mathrm{\Omega }^1(J^1E)`$.
Consider the canonical isomorphism $`𝒮_{\overline{y}}:\mathrm{T}_{\overline{\pi }^1(\overline{y})}^{}M\mathrm{V}_{\pi ^1(\overline{y})}(\pi )\mathrm{V}_{\overline{y}}(\pi ^1)`$ which consists in associating to an element $`\alpha v\mathrm{T}_{\overline{\pi }^1(\overline{y})}^{}M\mathrm{V}_{\pi ^1(\overline{y})}(\pi )`$ the directional derivative in $`\overline{y}`$ with respect to $`\alpha v`$. Taking into account that $`\alpha v`$ acts in $`J_y^1E`$ by translation, we have
$$𝒮_{\overline{y}}(\alpha v):=D_{\alpha v}(\overline{y}):f\underset{t0}{lim}\frac{f(\overline{y}+t(\alpha v))f(\overline{y})}{t}$$
for $`f\mathrm{C}^{\mathrm{}}(J_y^1E)`$. Then we have the following isomorphism of $`\mathrm{C}^{\mathrm{}}(J^1E)`$-modules
$$𝒮:\mathrm{\Gamma }(J^1E,\overline{\pi }^1\mathrm{T}^{}M\pi ^1\mathrm{V}(\pi ))\mathrm{\Gamma }(J^1E,\mathrm{V}(\pi ^1))$$
which is called the vertical endomorphism $`𝒮`$. ($`\mathrm{\Gamma }(A,B)`$ denotes the set of sections of the projection $`AB`$). Notice that $`𝒮\mathrm{\Gamma }(J^1E,(\pi ^1\mathrm{V}(\pi ))^{}\mathrm{V}(\pi ^1)\overline{\pi }^1\mathrm{T}M)`$ (where all the tensor products are on $`\mathrm{C}^{\mathrm{}}(J^1E)`$). Then, another vertical endomorphism $`𝒱`$ arises from the natural contraction between the factor $`\mathrm{\Gamma }(J^1E,(\pi ^1\mathrm{V}(\pi ))^{})`$ of $`𝒮`$ and the factor $`\mathrm{\Gamma }(J^1E,\pi ^1\mathrm{V}(\pi ))`$ of $`\theta `$:
$$𝒱=𝑖(𝒮)\theta \mathrm{\Omega }^1(J^1E)\mathrm{\Gamma }(J^1E,\mathrm{V}(\pi ^1)\overline{\pi }^1\mathrm{T}M)$$
so it is a morphism
$$𝒱:\mathrm{\Gamma }(J^1E,\mathrm{V}^{}(\pi ^1)\overline{\pi }^1\mathrm{T}^{}M)\mathrm{\Omega }^1(J^1E)$$
$`𝒮`$ can also be thought of as a morphism
$$𝒮:\mathrm{\Gamma }(J^1E,\mathrm{V}^{}(\pi ^1)\overline{\pi }^1\mathrm{T}^{}M)\mathrm{\Gamma }(J^1E,(\overline{\pi }^1\mathrm{V}(\pi ))^{})$$
As every connection $``$ on $`\pi :EM`$ gives an injection $`^v:\mathrm{\Gamma }(J^1E,(\overline{\pi }^1\mathrm{V}(\pi ))^{})\mathrm{\Omega }^1(J^1E)`$, then it makes sense to define
$$𝒮^{}:=^v𝒮:\mathrm{\Gamma }(J^1E,\mathrm{V}^{}(\pi ^1)\overline{\pi }^1\mathrm{T}^{}M)\mathrm{\Omega }^1(J^1E)$$
As a consequence of the foregoing, the operation $`𝒮^{}𝒱`$ is meaningful.
In a natural system of coordinates the local expressions of all these elements are
$`𝒮`$ $`=`$ $`\zeta ^A{\displaystyle \frac{}{v_\nu ^A}}{\displaystyle \frac{}{x^\nu }}`$
$`𝒱`$ $`=`$ $`\left(\mathrm{d}y^Av_\nu ^A\mathrm{d}x^\nu \right){\displaystyle \frac{}{v_\eta ^A}}{\displaystyle \frac{}{x^\eta }}`$
$`𝒮^{}`$ $`=`$ $`\left(\mathrm{d}y^A\mathrm{\Gamma }_\nu ^A\mathrm{d}x^\nu \right){\displaystyle \frac{}{v_\eta ^A}}{\displaystyle \frac{}{x^\eta }}`$
where $`\{\zeta ^A\}`$ is the local basis of $`\mathrm{\Gamma }(J^1E,\pi ^1\mathrm{V}(\pi ))^{}`$ which is dual of $`\left\{{\displaystyle \frac{}{y^A}}\right\}`$ , and $`\mathrm{\Gamma }_\nu ^A`$ are the component functions of the connection $``$.
### Acknowledgments
We are grateful for the financial support of the CICYT TAP97-0969-C03-01 and the CICYT PB98-0821. We wish to thank Mr. Jeff Palmer for his assistance in preparing the English version of the manuscript.
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# Naked Singularity and Gauss-Bonnet Term in Brane World Scenarios
## I Introduction
It was suggested long time ago that the cosmological constant problem could be solved if one is willing to go to higher dimensional spacetime. Recently, a number of authors have outlined possibilities for solving this long standing problem within the brane world scenario . In , our 4-dimensional world is embedded in a 5-dimensional universe and remains flat in the presence of an arbitrary vacuum energy density $`V.`$ The price one pays for this remarkable phenomenon is the appearance of naked singularities in the 5-dimensional universe.
In this note, we consider adding to the action higher derivative terms, in particular the Gauss-Bonnet term, and/or a potential for the dilaton field. A simple anti-de Sitter solution and inflationary solutions to the action with Gauss-Bonnet term were studied in . Models with gravity coupled to scalars are considered in , and in the context of five dimensional supergravity as well. Here we explore the question of how the solution of the cosmological constant problem might be modified by the addition of these terms.
To introduce the subject and to set the stage for our discussion, let us review the proposed solution reduced to its simplest version. We will follow which we will refer to as KSS. The low energy effective action is taken to be
$$S=d^5x\sqrt{G}\left[R\frac{4}{3}(\phi )^2\right]+d^4x\sqrt{g}(V).$$
(1)
Gravity and a dilaton field $`\phi `$ live in the 5-dimensional world and is coupled to a thin 4-dimensional brane whose position is taken to be at $`y=0.`$ Here $`M,N=0,1,2,3,5`$ and $`\mu ,\nu =0,1,2,3.`$ The metric $`g_{\mu \nu }=\delta _\mu ^M\delta _\nu ^NG_{MN}(y=0)`$ is the 4-dimensional metric on the brane. (We use the convention in which $`R_{NPQ}^M=+\mathrm{\Gamma }_{NQ,P}^M\mathrm{}`$, $`R_{MN}=R_{MPN}^P`$, $`R=R_M^M`$, and the signature $`(,+,+,+,+).`$)
Einstein’s equations read
$$R_{MN}\frac{1}{2}G_{MN}R\frac{4}{3}\left[_M\phi _N\phi \frac{1}{2}G_{MN}(\phi )^2\right]+\frac{1}{2}\sqrt{\frac{g}{G}}Vg_{\mu \nu }\delta _M^\mu \delta _N^\nu \delta (y)=0$$
(2)
with $`g`$ and $`G`$ the determinant of $`g_{\mu \nu }`$ and $`G_{MN}`$ respectively. The equation of motion for the dilaton field is given by
$$^2\phi =0.$$
(3)
A theorist living in the 4-dimensional brane world would notice that the action contains a cosmological constant given by the vacuum energy density $`V.`$ KSS showed, however, that for any $`V`$ there exists a solution to (2) and (3) with the metric taking the form
$$ds^2=e^{2A(y)}(dt^2+dx_1^2+dx_2^2+dx_3^2)+dy^2.$$
(4)
In other words, the 4-dimensional world is flat in spite of the presence of $`V.`$ The dilaton field adjusts itself so that this solution exists. KSS referred to this as a self-tuning solution of the cosmological constant problem.
Arithmetically, this is possible because the equations of motion are solved separately for $`y>0`$ and $`y<0`$ and we have enough arbitrary integration constants that we can adjust in order to get the flat 4-dimensional world (4) we want. The equations of motion are, where denotes differentiation with respect to $`y`$,
$$\frac{4}{9}(\phi ^{})^2+A^{\prime \prime }=\frac{1}{6}V\delta (y),$$
(5)
$$(A^{})^2\frac{1}{9}(\phi ^{})^2=0,$$
(6)
and
$$\phi ^{\prime \prime }+4A^{}\phi ^{}=0.$$
(7)
(5) is the difference between the $`\mu \mu `$ and $`55`$ components of Einstein’s equation, (6) the $`55`$ component (which does not receive a brane contribution proportional to $`\delta (y)`$), and (7) the dilaton equation of motion. The 5-dimensional Bianchi identity states that only two of the three equations are independent. For example, (5) and (6) imply (7). Solving (6) for $`A^{}`$ and inserting into (7) we see immediately that $`\phi `$ is given by the logarithm of $`y.`$
Thus, the solution may be chosen to be, for $`y_+>y>0`$
$$\phi (y)=\frac{3}{4}\mathrm{log}(y_+y)+d_+$$
(8)
and
$$A(y)=\frac{1}{4}\mathrm{log}(y_+y)+a_+,$$
(9)
and for $`y_{}<y<0`$
$$\phi (y)=\frac{3}{4}\mathrm{log}(y+y_{})+d_{}$$
(10)
and
$$A(y)=\frac{1}{4}\mathrm{log}(y+y_{})+a_{}$$
(11)
Here $`a_+,a_,d_{+,}`$ $`d_{},`$ $`y_+`$ and $`y_{}`$are integration constants. The continuity of $`\phi `$ and $`A`$ at $`y=0`$ determines $`d_+`$ and $`a_+`$ in terms of the other constants. Integrating (5) and (7) across $`y=0`$ gives the jump conditions that
$$A^{}(y=0^+)A^{}(y=0^{})=\frac{1}{4}\left(\frac{1}{y_+}\frac{1}{y_{}}\right)=\frac{1}{6}V$$
(12)
and
$$\phi ^{}(y=0^+)\phi ^{}(y=0^{})=\frac{3}{4}\left(\frac{1}{y_+}\frac{1}{y_{}}\right)=0$$
(13)
These two equations merely fix $`y_+`$ and $`y_{}`$ in terms of $`V.`$ Thus, KSS obtained a flat space solution, type I in their classification, for any $`V`$.
One can say that, in some sense, the cosmological constant problem is solved because we have a lot of integration constants. Of course, the framework of embedding our universe in a larger universe is of crucial importance.
The heavy price that one pays in the KSS solution is the appearance of naked singularities at $`y_+`$ and $`y_{}.`$ Near $`y_+`$ for example, the metric components $`g_{00}=g_{ii}=e^{2A(y)}`$ vanish like $`\sqrt{y_+y}`$. Various curvature invariants, for example the Ricci scalar $`R=20(A^{})^28A^{\prime \prime }`$ diverge. These singularities are not clothed by an event horizon. That some types of naked singularity are not acceptable is known as cosmic censorship. Recently, Gubser has studied in detail various singularities.
For metrics of the form $`ds^2=e^{2A(y)}g_{\mu \nu }(x)dx^\mu dx^\nu +dy^2`$ we can see easily that $`\sqrt{G}R=e^{2A(y)}\sqrt{g}R_{(4)}+\mathrm{}`$ where the 4-dimensional scalar curvature $`R_{(4)}`$ is constructed out of $`g_{\mu \nu }(x).`$ Thus, the effective 4-dimensional Planck mass squared is proportional to $`𝑑ye^{2A(y)}`$. In order to have a finite 4-dimensional Planck mass, KSS were forced to choose $`y_+`$ and $`y_{}`$ positive. They (rather arbitrarily) postulated that the universe is cutoff at $`y_+`$ and $`y_{}`$ and thus obtained a finite Planck mass squared proportional to
$$𝑑ye^{2A(y)}=_0^{y_+}𝑑ye^{2A(y)}+_y_{}^0𝑑ye^{2A(y)}.$$
(14)
Once these integrals are thus cutoff to give finite values, we can obtain any desired value for the 4-dimensional Planck mass squared simply by shifting the additive integration constant allowed by (9) and (11) in the solution for $`A(y).`$
Notice that while changing $`V`$ does not appear to affect our 4-dimensional universe, which remains flat whatever the value of $`V,`$ it does affect the 5-dimensional universe. In particular, it moves the naked singularities around.
This completes a necessarily brief review of KSS to which we refer for further details. We should perhaps mention here that KSS showed that more generally we can write $`f(\phi )`$ instead of $`V`$ and the same conclusion continues to hold. This is easy to see: in (5) $`V`$ is replaced by $`f(\phi )`$ and in (7) the term $`+\frac{3}{8}\delta (y)\frac{f}{\phi }`$ is added to the right hand side. A particularly popular choice is $`f(\phi )=Ve^{b\phi }`$ as inspired by string theory.
At this point, it may be useful to say a few words about fine tuning and self tuning. Traditionally, one is to write down in a field theory Lagrangian all terms with dimension less than or equal to 4 allowed by symmetry. The coefficients of all such terms are to be regarded as free parameters. If we arbitrarily set one coefficient equal to another or to zero without being able to invoke a symmetry principle, then we are said to have fine tuned. We refer to this as the strong definition of fine tuning. Here we are dealing with a non-renormalizable 5-dimensional action (1) which nobody would regard as fundamental. Presumably, this action is to be regarded as the low energy effective action of some more fundamental theory such as string theory. A traditionalist would say that KSS, starting with the action (1), has already committed fine tuning according to the strong definition by excluding all sorts of possible terms from (1). We feel that it is necessary to formulate a weak definition of fine tuning, according to which one can exclude or include any sorts of terms from the higher dimensional action without being accused of fine tuning. Thus, we would regard the KSS solution as not a fine tuning solution. Rather, KSS described their solution as self tuning, in the sense that the 4-dimensional world remains flat regardless of the value of $`V.`$
For example, the 5-dimensional action (1) could perfectly well contain the term $`\mathrm{\Lambda }e^{a\phi }`$ corresponding to a bulk cosmological constant $`\mathrm{\Lambda }`$, which KSS set to $`0`$ in their self-tuning solutions. A traditionalist using the strong definition of fine tuning would definitely call setting $`\mathrm{\Lambda }`$ to $`0`$ fine tuning; in the language of this subject, however, this is not called fine tuning. If $`\mathrm{\Lambda }`$ is not equal to $`0,`$ then as KSS themselves showed, the self tuning feature of their solution is lost. To obtain a flat 4-dimensional world, one has to adjust $`V`$ to have a value determined by $`\mathrm{\Lambda }`$. KSS called this fine tuning, as we think anybody would.
## II Action with Gauss-Bonnet term
We ask whether we could avoid naked singularities by adding higher derivative terms such as $`R^2`$ to the Einstein-Hilbert action. At the same time, of course we still have the highly non-trivial constraint that $`_0^{\mathrm{}}𝑑ye^{2A(y)}`$ and $`_{\mathrm{}}^0𝑑ye^{2A(y)}`$ have to be finite.
We replace the bulk action in (1) by
$$S_{bulk}=d^5x\sqrt{G}\left[R+aR^24bR^{MN}R_{MN}+cR^{MNPQ}R_{MNPQ}\frac{4}{3}(\phi )^2𝒱(\phi )\right],$$
(15)
We have also included a potential $`𝒱(\phi )`$ for the dilaton field.
We will proceed in stages. We will first study the effect of the higher derivative terms with $`𝒱(\phi )=0.`$ Then in section III we will include $`𝒱(\phi ).`$
Varying $`S_{bulk}`$ with respect to the metric tensor we obtain the equation of motion
$`\{R_{MN}{\displaystyle \frac{1}{2}}G_{MN}R{\displaystyle \frac{1}{2}}G_{MN}[aR^24bR^{PQ}R_{PQ}+cR^{PQST}R_{PQST}]`$ (18)
$`+2aRR_{MN}4cR_{MP}R_N^P+2cR_{MPQS}R_N^{PQS}+4(2bc)R^{PQ}R_{MPQN}`$
$`+2(ab)G_{MN}R_{;P}^{;P}2(a2b+c)R_{;M;N}4(bc)R_{MN;P}^{;P}\}`$
$`=`$ $`T_{MN},`$ (19)
where
$$T_{MN}=\frac{4}{3}\left[_M\phi _N\phi \frac{1}{2}G_{MN}(\phi )^2\right]\frac{1}{2}G_{MN}𝒱(\phi ).$$
(20)
These higher derivative terms naturally arise in the low energy effective action of string theory. For our purposes, we do not inquire of their deeper origin but simply treat the action as “phenomenological.” (For $`a=b=0`$, (18) agrees with .)
We will choose $`a=b=c=\lambda `$ so that the added term $`\lambda (R^24R^{MN}R_{MN}+R^{MNPQ}R_{MNPQ})`$ is of the Gauss-Bonnet form. One sees from the equation of motion that all terms involving fourth derivative vanish in this case and (18) reduces to the result given in . As is well known, the Gauss-Bonnet combination is a topological invariant in $`4`$dimensional spacetime. In higher dimensional spacetime, it is not a topological invariant but nevertheless has particularly attractive properties, as discovered by Zwiebach and explained by Zumino. We proceed in an exploratory spirit we believe appropriate for this stage of development of this nascent subject and do not apologize further for this specific choice.
With the ansatz for the metric in (4), we obtain
$$\frac{4}{9}(\phi ^{})^2+\left[14\lambda (A^{})^2\right]A^{\prime \prime }=\frac{1}{6}f(\phi )\delta (y),$$
(21)
$$(A^{})^2\frac{1}{9}(\phi ^{})^22\lambda (A^{})^4=0,$$
(22)
and
$$\phi ^{\prime \prime }+4A^{}\phi ^{}=\frac{3}{8}\delta (y)\frac{f}{\phi }(\phi ).$$
(23)
The matching condition at the location of the 3-brane, $`y=0`$, becomes
$`{\displaystyle \frac{8}{3}}\phi ^{}(y)|_0^{}^{0^+}`$ $`=`$ $`{\displaystyle \frac{f}{\phi }}\left(\phi (0)\right),`$ (24)
$`6\left[A^{}(y){\displaystyle \frac{4}{3}}\lambda A^{}(y)^3\right]|_0^{}^{0^+}`$ $`=`$ $`f(\phi (0)),`$ (25)
and the continuity condition for $`\phi (y)`$ and $`A(y)`$ at $`y=0`$. In the limit $`\lambda 0`$ these equations reduce to the ones studied in KSS of course. Again, the Bianchi identity assures us that only two out of three bulk equations are independent. After solving two equations we can use the third one as a convenient check.
We solve these equations in the bulk, for $`y>0`$ say. The scalar equation (23) gives
$$\phi ^{}=de^{4A(y)}$$
(26)
with $`d`$ an integration constant. Inserting this into (22) we obtain
$$(A^{})^2=\frac{1}{4\lambda }\left(1\pm \sqrt{1\frac{8d^2}{9}\lambda e^{8A}}\right)$$
(27)
From (27) we see immediately that, for $`\lambda >0`$, we must have $`e^{8A}>(8d^2/9)\lambda `$ and hence the 4-D Planck mass squared is infinite unless spacetime is cut off by some singularity. To solve (27), we change variable to $`\kappa (y)=e^{4A}=1/G_{tt}^2`$ and (27) becomes
$$(\kappa ^{})^2=\frac{4\kappa ^2}{\lambda }\left(1\pm \sqrt{1\frac{8d^2}{9}\lambda \kappa ^2}\right),$$
(28)
which can be readily solved.
For $`\lambda >0`$, we obtain
$`y(\kappa )`$ $`=`$ $`y_0\pm {\displaystyle \frac{1}{\sqrt{8\lambda }}}\left[\mathrm{log}\left(\mathrm{cot}{\displaystyle \frac{\theta +\pi }{4}}\right)\mathrm{csc}{\displaystyle \frac{\theta }{2}}\right],`$ (29)
$`y(\kappa )`$ $`=`$ $`y_0\pm {\displaystyle \frac{1}{\sqrt{8\lambda }}}\left[\mathrm{log}\left(\mathrm{tan}{\displaystyle \frac{\theta }{4}}\right)\mathrm{sec}{\displaystyle \frac{\theta }{2}}\right],`$ (30)
where
$$\mathrm{sin}\theta =\frac{2d}{3}\sqrt{2\lambda }\kappa =\frac{2d}{3}\sqrt{2\lambda }\frac{1}{G_{tt}^2},$$
(31)
$`\pi /2\theta \pi /2`$ and $`y_0`$ is an integration constant. These four solutions correspond to the four different choices of signs in (28). In order to have finite 4-D Planck mass, there are again naked singularities in the fifth dimension.
For $`\lambda =|\lambda |<0`$, only the plus sign in (28) is allowed and we have
$$y(\kappa )=y_0\pm \frac{1}{\sqrt{8|\lambda |}}\left[2\mathrm{tan}^1\left(\mathrm{tanh}\frac{\psi }{4}\right)\text{csch}\left(\frac{\psi }{2}\right)\right],$$
(32)
where
$$\mathrm{sinh}\psi =\frac{2d}{3}\sqrt{2|\lambda |}\kappa .$$
(33)
In this case, the naked singularity persists if we demand finite 4-D Planck mass. In conclusion, while we still have the nice self tuning solution of the cosmological constant problem as in KSS we are also stuck with the naked singularities.
## III ADDING A DILATON POTENTIAL
We next explore what happens if in the action we add a potential $`𝒱(\phi )`$ for the scalar field. A potential can arise in string theory from higher order corrections, but again in the “phenomenological” spirit of this paper we do not concern ourselves with its origin. The equations of motion are modified to
$$\frac{4}{9}(\phi ^{})^2+\left[14\lambda (A^{})^2\right]A^{\prime \prime }=\frac{1}{6}f(\phi )\delta (y),$$
(34)
$$(A^{})^2\frac{1}{9}(\phi ^{})^22\lambda (A^{})^4=\frac{1}{12}𝒱(\phi ),$$
(35)
and
$$\phi ^{\prime \prime }+4A^{}\phi ^{}=\frac{3}{8}\frac{𝒱(\phi )}{\phi }+\frac{3}{8}\delta (y)\frac{f}{\phi }(\phi )$$
(36)
When $`\lambda =0`$, these equations have been studied previously and a solution-generating method inspired by supergravity was suggested in . It is shown that, for $`\lambda =0`$, if $`𝒱(\phi )`$ takes the special form
$$𝒱(\phi )=\frac{27}{4}\left(\frac{W(\phi )}{\phi }\right)^212W(\phi )^2,$$
(37)
then a solution to
$`\phi ^{}`$ $`=`$ $`{\displaystyle \frac{9}{4}}{\displaystyle \frac{W(\phi )}{\phi }}`$ (38)
$`A^{}`$ $`=`$ $`W(\phi ),`$ (39)
is also a solution to the equations of motion. Furthermore, by counting the number of integration constants one can show that the solution space of (37)-(39) coincides with the solution space of the equations of motion following from (34)-(36) for $`\lambda =0`$. $`W(\phi )`$ is sometimes called the “superpotential” for obvious reason, though no supersymmetry is involved here. Note that in this case
$$A^{\prime \prime }=\frac{9}{4}\left(\frac{W(\phi )}{\phi }\right)^20.$$
(40)
From (34) and (35), staying in the bulk and thus ignoring $`\delta (y)`$ for now, we obtain similar first order equations for nonzero $`\lambda `$:
$`𝒱(\phi )`$ $`=`$ $`\left[{\displaystyle \frac{27}{4}}\left({\displaystyle \frac{W(\phi )}{\phi }}\right)^2+{\displaystyle \frac{3}{2\lambda }}\right]\left(14\lambda W(\phi )^2\right)^2{\displaystyle \frac{3}{2\lambda }},`$ (41)
$`A^{}`$ $`=`$ $`W(\phi ),`$ (42)
$`\phi ^{}`$ $`=`$ $`{\displaystyle \frac{9}{4}}\left(14\lambda W(\phi )^2\right){\displaystyle \frac{W(\phi )}{\phi }}.`$ (43)
In terms of the superpotential $`W(\phi )`$, the matching condition at the location of the 3-brane, $`y=0`$, is now
$`6{\displaystyle \frac{}{\phi }}\left[W(\phi (y)){\displaystyle \frac{4}{3}}\lambda W(\phi (y))^3\right]|_{y=0^{}}^{y=0^+}`$ $`=`$ $`{\displaystyle \frac{f}{\phi }}\left(\phi (0)\right),`$ (44)
$`6\left[W(\phi (y)){\displaystyle \frac{4}{3}}\lambda W(\phi (y))^3\right]|_{y=0^{}}^{y=0^+}`$ $`=`$ $`f(\phi (0)).`$ (45)
If, in a specific model, the 3-brane tension $`f(\phi )`$ is given by $`12(W\frac{4}{3}\lambda W^3)`$, these jump conditions can be satisfied identically.
It is simple to check that (36) is satisfied automatically. Generalization to $`n`$ scalar fields $`𝝋`$$`=(\phi _1,\mathrm{},\phi _n)`$ is achieved by the following replacement:
$`W(\phi )`$ $``$ $`W(𝝋),`$ (46)
$`\left({\displaystyle \frac{W(\phi )}{\phi }}\right)^2`$ $``$ $`{\displaystyle \frac{W(𝝋)}{𝝋}}{\displaystyle \frac{W(𝝋)}{𝝋}}.`$ (47)
In the limit $`\lambda 0`$, these equations reduces to (37)-(39). Note that the dilaton potential $`𝒱(\phi )`$ is now bounded below if $`\lambda `$ is positive. The second derivative of $`A`$ now becomes
$$A^{\prime \prime }=\frac{9}{4}\left(\frac{W(\phi )}{\phi }\right)^2\left(14\lambda W(\phi )^2\right)$$
(48)
and is no longer to be non-positive always. In five dimensional gauged supergravity, (38) and (39) arise as conditions for unbroken supersymmetry, and (40) is used to prove a $`c`$theorem. It would be interesting to study the implications of (41)-(43) in the context of these models.
The reason that we are able to obtain the first order Bogomol’nyi equations in the presence of higher derivative terms is that we have chosen the particularly nice Gauss-Bonnet combination. For our purposes we regard the method as simply a method for solving coupled differential equations. Thus, for a given choice of $`W`$ we generate a solution for some $`𝒱(\phi ).`$ Note that if we did not use the Gauss-Bonnet combination we would have third and fourth derivatives of $`A`$ appearing in (34).
We now choose $`W(\phi )=s\phi `$ to be a linear function of $`\phi `$. Note that $`s`$ has dimension of an inverse length. We find it convenient to define the length scale $`l=2/(9\sqrt{\lambda }s^2)`$ and the dimensionless ratio $`\sigma =(l/\sqrt{\lambda })^{\frac{1}{2}}.`$
Following the steps outlined above, we generate the solution for $`y>0`$
$$\phi (y)=\frac{3}{2\sqrt{2}}\sigma \mathrm{tanh}\left(\frac{yy_+}{l}\right),$$
(49)
$$A(y)=\frac{1}{2}\sigma ^2\mathrm{log}\left[\mathrm{cosh}\left(\frac{yy_+}{l}\right)\right]+k_+$$
(50)
and for $`y<0`$
$$\phi (y)=\epsilon \frac{3}{2\sqrt{2}}\sigma \mathrm{tanh}\left(\frac{yy_{}}{l}\right),$$
(51)
$$A(y)=\frac{1}{2}\sigma ^2\mathrm{log}\left[\mathrm{cosh}\left(\frac{yy_{}}{l}\right)\right]+k_{}$$
(52)
The potential is a familiar double well
$$𝒱(\phi )=\frac{3}{2\lambda }\left(\frac{1}{\sigma ^2}+1\right)\left(1\frac{8}{9}\frac{\phi ^2}{\sigma ^2}\right)^2\frac{3}{2\lambda },$$
(53)
Note from (36) that we have freedom in choosing the sign of $`\phi `$ represented by $`\epsilon =\pm 1`$ in (51). The $`\phi `$ solution we have is a kink in the fifth dimension interpolating two vacua of the potential. The spacetime is asymptotically $`AdS`$ and there is no singularity at all. Moreover, the 4-D Planck mass is finite.
Unfortunately, in this particular example, the self tuning feature of KSS is also lost. To see this, take $`f(\phi )=V`$ for simplicity. The continuity of $`\phi `$ and $`\phi ^{}`$ fixes $`\epsilon =+1`$ and $`y_+=y_{}y_{}.`$ The continuity of $`A`$ fixes $`k_+=`$ $`k_{}`$ while (24) tells us that the jump in $`[1\frac{4}{3}\lambda (A^{})^2]A^{}`$ across $`y=0`$ is equal to $`\frac{1}{6}V.`$ But this cannot be if $`V`$ is not zero since $`A(y)=\frac{1}{2}\sigma ^2\mathrm{log}\mathrm{cosh}\frac{1}{l}(yy_{})`$ is perfectly smooth across $`y=0.`$ The crucial point here is that we no longer have the freedom of including an additive constant in the solution for $`\phi `$ in this example.
We can also choose a more general superpotential of the form $`W(\phi )=s\phi +r`$. We find that, by choosing $`\lambda <0`$ and $`s=\sqrt{2/(9|\lambda |)}`$, the dilaton potential is a constant
$$𝒱(\phi )=\frac{3}{2|\lambda |}$$
(54)
which acts like a bulk cosmological constant and is independent of $`r`$ in the superpotential $`W(\phi )`$. Thus $`r`$ plays the role of an integration constant and we recover the self tuning feature in KSS. Unfortunately, with negative $`\lambda `$ the hyperbolic functions in (49)-(52) turn into trigonometric functions and we again have naked singularities in the bulk.
In these two examples we constructed, we need either fine tuning to avoid the naked singularities or naked singularities to maintain the self tuning feature.
## Note Added
After this paper was submitted a paper appeared in which the authors proved a no-go theorem which states that, in the scenario of , one needs either fine-tuning or naked singularities to achieve the flatness of our universe. Although the two examples we constructed here, in the presence of Gauss-Bonnet term, are consistent with this no-go theorem, we would like to point out that a crucial ingredient of the proof given in , $`A^{\prime \prime }0`$, is not true in our case, as can be seen from (48). Therefore there might still be hope of retaining self-tuning feature without invoking naked singularities.
###### Acknowledgements.
We thank Lisa Randall and Eva Silverstein for informative discussions and for encouragement, and David Gross, Steve Gubser, Gary Horowitz, and Joseph Polchinski for helpful comments. I. L. also benefitted from discussions with Aki Hashimoto and Michael Quist. This work was supported in part by the NSF under grant number PHY 89-04035 at ITP and by Department of Energy under grant number DOE-ER-40682-143 at CMU. I. L. is supported in part by an ITP Graduate Fellowship.
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# Magnetotransport study of the charged stripes in high-𝑇_𝑐 cuprates
## 1 INTRODUCTION
The conducting state of the high-$`T_c`$ cuprates appears as a result of hole or electron doping of the parent antiferromagnetic (AF) insulator. In general, there is a tendency of doped holes in the AF environment to phase-segregate, which may give rise to an intriguing microscopic state with carriers gathered within an array of quasi-1D “stripes” separating AF domains . An ordered striped structure has been observed in La<sub>2</sub>NiO<sub>4.125</sub> (Ref. ) and in La<sub>1.6-x</sub>Nd<sub>0.4</sub>Sr<sub>x</sub>CuO<sub>4</sub> (Nd-LSCO, Ref. ), while most superconducting cuprates demonstrate incommensurate magnetic fluctuations which can be considered as dynamical stripe correlations . The dynamical stripes might be responsible for the peculiar normal state of cuprates as well as for the occurrence of superconductivity , but still very little is known about the electron dynamics in the stripes.
In this paper, we report an extraordinary behavior of the magnetoresistance (MR) in antiferromagnetic YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> (YBCO), which provides evidences that conducting stripes actually exist in CuO<sub>2</sub> planes and have a considerable impact on the electron transport. We note that a recent study of the Hall effect in Nd-LSCO also provides insights into the electron transport in static stripes .
## 2 EXPERIMENTAL
The high-quality YBCO single crystals were grown by the flux method in Y<sub>2</sub>O<sub>3</sub> crucibles, and a high-temperature annealing was used to reduce their oxygen content. The MR was measured by sweeping the magnetic field at fixed temperatures stabilized by a capacitance sensor with an accuracy of $``$1 mK. The angular dependence of the MR was determined by rotating the sample within a 100 range under constant magnetic fields up to 16 T.
## 3 RESULTS
The heavily-underdoped YBCO crystals, though located deep in the AF range of the phase diagram ($`x`$$``$0.3), are far from conventional insulators: the in-plane resistivity $`\rho _{ab}`$ remains “metallic” at high $`T`$ and it grows slower than expected for the hopping electron transport at low $`T`$ . These AF crystals demonstrate an unusual behavior of the in-plane MR, $`\mathrm{\Delta }\rho _{ab}/\rho _{ab}`$, when the magnetic field $`H`$ is applied along the CuO<sub>2</sub> planes, as shown in Fig. 1. At weak fields, the longitudinal in-plane MR \[$`H`$$``$$`I`$$``$$`ab`$\] is negative and follows roughly a $`T`$-independent $`\zeta H^2`$ curve, and then abruptly saturates above some threshold field. The threshold field and the saturated MR value gradually increase with decreasing temperature. The MR anomaly becomes noticeable near the Néel temperature $`T_N`$$``$230 K ($`T_N`$ can be obtained from the $`\rho _c(T)`$ data, as reported in Ref. ), but evolves rather smoothly through $`T_N`$, which indicates that the long-range AF order itself is not responsible for its origin .
When the magnetic field is turned in the plane to become perpendicular to the current \[$`H`$$``$$`ab`$; $`H`$$``$$`I`$\], the low-field MR term just switches its sign, retaining its magnitude and the threshold-field value . This can be graphically shown in the MR data taken upon rotating $`H`$ within the $`ab`$ plane, which revealed a striking anisotropy with a “$`d`$-wave”-like symmetry; i.e. $`\mathrm{\Delta }\rho _{ab}/\rho _{ab}`$ changes from negative at $`\alpha `$=0 to positive at $`\alpha `$=90, being zero at about 45 ($`\alpha `$ is the angle between $`H`$ and $`I`$), see Fig. 2. It is worth noting that the low-field MR feature is not observed at all when the magnetic field is applied along the $`c`$-axis.
The most intriguing peculiarity of the low-field MR appears at temperatures below $``$25 K, where the $`H`$-dependence of $`\rho _{ab}`$ becomes irreversible. Figure 3 shows the low-field MR term measured for $`H`$$``$$`I`$ (for clarity, the background MR, $`\gamma H^2`$, determined at high fields, is subtracted: $`\mathrm{\Delta }\rho _{ab}/\rho _{ab}`$ = $`(\mathrm{\Delta }\rho _{ab}/\rho _{ab})^{}+\gamma H^2)`$. Initially the irreversibility appears as a small hysteresis on the MR curve; however, upon cooling to 10 K it becomes much more pronounced (the MR peaks are shifted from $`H`$=0 and strongly suppressed). We note that the first field sweep, which starts at $`\mathrm{\Delta }\rho _{ab}/\rho _{ab}`$=0, differs significantly from the subsequent ones. The salient point here is that the resistivity does not return to its initial value after removing the magnetic field; hence, the system acquires a memory. In other words, the application of the magnetic field at low $`T`$ induces a persistent resistivity anisotropy in the CuO<sub>2</sub> planes.
The picture would be incomplete without the data on the $`c`$-axis transport (across the CuO<sub>2</sub> planes). It was shown that in antiferromagnetic YBCO below $`T_N`$, the suppression of spin fluctuations by the magnetic field results in a large positive out-of-plane MR . Figure 4 shows an intriguing MR behavior produced by a superposition of the negative low-field MR feature on top of the large positive $`\gamma H^2`$ background in a sample with $`T_N300`$ K.
## 4 DISCUSSIONS
It is very difficult to understand the MR anomalies presented here, especially the memory effect, without considering an inhomogeneous state or a superstructure in the CuO<sub>2</sub> planes instead of a uniform AF state. The picture of charged “stripes” in the CuO<sub>2</sub> planes allows one to account for all the observed MR peculiarities, by assuming that the magnetic field gives rise to a directional ordering of the stripes . Actually, the aligning of stripes with confined carriers moving along would change the current paths and introduce the in-plane anisotropy. The rotation of stripes by the magnetic field gives a reasonable explanation for the in-plane MR with the $`d`$-wave-shaped angular dependence. Within this picture, the threshold field of several Tesla is presumably coming from the establishment of the directional order of the stripes.
As the temperature is lowered, it is expected that the stripe dynamics slows down and the magnetic domain structure in the CuO<sub>2</sub> planes is frozen, forming a cluster spin glass. The spin-glass transition temperature has been reported to be about 20-25 K for the AF compositions , which is in good agreement with the temperature where the hysteretic MR behavior is found.
Though an explanation of the out-of-plane MR feature is not so straightforward, one can imagine that by adjusting the direction of the stripes in neighboring CuO<sub>2</sub> planes, the magnetic field increases the overlapping both between the stripes in the real space and between their quasi-1D carriers in the $`k`$-space, thereby enhancing the probability of the electron hopping.
One may wonder how the MR anomaly evolves when the carrier density is increased, bringing the system to more metallic region. Figure 5 shows the comparison of the MR data for two different carrier densities; the data were taken on the same sample, where the $``$10% increase in the carrier density was achieved by keeping the sample at room temperature, which causes the oxygen reordering. It is clear that the threshold field for the stripe ordering increases with increasing carrier density. The data in Fig. 5 suggest that the threshold field becomes inaccessibly high when the carrier density is increased to the superconducting region ($`x>0.4`$). This is probably the reason why the MR anomaly reported here has never been observed in superconducting samples.
## 5 CONCLUSION
A variety of unusual MR features found in heavily underdoped YBCO have provided new information on the conducting charged stripes in the CuO<sub>2</sub> planes. The MR behavior implies that the stripes couple to the external magnetic field and undergo topological ordering at fields of the order of a few T, although the actual mechanism that couples the stripes to the magnetic field is not clear yet. Upon cooling the sample below $``$20 K, the dynamics of stripes slows down and the directional order of the stripes becomes persistent, giving rise to a “memory effect” in the resistivity. These findings show that the magnetic field can be used as a tool to manipulate the striped structure and open a possibility to clarify the electron dynamics within the stripes.
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# IS OUR VACUUM STABLE ? 11footnote 1Based on a lecture given by E. I. Guendelman in the symposium ”The Future of the Universe and the Future of our Civilization” July 2-6, 1999
## 1 On Vacuum Stability in the Weinberg-Salam model.
In modern theories of elementary particle interactions, the vacuum is defined by the expectation value of a scalar field. In the Weinberg-Salam model such scalar field (the Higgs field) is responsible for the gauge symmetry breaking and fermion masses production. Let us begin our studies of vacuum instability by considering a model containing just a scalar field with a potential $`V(\varphi )`$ as depicted in fig 1.
The theory is governed by an action in the form $`S=d^4x(\frac{1}{2}_\mu \varphi ^\mu \varphi V(\varphi ))`$, where $`_\mu \varphi ^\mu \varphi =(_0)^2(\varphi )^2`$
As it is obvious, there is no classical (in Minkowski space) solution connecting $`\varphi _+`$ and $`\varphi _{}`$. There is however a Euclidean solution connecting a point close to $`\varphi _{}`$ to $`\varphi _+`$, the tunneling solution. This is a solution in imaginary time. Defining $`\tau =it`$, $`\rho =\sqrt{\stackrel{}{x}^2+\rho ^2}`$ and considering solutions of the form $`\varphi =\varphi (\rho )`$, we obtain the following equation of motion
$$\frac{d^2\varphi }{d\rho ^2}+\frac{3}{\rho }\frac{d\varphi }{d\rho }+V^{^{}}(\varphi )=0$$
(1)
Equation (1) is like that of a particle, $`\varphi `$ playing the role of ”position”, $`\rho `$ ”time”, $`\frac{3}{\rho }\frac{d\varphi }{d\rho }`$ friction, $`V`$, the mechanical potential (see$`V`$ in Fig.2).
To avoid a singularity at $`\rho =0`$, $`\frac{d\varphi }{d\rho }`$ must vanish for $`\rho =0`$. If we set as initial condition at $`\rho =0`$, $`\varphi `$ close to $`\varphi =\varphi _{}`$ the field stays close to $`\varphi =\varphi _{}`$ for a long $`\rho `$ interval, the friction term becomes negligible and $`\varphi `$ arrives at $`\varphi =\varphi _+`$ with $`\frac{d\varphi }{d\rho }0`$, i.e. ”overshoots”. If $`\varphi =\varphi _0`$, because of the friction the field does not make it to $`\varphi =\varphi _+`$. By continuity of a value $`\varphi =\varphi _A`$ exists such that $`\varphi \varphi _+`$ as $`\rho \mathrm{}`$. This is the tunneling solution. This solution has a classical Euclidean action $`S_E`$ and the tunneling probability per unit time per unit volume is proportional to $`e^{S_E}`$. $`\varphi _A\varphi _{}`$ takes place classically.
In the quantum theory one finds zero-point fluctuations that give rise to vacuum energy $`\frac{1}{2}\mathrm{}\omega `$ for each boson mode $`(\omega =\sqrt{\stackrel{}{p}^2+m^2})`$, since masses depend on the Higgs field ($`m^2=V^\mathrm{"}(\varphi )`$ for the Higgs field itself, $`m^2e^2\varphi ^2`$ for gauge vector fields, etc.) and $`\frac{1}{2}\mathrm{\Sigma }\mathrm{}\omega \frac{1}{2}V\frac{d^3k}{(2\pi )^3}\mathrm{}\sqrt{\stackrel{}{k}^2+m^2(\varphi )}`$, we see that an infinite $`\varphi `$ dependent vacuum energy appears. The Dirac fermions also contribute $`\mathrm{}\omega `$ for each occupied negative energy state. The resulting energy density $`V_{eff}(\varphi )`$ makes sense after a renormalization process which get rid of the infinities. The process of renormalization introduces an arbitrary scale into the problem. If we demand that such scale will not affect physical quantities like $`V_{eff}(\varphi )`$ itself, we arrive at what is called renormalization group equation for $`V_{eff}(\varphi )`$. Such $`V_{eff}(\varphi )`$ holds for small coupling although $`\varphi `$ itself can be large.
One can then ask: for which values $`V_{eff}(\varphi )`$ our vacuum is unstable? One can check that for the known values of $`m_{Higgs}>90Gev`$,$`V_{eff}(\varphi =0)>V_{eff}(\varphi _+)`$ ($`\varphi _+`$ is the value of the Higgs for our vacuum), but for large values of $`\varphi `$, $`V_{eff}(\varphi )`$ can be negative (conventionally we set $`V_{eff}(\varphi _+)0`$). The condition that does not happen gives $`m_{top}95Gev+0.60m_{Higgs}`$
Notice that for $`m_{top}175Gev`$, this implies that $`m_{Higgs}135Gev`$, while the experimental data only tell us that $`95Gevm_{Higgs}190Gev`$. So that we cannot tell for sure whether the standard model predicts a stable or unstable vacuum.
Even assuming the vacuum is unstable, which is possible according to what we have seen before, can ask: was the age of the universe long enough so that the probability of some nucleation took place in our past light cone? The space-time volume of our past light cone $`t^4`$, $`t`$ is the age of the universe and the probability of nucleating per unit time per unit volume $`e^{S_E}`$. The probability of nucleation is $`t^4e^{S_E}`$ and the age of the universe in electroweak units is $`te^{101}`$ so that $`t^4e^{S_E}e^{S_E+404}`$. So that $`S_E404`$ gives the middle region Fig. 3 where the rate is not high enough so it is unlikely that we noticed already any phase transition.
The value $`S_E404`$ is consistent with the known bounds on $`m_H`$ and the known value of $`m_{top}`$ (See Fig. 3 which contains the inserted horizontal line $`m_{top}=175Gev,90Gev<m_H<190Gev`$) .
## 2 The Cosmological Constant Problem
Calculations of $`V_{eff}(\varphi )`$ are generally a calculation of the vacuum energy of the universe and how it depends on $`\varphi `$. In flat space there is no meaning to the constant part of $`V_{eff}(\varphi )`$ if gravity is ignored. Once gravity taken into account is very important. The naive prediction for $`V_{eff}(\varphi )|_{\varphi =\varphi _{min}}=\mathrm{}`$ but we really need $`V_{eff}(\varphi )|_{\varphi =\varphi _{min}}=0`$.
One has the suspicion that the calculations of $`V_{eff}(\varphi )`$ mentioned above may miss some physics. Fine tuning $`V_{eff}(\varphi )`$, by adding a constant so that $`V_{eff}(\varphi )=0`$ may not be enough as we will argue, since the physics that sets $`V_{eff}(\varphi _+)0`$ can act again when $`V_{eff}`$ ”tries” to cross once again the zero value and this could affect seriously our notions concerning stability of the vacuum.
## 3 A Model that sets $`V_{eff}(\varphi _+)=0`$ naturally
Observation: Usually Generally Covariant theory is build from the action $`S=L\sqrt{g}d^4x`$, $`g=det(g_{\mu \nu })`$, $`L=scalar`$, ($`L^{}=L`$ under general coordinate transformations). Notice that $`\sqrt{g}d^4x`$ is an invariant volume element. It is possible to build another volume element, independent of $`g_{\mu \nu }`$. Take for example, given 4-scalars $`\phi _a`$ (a = 1,2,3,4), the density
$$\mathrm{\Phi }=\epsilon ^{\mu \nu \alpha \beta }\epsilon _{abcd}_\mu \phi _a_\nu \phi _b_\alpha \phi _c_\beta \phi _d$$
(2)
$`\mathrm{\Phi }`$ transform like $`\sqrt{g}`$, so $`\mathrm{\Phi }d^4x`$ is also an invariant.
One can allow both geometrical objects to enter the theory and consider
$$S=L_1\mathrm{\Phi }d^4x+L_2\sqrt{g}d^4x$$
(3)
Here $`L_1`$ and $`L_2`$ are $`\phi _a`$ independent. There is a good reason not to consider mixing of $`\mathrm{\Phi }`$ and $`\sqrt{g}`$ , like for example using $`\frac{\mathrm{\Phi }^2}{\sqrt{g}}`$. This is because (3) is invariant (up to the integral of a total divergence) under the infinite dimensional symmetry $`\phi _a\phi _a+f_a(L_1)`$ where $`f_a(L_1)`$ is an arbitrary function of $`L_1`$ if $`L_1`$ and $`L_2`$ are $`\phi _a`$ independent. Such symmetry (up to the integral of a total divergence) is absent if mixed terms are present.
We will study now the dynamics of a scalar field $`\varphi `$ interacting with gravity as given by the action (3) with
$$L_1=\frac{1}{\kappa }R(\mathrm{\Gamma },g)+\frac{1}{2}g^{\mu \nu }_\mu \varphi _\nu \varphi V(\varphi ),L_2=U(\varphi )$$
(4)
$$R(\mathrm{\Gamma },g)=g^{\mu \nu }R_{\mu \nu }(\mathrm{\Gamma }),R_{\mu \nu }(\mathrm{\Gamma })=R_{\mu \nu \lambda }^\lambda ,R_{\mu \nu \sigma }^\lambda (\mathrm{\Gamma })=\mathrm{\Gamma }_{\mu \nu ,\sigma }^\lambda \mathrm{\Gamma }_{\mu \sigma ,\nu }^\lambda +\mathrm{\Gamma }_{\alpha \sigma }^\lambda \mathrm{\Gamma }_{\mu \nu }^\alpha \mathrm{\Gamma }_{\alpha \nu }^\lambda \mathrm{\Gamma }_{\mu \sigma }^\alpha .$$
(5)
In the variational principle $`\mathrm{\Gamma }_{\mu \nu }^\lambda ,g_{\mu \nu }`$, the scalar fields $`\phi _a`$ and the scalar field $`\varphi `$ are to be treated as independent variables.
We can require the scale invariance of the theory. If we perform the global scale transformation ($`\theta `$ = constant) $`g_{\mu \nu }e^\theta g_{\mu \nu }`$ then (2), with the definitions (3), (4), is invariant. $`V(\varphi )`$ and $`U(\varphi )`$ are in the form $`V(\varphi )=f_1e^{\alpha \varphi },U(\varphi )=f_2e^{2\alpha \varphi }`$ and $`\phi _a`$ is transformed according to $`\phi _a\lambda _a\phi _a`$ (no sum on a) which means $`\mathrm{\Phi }\left(_a\lambda _a\right)\mathrm{\Phi }\lambda \mathrm{\Phi }`$ such that $`\lambda =e^\theta `$ and $`\varphi \varphi \frac{\theta }{\alpha }`$. In this case we call the scalar field $`\varphi `$ needed to implement the scale invariance as ”dilaton”.
Now,in the general case, let us consider the equations which are obtained from the variation of the $`\phi _a`$ fields. We obtain then $`A_a^\mu _\mu L_1=0`$ where $`A_a^\mu =\epsilon ^{\mu \nu \alpha \beta }\epsilon _{abcd}_\nu \phi _b_\alpha \phi _c_\beta \phi _d`$. Since det $`(A_a^\mu )=\frac{4^4}{4!}\mathrm{\Phi }^30`$ if $`\mathrm{\Phi }0`$. Therefore if $`\mathrm{\Phi }0`$ we obtain that $`_\mu L_1=0`$, or that $`L_1=M`$, where M is constant. This constant M appears in a self-consistency condition of the equations of motion that allows us to solve for $`\chi \frac{\mathrm{\Phi }}{\sqrt{g}}`$
$$\chi =\frac{2U(\varphi )}{M+V(\varphi )}.$$
(6)
To get the physical content of the theory, it is convenient to go to the Einstein conformal frame where
$$\overline{g}_{\mu \nu }=\chi g_{\mu \nu }$$
(7)
and $`\chi `$ given by (6). In terms of $`\overline{g}_{\mu \nu }`$ the non Riemannian contribution (defined as $`\mathrm{\Sigma }_{\mu \nu }^\lambda =\mathrm{\Gamma }_{\mu \nu }^\lambda \{_{\mu \nu }^\lambda \}`$ where $`\{_{\mu \nu }^\lambda \}`$ is the Christoffel symbol), disappears from the equations, which can be written then in the Einstein form ($`R_{\mu \nu }(\overline{g}_{\alpha \beta })`$ = usual Ricci tensor)
$$R_{\mu \nu }(\overline{g}_{\alpha \beta })\frac{1}{2}\overline{g}_{\mu \nu }R(\overline{g}_{\alpha \beta })=\frac{\kappa }{2}T_{\mu \nu }^{eff}(\varphi )$$
(8)
where
$$T_{\mu \nu }^{eff}(\varphi )=\varphi _{,\mu }\varphi _{,\nu }\frac{1}{2}\overline{g}_{\mu \nu }\varphi _{,\alpha }\varphi _{,\beta }\overline{g}^{\alpha \beta }+\overline{g}_{\mu \nu }V_{eff}(\varphi ),V_{eff}(\varphi )=\frac{1}{4U(\varphi )}(V+M)^2.$$
(9)
If $`V(\varphi )=f_1e^{\alpha \varphi }`$ and $`U(\varphi )=f_2e^{2\alpha \varphi }`$ as required by scale invariance, we obtain from (10) $`V_{eff}=\frac{1}{4f_2}(f_1+Me^{\alpha \varphi })^2`$
Since we can always perform the transformation $`\varphi \varphi `$ we can choose by convention $`\alpha >O`$. We then see that as $`\varphi \mathrm{},V_{eff}\frac{f_1^2}{4f_2}=`$ const. providing an infinite flat region. Also a minimum is achieved at zero cosmological constant for the case $`\frac{f_1}{M}<O`$ at the point $`\varphi _{min}=\frac{1}{\alpha }ln\frac{f_1}{M}`$. Finally, the second derivative of the potential $`V_{eff}`$ at the minimum is $`V_{eff}^{\prime \prime }=\frac{\alpha ^2}{2f_2}f_1^2>O`$ if $`f_2>O`$,
There are many interesting issues that one can raise here. The first one is of course the fact that a realistic scalar field potential, with massive excitations when considering the true vacuum state, is achieved in a way which is consistent with the idea of the scale invariance. The second point to be raised is that since there is an infinite region of flat potential for $`\varphi \mathrm{}`$, we expect a slow rolling new inflationary scenario to be viable, provided the universe is started at a sufficiently large value of the scalar field $`\varphi `$. Furthermore, one can consider this model as suitable for the present day universe rather than for the early universe, after we suitably reinterpret the meaning of the scalar field $`\varphi `$. This can provide a long lived almost constant vacuum energy for a long period of time, which can be small if $`f_1^2/4f_2`$ is small.
Such small energy density will eventually disappear when the universe achieves its true vacuum state.
Notice that for generic functions $`V(\varphi )`$, $`U(\varphi )`$ the minimum of $`V_{eff}(\varphi )`$, as given from (9), is at zero if $`V+M=0`$ at some point and if $`V^{}`$ is finite there (also $`U>0`$ there). So $`V_{eff}^{}=0`$ and $`V_{eff}=0`$ is achieved generally without fine tuning! If in the neighborhood of $`\varphi _+`$ $`V+M=0`$, $`U(\varphi )>0`$ and $`V+M`$ as a function of $`\varphi `$ that goes through zero. Then $`V_{eff}\frac{(V+M)^2}{4U(\varphi _+)}`$ has a local minimum at zero. That is $`V(\varphi _+)=0`$ and $`V(\varphi _+)=0`$ automatically without fine tuning. Therefore zero vacuum energy state is obtained naturally!
Going back to the general $`V(\varphi )`$, $`U(\varphi )`$ case we can ask the question: given the classically stable state $`\varphi _+`$, where $`\varphi =\varphi _+`$, $`U(\varphi _+)>0`$, $`V(\varphi _+)+M=0`$ can we make this into an unstable state?. Remember that $`V_{eff}=\frac{(V+M)^2}{4U(\varphi _+)}`$ and that we obtained $`V+M=0`$ as a stable (classically) state under the conditions that $`U`$ is $`>0`$ at this point and it is a regular function there. We take to be true that $`U(\varphi )`$ is a nice function everywhere. The only way $`V_{eff}=\frac{(V+M)^2}{4U}`$ can change sign is for $`U(\varphi )`$ to change sign. For $`U(\varphi )`$ being a nice function, this can only happen if $`U(\varphi )`$ goes to zero. If no fine tuning is invoked at that point, $`V+M0`$ (if it is true, for other $`M`$, i.e. for another initial condition of the universe it will not be true). Then $`V_{eff}`$ looks like in Fig.4
Remember that in the euclidean solution relevant for nucleation -$`V_{eff}`$ is the relevant potential.
It is clear from Fig.5 that no tunneling is possible now!
From eq.(6) we have $`V+M0|\mathrm{\Phi }|>>\sqrt{g}`$ and $`U0|\mathrm{\Phi }|<<\sqrt{g}`$.
It is interesting to see what the volume element $`\sqrt{\overline{g}}d^4x`$ where $`\overline{g}_{\mu \nu }=\chi g_{\mu \nu }=`$ Einstein metric compares with the volume element $`\mathrm{\Phi }d^4x`$. Indeed, $`\sqrt{\overline{g}}=\chi ^2\sqrt{g}=\chi \mathrm{\Phi }`$, so that $`\sqrt{\overline{g}}>>|\mathrm{\Phi }|`$ for $`V+M0`$, i.e. Einstein energy density gets diluted $`V_{eff}0`$. On other limit, $`U0`$, $`\sqrt{\overline{g}}<<|\mathrm{\Phi }|`$, implies that Einstein energy density gets concentrated.This is the physical reason that leads to the energy barrier against crossing from $`V_{eff}>0`$ to $`V_{eff}<0`$.
## 4 The Decay of an Elementary Particle into a Universe
Finally we consider a different type of instability of our vacuum. In this case a region which contains a false vacuum inside, like an elementary particle, can decay into a universe!
We considered a model of an elementary particle as a $`2+1`$ dimensional brane evolving in a $`3+1`$ dimensional space. The introduction of a gauge field which takes place in the brane as well as a normal surface tension, following the standard approach to the theory of extended objects , can lead to a stable ”elementary particle” configuration. The simplest form of the action that permits a stable configuration is:
$$S=\sigma _0\sqrt{h}d^3y+\lambda \sqrt{h}F_{\alpha \beta }F^{\alpha \beta }d^3y$$
(10)
Where $`h_{\alpha \beta },\alpha \beta =0,1,2`$ is the induced metric on the surface of the membrane, $`h=det(h_{\alpha \beta })`$, and $`F_{\alpha \beta }F^{\alpha \beta }`$ the Lagrangian of the gauge field. If we assume a spherically symmetric vector potential in the brane (up to a gauge transformation) of the simplest form (a monopole potential), we receive the general form of the surface tension as; $`\sigma =\sigma _0+\frac{\sigma _1}{r^4}`$ ($`\sigma _1`$ being $`2\lambda f^2`$, f being the strength of the monopole configuration defined by the vector potential in the brane). The energy of the static wall, $`4\pi r^2\sigma `$, has a non trivial minimum for any $`\lambda >0`$ that permits a stable configuration. This is the simplest possible model, below we shall consider the effect of gravity and of an internal vacuum energy. A membrane as discussed above, defines boundaries between different phases with different values for their energy densities.
We took the metric for the inside of the membrane phase, a false vacuum one, as a de Sitter metric, i.e., $`ds^2=(1\chi ^2r^2)dt^2+\frac{dr^2}{(1\chi ^2r^2)}+r^2(d\theta ^2+sin^2\theta d\varphi ^2)`$ where $`\chi ^2=\frac{8\pi \rho G}{3}`$, $`\rho `$ is the energy density. Outside, in the empty space, we can have only a Schwarzschild space time according to Birkhoff’s theorem, i.e., $`ds^2=(1\frac{2GM}{r})dt^2+\frac{dr^2}{(1\frac{2GM}{r})}+r^2(d\theta ^2+sin^2\theta d\varphi ^2)`$ and in the membrane, we have a singular energy momentum tensor.
Demanding that Einstein’s equations to be satisfied not only inside and outside but also in the membrane we get $`\sqrt{1\chi ^2r^2+\dot{r}^2}\sqrt{1\frac{2GM}{r}+\dot{r}^2}=4\pi G\sigma r`$ where $`\dot{r}\frac{dr}{d\tau }`$ and $`\sigma `$ is the one discussed before. Following (7), we take as the Hamiltonian the mass of the system, which gives us:
$$H=\frac{\chi ^2r^3}{2G}\frac{(4\pi G)^2\sigma ^2r^3}{2G}+4\pi \sigma r^2(1\chi ^2r^2+\dot{r}^2)^{1/2}$$
(11)
Having obtained the Hamiltonian, all the others classical dynamical variables can be obtained as was done in . The conjugate momentum p will be equal to
$`p=\frac{L}{\dot{r}}`$, the Lagrangian will be equal to $`L=\dot{r}H\frac{d\dot{r}}{\dot{r}^2}`$ This give for the conjugate momentum $`p=\frac{H}{\dot{r}}\frac{d\dot{r}}{\dot{r}}`$ Using H as before we arrive at the value of p, which is equal to $`p=4\pi \sigma r^2arcsinh(\frac{\dot{r}}{\sqrt{1\chi ^2r^2}})`$ inside horizon and $`p=4\pi \sigma r^2arccosh(\frac{\dot{r}}{\sqrt{\chi ^2r^21}})`$ outside. An arbitrary function of r can be added in the definition of p. Classically it corresponds to an additional total derivative of a function of r in the Lagrangian, while in Quantum Mechanics it corresponds to a redefinition of the wavefunction $`\mathrm{\Psi }^{}=e^{if(r)}\mathrm{\Psi }`$ This means that the Hamiltonian can be taken as
$$H=\frac{\chi ^2r^3}{2G}\frac{(4\pi G)^2\sigma ^2r^3}{2G}+4\pi \sigma r^2\sqrt{1\chi ^2r^2}K(\frac{p}{4\pi \sigma r^2})$$
(12)
where $`K=cosh`$ inside the horizon and $`K=sinh`$ outside it. In order to achieve a quantum mechanical approach we shall assume that $`p=i\frac{}{r}`$ and from this $`e^{ia\frac{}{r}}\mathrm{\Psi }(r)=\mathrm{\Psi }(ria)`$. The Schroedinger equation is $`H\mathrm{\Psi }=m\mathrm{\Psi }`$ in which m is the mass parameter of the external Schwarzschild solution. Defining the dimensionless variable (in units where $`\mathrm{}`$ = c = 1) $`x=\frac{4\pi r^3\sigma _0}{3}4\pi \frac{\sigma _1}{r}`$ we receive the following difference equation for $`\mathrm{\Psi }`$, interpreting the order of operators in $`\frac{p}{4\pi \sigma r^2}`$ as $`\frac{1}{4\pi \sigma r^2}p`$ $`f(x)\mathrm{\Psi }(x)+g(x)[\mathrm{\Psi }(x+i)+\mathrm{\Psi }(xi)]=0`$ f and g are real functions of x, inside the horizon, and
$$f(x)\mathrm{\Psi }(x)+g(x)[\mathrm{\Psi }(x+i)\mathrm{\Psi }(xi)]=0$$
(13)
outside. Expanding the equation for $`\mathrm{\Psi }`$ outside the horizon, taking $`x>>1`$ (setting $`r\frac{1}{\chi }`$ and $`\chi G^{1/2}\rho _0^{1/2}`$, we see that $`x>>1`$ is satisfied if the typical energy scales determining $`\sigma _0,\sigma _1and\rho _0`$ are $`<<`$ Planck scale) and keeping the first nonvanishing contribution only, we receive the equation:
$$\frac{f}{2g}\mathrm{\Psi }(x)=i\frac{\mathrm{\Psi }}{x}$$
(14)
It has the form of a Schroedinger equation (x is time-like outside the horizon). The solution is $`\mathrm{\Psi }=Ce^{i{\scriptscriptstyle ({\scriptscriptstyle \frac{f}{2g}})𝑑x}}`$ where C=constant. This means that once a bubble passes the horizon it will expand indefinitely, since $`|\mathrm{\Psi }|^2`$ = constant and therefore the modulus of the amplitude for the bubble being at $`r=\frac{1}{\chi }+ϵ`$ ($`ϵ>0`$ is very small) is the same as the amplitude for the membrane being at r $`\mathrm{}`$ with probability equal 1. Therefore if the wave function of the bubble has a tail long enough, so it can get the horizon $`r=\frac{1}{\chi }`$, we have the possibility of the formation of an infinite size bubble i.e. the formation of a universe. The resulting universe becomes large not by expanding and displacing an exterior region. This cannot do, since the interior has negative pressure and the outside zero pressure. Really the bubble expands forming a wormhole region that disconnects from the outside creating a ”baby” universe in this case.
## References
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# Analysis of Oscillator Neural Networks for Sparsely Coded Phase Patterns
## 1 Introduction
One of the important but unsolved problems in neuroscience is to determine how information is coded in neuronal activities. Most of the traditional neural network models consisting of binary units are constructed on the assumption that information is coded only in the mean firing rate of the neurons. Although these models have provided us with theoretically interesting information, they ignore many dynamical aspects of neuronal activities. In fact, oscillatory activity appears to be ubiquitous in many neuronal systems. For example, some recent biological experiments have revealed that spatially synchronous oscillations of neuronal ensembles are dependent on the global properties of the external stimulus. It has been suggested that such synchronization is computationally significant in information processing . The hippocampus is also one of the areas in which neuronal synchronization is observed and is believed to play an important role in memory processing.
With these new findings, many models concerning the dynamical aspects of neurons and memory processing have been proposed . Among these, models consisting of networks of oscillator components are particularly attractive, owing to their mathematical tractability. Like the Hopfield model, the simplicity of these models allows us to obtain useful analytic results. In fact, it has been reported that under suitable conditions, the oscillator network models model associative memory, and we can theoretically evaluate their maximum storage capacities and basins of attraction .
However, such phase oscillator models are based on the unrealistic situation that all the components in the network are always in the firing state. In fact, real neurons can be in either the firing or non-firing state. In addition, it is known that the level of activity in our brain is very low. (i.e., at any given time only a small percentage of neurons are in the firing state.) This situation is termed sparse coding. From the results of theoretical studies of associative memory with binary units, it has been found that the storage capacity diverges as $`1/a\mathrm{ln}a`$ as the firing rate $`a`$ becomes small .
It thus seems that to faithfully to capture the essential dynamics of real oscillatory neuronal systems, it is necessary to extend the oscillator model to treat the non-firing state as well as the firing state. In this paper, we present a simple extended oscillator neural network model of this nature. Using this extended model, we study its associative memory ability to recall sparsely coded phase patterns in which some neurons are in the non-firing state and others encode information in the timings of spikes.
In the next section, we first review the theoretical basis of the phase oscillator model and propose an extend version of the oscillator model to treat non-firing states. In the analysis of this model, we estimate the maximum storage capacity, derive the basin of attraction, and evaluate the quality of recalled memories. Then, we find that when we define the threshold as a dynamical variable in a certain manner, the size of the basin of attraction can be increased. Embedding patterns with different activity levels and the influence of synaptic dilution are also studied.
## 2 Theories of Oscillator Neural Networks
### 2.1 Phase Oscillator Model
Let us first start with a survey of the theoretical basis of the phase oscillator model. We assume that neurons fire periodically and interact weakly with each other. In general, although such a neural system can be described in terms of many internal dynamical variables, it is well known that it can be reduced to a system of simple coupled phase oscillators . In this form, we can characterize the state of the i-th neuron by a single variable, $`\varphi _i(t)`$, which is referred to as the phase. This variable represents the timing of the neuronal spikes at time $`t`$. A typical reduced equation takes the form
$$\frac{d\varphi _i(t)}{dt}=\omega _i(t)+\underset{j=1}{\overset{N}{}}J_{ij}\mathrm{sin}(\varphi _j(t)\varphi _i(t)+\beta _{ij}),$$
(1)
where $`J_{ij}`$ and $`\beta _{ij}`$ characterize the interaction between the i-th and j-th neurons. Assuming that all natural frequencies $`\omega _i(t)`$ are equal to some fixed value $`\omega _0`$, we can eliminate the $`\omega _i(t)`$ term in Eq.(1) by redefining $`\varphi _i(t)`$ through $`\varphi _i(t)\varphi _i(t)+\omega _0t`$. When we represent the state of the i-th neuron by the complex form $`W_i(t)=\mathrm{exp}(i\varphi _i(t))`$, Eq.(1) can be written in the alternative form
$$\frac{dW_i(t)}{dt}=\frac{1}{2}(h_i(t)\stackrel{~}{h_i}(t)W_i^2(t)),h_i(t)=\underset{j=1}{\overset{N}{}}C_{ij}W_j(t),$$
(2)
where the complex variable $`C_{ij}=J_{ij}\mathrm{exp}(i\beta _{ij})`$ represents the effect of the interaction between the i-th and j-th neurons and $`\stackrel{~}{h_i}(t)`$ is the complex conjugate of $`h_i(t)`$. Considering the fact that all neurons relax toward the equilibrium state satisfying the relation $`W_i=h_i/|h_i|`$, we can simplify the above to the following discrete time system:
$$W_i(t+1)=\frac{h_i(t)}{|h_i(t)|},h_i(t)=\underset{j=1}{\overset{N}{}}C_{ij}W_j(t).$$
(3)
This is known as a phase oscillator network model. It can be thought of as a synchronous updated version of the phase oscillator neural network.
### 2.2 Extended Model
The weakness of the model described by Eq.(3) is that it can only be used to treat the firing state. We wish to extend this model so that it has the ability of retrieving sparsely coded phase patterns. In Eq.(3), $`h_i(t)`$ can be regarded as the local field produced by all other neurons. This field determines the state of the i-th neuron at the next time step. In a real neuron, when the membrane voltage is less than some threshold, the neuron generates no neuronal spikes. Therefore, it is reasonable to extend the above model so that the generation of spikes depends on the strength of the local field. Based on this consideration, for oscillator neural networks we propose the following generalized model:
$$W_i(t+1)=f(|h_i(t)|)\frac{h_i(t)}{|h_i(t)|},h_i(t)=\underset{ji}{\overset{N}{}}C_{ij}W_j(t).$$
(4)
In this paper, we assume that $`f(x)=\mathrm{\Theta }(xH)`$, where $`\mathrm{\Theta }(x)`$ is a step function defined by
$$\mathrm{\Theta }(x)=\{\begin{array}{cc}1\hfill & \text{for }x0\hfill \\ 0\hfill & \text{for }x<0\text{,}\hfill \end{array}$$
(5)
and $`H`$ is a threshold parameter controlling the activity of the network.
Figure1(a) displays the function $`\mathrm{\Theta }(xH)`$, and (b) illustrates the dynamical change at updating. Since the amplitude $`|W_i(t+1)|`$ depends on the threshold $`H`$, it has a strong influence on the activity of the system.
Now, we define a set of complex patterns denoted by $`\xi _i^\mu =A_i^\mu \mathrm{exp}(i\theta _i^\mu )`$ $`(\mu =1,\mathrm{},P`$ ; $`i=1,\mathrm{},N)`$, where $`P`$ is the total number of patterns and $`N`$ is the total number of neurons. The variables $`\theta _i^\mu `$ and $`A_i^\mu `$ represent the phase and the amplitude of the i-th neuron in the $`\mu `$-th pattern, respectively. For theoretical simplicity, we choose $`A_i^\mu `$ independently with the probability distribution:
$$A_i^\mu =\{\begin{array}{cc}1\hfill & \text{for firing state with probability }a\hfill \\ 0\hfill & \text{for silent state with probability }1a\text{.}\hfill \end{array}$$
(6)
For the firing state, $`\theta _i^\mu `$ is chosen at random from a uniform distribution between $`0`$ and $`2\pi `$. Note that all the patterns have the same mean firing rate, $`a`$. As the learning rule, we adopt a generalized Hebb rule taking the form
$$C_{ij}=\frac{1}{aN}\underset{\nu =1}{\overset{P}{}}\xi _i^\nu \stackrel{~}{\xi _j^\nu }.$$
(7)
This relation is based on information obtained from experiments on biological systems that when the i-th and j-th neurons are firing simultaneously, the connection between them is enhanced, and otherwise no modification occurs.
### 2.3 Theoretical Analysis
To analyze the recalling process theoretically, we must introduce several macroscopic order parameters. The load parameter $`\alpha `$ is defined by $`\alpha =P/N`$. The overlap $`M^\mu (t)`$ between $`W_i(t)`$ and $`\xi _i^\mu `$ at time $`t`$ is defined by
$$M^\mu (t)=m^\mu (t)\mathrm{exp}(i\phi ^\mu (t))=\frac{1}{aN}\underset{j=1}{\overset{N}{}}\stackrel{~}{\xi }_j^\mu W_j(t).$$
(8)
In practice, owing to rotational symmetry, the similarity between the state of the system $`W_i(t)`$ and the $`\mu `$-th pattern $`\xi _i^\mu `$ can be measured by $`m^\mu (t)`$, the amplitude of $`M^\mu (t)`$. Now, let us consider the situation in which the system is retrieving the pattern $`\xi _i^1`$, that is,
$$m(t)m^1(t)O(1),m^\mu (t)O\left(\frac{1}{\sqrt{N}}\right)(\mu 1).$$
(9)
The local field $`h_i(t)`$ can be separated as
$$h_i(t)=M^1(t)\xi _i^1+z_i(t)=|\xi _i^1|m(t)e^{i(\phi ^1(t)+\theta _i^1)}+z_i(t),$$
(10)
where $`z_i(t)`$ is defined by
$$z_i(t)=\frac{1}{aN}\underset{j,\nu 1}{}\xi _i^\nu \stackrel{~}{\xi _j^\nu }W_j(t).$$
(11)
The first term in Eq.(10) is the signal driving the system to recall the pattern, and the second term can be regarded as the noise arising from the other learned patterns. It is the essence of our analysis that we treat the second noise term as a complex Gaussian noise characterized by
$$<z_i(t)>_i=0,<|z_i(t)|^2>_i=2\sigma ^2(t).$$
(12)
Under the above assumptions, we study the properties of this network by applying the methods of statistical neurodynamics . To begin with, we calculate the overlap at time $`t+1`$. From Eqs.(8), (4) and (10), $`M^1(t+1)`$ is given by
$`M^1(t+1)`$ $`=`$ $`m(t+1)\mathrm{exp}(i\phi ^1(t+1))={\displaystyle \frac{1}{aN}}{\displaystyle \underset{j=1}{\overset{N}{}}}\stackrel{~}{\xi }_j^1f(|h_j(t)|){\displaystyle \frac{h_j(t)}{|h_j(t)|}}`$
$`=`$ $`{\displaystyle \frac{1}{aN}}{\displaystyle \underset{j}{}}\stackrel{~}{\xi _j^1}f(|\xi _j^1M^1(t)+z_j(t)|){\displaystyle \frac{\xi _j^1M^1(t)+z_j(t)}{|\xi _j^1M^1(t)+z_j(t)|}}.`$
Now, we assume that the phase of $`M^1(t+1)`$ is almost constant, that is, $`\phi ^1(t+1)\phi _0`$. The validity of this assumption is supported by the results of preliminary numerical simulations. Owing to the rotational symmetry of the complex Gaussian noise, we can replace $`z_j(t)`$ with $`z_j(t)\mathrm{exp}[i(\phi _0+\theta _j^1)]`$. After some calculations, in the limit $`N\mathrm{}`$, we obtain
$`m(t+1)`$ $`=`$ $`{\displaystyle \frac{1}{aN}}{\displaystyle \underset{j}{}}|\xi _j^1|f(||\xi _j^1|m(t)+z_j(t)|){\displaystyle \frac{|\xi _j^1|m(t)+z_j(t)}{||\xi _j^1|m(t)+z_j(t)|}}`$ (13)
$`=`$ $`f(|m(t)+z(t)|){\displaystyle \frac{m(t)+z(t)}{|m(t)+z(t)|}}_{z(t)},`$
where $`_{z(t)}`$ represents the average over the complex Gaussian noise $`z(t)`$ defined by Eq.(12). To calculate Eq.(13) numerically, we need the value of $`\sigma (t)`$, which is the variance of $`z(t)`$. Thus, in the next step, we consider the relation between $`z_i(t+1)`$ and $`z_i(t)`$, from which we can obtain the relation between $`\sigma (t+1)`$ and $`\sigma (t)`$. From Eq.(11), the noise at time $`t+1`$ can be written as
$`z_i(t+1)`$ $`=`$ $`{\displaystyle \frac{1}{aN}}{\displaystyle \underset{j,\nu 1}{}}\xi _i^\nu \stackrel{~}{\xi _j^\nu }W_j(t+1)`$ (14)
$`=`$ $`{\displaystyle \frac{1}{aN}}{\displaystyle \underset{j,\nu 1}{}}\xi _i^\nu \stackrel{~}{\xi _j^\nu }f(|h_j(t)|){\displaystyle \frac{h_j(t)}{|h_j(t)|}}`$
$`=`$ $`{\displaystyle \frac{1}{aN}}{\displaystyle \underset{j,\nu 1}{}}\xi _i^\nu \stackrel{~}{\xi _j^\nu }f(|h_j^\nu (t)+h_j^\nu (t)|){\displaystyle \frac{h_j^\nu (t)+h_j^\nu (t)}{|h_j^\nu (t)+h_j^\nu (t)|}},`$
where $`h_j^\nu (t)`$ and $`h_j^\nu (t)`$ are defined by
$$h_j^\nu (t)=\frac{1}{aN}\underset{k,\mu \nu }{}\xi _j^\mu \stackrel{~}{\xi _k^\mu }W_k(t),h_j^\nu (t)=\frac{1}{aN}\underset{k}{}\xi _j^\nu \stackrel{~}{\xi _k^\nu }W_k(t).$$
(15)
Note that these functions satisfy the relation $`h_j(t)=h_j^\nu (t)+h_j^\nu (t)`$. We can carry out the summation in Eq.(14) under the assumption that $`h_j^\nu (t)`$ is independent of $`\xi _i^\nu `$. Doing so, we obtain
$`z_i(t+1)`$ $``$ $`K(m(t),\sigma (t))+z_i(t)G(m(t),\sigma (t))`$ (16)
$`K(m(t),\sigma (t))`$ $`=`$ $`{\displaystyle \frac{1}{aN}}{\displaystyle \underset{j,\nu 1}{}}\xi _i^\nu \stackrel{~}{\xi _j^\nu }f(|h_i(t)|){\displaystyle \frac{h_i(t)}{|h_i(t)|}}`$ (17)
$`G(m(t),\sigma (t))`$ $`=`$ $`a({\displaystyle \frac{f^{}(|m(t)+z(t)|)}{2}}+{\displaystyle \frac{f(|m(t)+z(t)|)}{2|m(t)+z(t)|}})`$
$`+(1a)({\displaystyle \frac{f^{}(|z(t)|)}{2}}+{\displaystyle \frac{f(|z(t)|)}{2|z(t)|}})_{z(t)}.`$
Note that Eqs.(16), (17) and (2.3) are needed to calculate $`\sigma `$ in both the equilibrium and the non-equilibrium cases we will describe in the next two sections.
### 2.4 Equilibrium State
In this section, we consider the equilibrium state of our model, in which $`z_i(t)`$ is constant. Applying $`z_i(t+1)=z_i(t)=z_i`$ to Eq.(16), we obtain
$$z_i=\frac{K}{1G}.$$
(19)
Using this equation, we immediately obtain
$$2\sigma ^2=<|z_i|^2>_i=\frac{|K|^2}{(1G)^2},$$
(20)
where
$`|K|^2`$ $`=`$ $`\alpha af(|m+z|)+(1a)f(|z|)_z`$ (21)
$``$ $`\alpha Q.`$ (22)
Consequently, we find that the equilibrium state satisfies the equations
$`m`$ $`=`$ $`f(|m+z|){\displaystyle \frac{m+z}{|m+z|}}_z`$ (23)
$`\sigma ^2`$ $`=`$ $`{\displaystyle \frac{\alpha }{2(1G)^2}}Q`$ (24)
$`G`$ $`=`$ $`a({\displaystyle \frac{f^{}(|m+z|)}{2}}+{\displaystyle \frac{f(|m+z|)}{2|m+z|}})`$
$`+(1a)({\displaystyle \frac{f^{}(|z|)}{2}}+{\displaystyle \frac{f(|z|)}{2|z|}})_z`$
$`Q`$ $`=`$ $`af(|m+z|)+(1a)f(|z|)_z.`$ (26)
Solving the above equilibrium equations, we can find two types of solutions. If the load parameter $`\alpha `$ is smaller than a certain value $`\alpha _c`$, there exists a solution for which the overlap $`m`$ between the system and the pattern is not zero. As $`m0`$ implies that the system has a correlation with the retrieved pattern, this solution corresponds to the retrieval state. Then, if $`\alpha `$ is larger than $`\alpha _c`$, there is only one solution, and for this solution $`m=0`$. This solution therefore corresponds to the non-retrieval state.
The critical load parameter $`\alpha _c`$ is called “the maximum storage capacity”. We now examine our theoretical results by comparing them with results from numerical simulations. In Fig.2(a), we show the dependence of $`\alpha _c`$ on several parameters, such as the threshold $`H=0.3,0.5,0.8`$ and the activity level $`a`$. In (a), We can see that as the activity level $`a`$ decreases, the storage capacity $`\alpha _c`$ increases for each $`H`$. In particular, in the limit $`a0`$, we numerically find that the storage capacity diverges in proportion to $`1/a\mathrm{ln}a`$, as in the case of the Hopfield model. The maximum storage capacity as a function of $`a`$ and $`H`$ is illustrated in (b). It is shown there that for any $`H`$, the storage capacity diverges as $`a`$ decreases to zero.
### 2.5 Non-equilibrium State
In this section, we derive dynamical equations to study the retrieving process. To obtain a recursion equation for $`\sigma `$, we start with Eq.(16). Squaring Eq.(16), we obtain
$`2\sigma ^2(t+1)`$ $`=`$ $`\alpha Q(m(t),\sigma (t))+2\sigma ^2(t)G^2(m(t),\sigma (t))`$ (27)
$`+2ReK(m(t),\sigma (t))\stackrel{~}{z}_i(t)\stackrel{~}{G}(m(t),\sigma (t))_{z(t)},`$
where $`K(m(t),\sigma (t))`$ and $`G(m(t),\sigma (t))`$ are given by Eqs.(17) and (2.3), and
$$Q(m(t),\sigma (t))=af(|m(t)+z(t)|)+(1a)f(|z(t)|)_{z(t)}.$$
(28)
To proceed with the calculation, we must estimate the time correlation of the noise. For the first-order approximation, we ignore the temporal correlation of $`z_i(t)`$. We then obtain
$`\sigma ^2(t+1)`$ $`=`$ $`{\displaystyle \frac{\alpha }{2}}Q(m(t),\sigma (t))+\sigma ^2(t)G^2(m(t),\sigma (t))`$ (29)
$`+\alpha a^2m(t+1)m(t)G(m(t),\sigma (t)).`$
This corresponds to the Amari-Maginu theory in the case of traditional neural networks.
For the second-order approximation, we take into account the fact that $`z_i(t)`$ is correlated only with $`z_i(t1)`$, while all correlations with $`z_i(t^{})`$ for $`t^{}<t1`$ are ignored. In this case, $`\sigma (t+1)`$ can be obtained from
$`\sigma ^2(t+1)`$ $`=`$ $`{\displaystyle \frac{\alpha }{2}}Q(m(t),\sigma (t))+\sigma ^2(t)G^2(m(t),\sigma (t))`$
$`+\alpha G(m(t),\sigma (t))X(t+1,t)`$
$`+\alpha a^2m(t+1)m(t1)G(m(t),\sigma (t))G(m(t1),\sigma (t1)),`$
where $`Q(m(t),\sigma (t))`$ is given by Eq.(28), $`G(m(t),\sigma (t))`$ is given by Eq.(2.3), and
$$X(t+1,t)=Re<W_j(t+1)\stackrel{~}{W}_j(t)>_j.$$
(31)
Here,
$$\rho (t,t1)=\frac{\alpha X(t,t1)}{2\sigma (t)\sigma (t1)}+\frac{\sigma (t1)}{\sigma (t)}G(m(t1),\sigma (t1))$$
(32)
is needed to evaluate Eq.(31). Note that $`\rho (t,t1)`$ is the correlation coefficient between $`z(t)`$ and $`z(t1)`$. For the second order approximation, consequently, the retrieving process of the network is described by the equations
$`m(t+1)`$ $`=`$ $`{\displaystyle \frac{1}{aN}}{\displaystyle \underset{j}{}}|\xi _j^1|f(||\xi _j^1|m(t)+z_j(t)|){\displaystyle \frac{|\xi _j^1|m(t)+z_j(t)}{||\xi _j^1|m(t)+z_j(t)|}}`$ (33)
$`=`$ $`f(|m(t)+z(t)|){\displaystyle \frac{m(t)+z(t)}{|m(t)+z(t)|}}_{z(t)}`$
$`\sigma ^2(t+1)`$ $`=`$ $`{\displaystyle \frac{\alpha }{2}}Q(m(t),\sigma (t))+\sigma ^2(t)G^2(m(t),\sigma (t))`$
$`+\alpha G(m(t),\sigma (t))X(t+1,t)`$
$`+\alpha a^2m(t+1)m(t1)G(m(t),\sigma (t))G(m(t1),\sigma (t1))`$
$`X(t+1,t)`$ $`=`$ $`Re<W_j(t+1)\stackrel{~}{W}_j(t)>_j`$ (35)
$`\rho (t,t1)`$ $`=`$ $`{\displaystyle \frac{\alpha X(t,t1)}{2\sigma (t)\sigma (t1)}}+{\displaystyle \frac{\sigma (t1)}{\sigma (t)}}G(m(t1),\sigma (t1))`$ (36)
$`G(m(t),\sigma (t))`$ $`=`$ $`a({\displaystyle \frac{f^{}(|m(t)+z(t)|)}{2}}+{\displaystyle \frac{f(|m(t)+z(t)|)}{2|m(t)+z(t)|}})`$
$`+(1a)({\displaystyle \frac{f^{}(|z(t)|)}{2}}+{\displaystyle \frac{f(|z(t)|)}{2|z(t)|}})_{z(t)}`$
$`Q(m(t),\sigma (t))`$ $`=`$ $`af(|m(t)+z(t)|)+(1a)f(|z(t)|)_{z(t)}.`$ (38)
For initial conditions, we choose $`\sigma ^2(0)=a\alpha /2`$, $`X(0,1)=0`$ and $`X(1,0)=a^2m(1)m(0)`$. We choose many values of the initial overlap $`m(0)`$ and carry out numerical calculations for each. In this way, we determine the lower bound, under which the network fails to retrieve the pattern.
For higher-order approximations, we could derive similar generalized equations describing the retrieval dynamical process. However, as shown in Fig.3, the theoretical prediction at second order gives reasonable agreement with the numerical simulation in contrast to that at first order. We thus conclude that it is sufficient to use the second-order approximation for the theoretical analysis in this paper.
Figure4 displays a phase diagram for $`m(0)_c`$ and $`m(\mathrm{})`$ at various mean activity levels $`a`$ and thresholds $`H`$. The solid lines and the data points here correspond to the theoretical and the numerical results, respectively. It is seen that even in the region satisfying $`\alpha <\alpha _c`$, the system cannot retrieve the pattern if the initial overlap $`m(0)`$ is smaller than $`m(0)_c`$. The boundary corresponding to $`m(0)_c`$ is represented by the lower curves in Fig.4. Thus, for $`m(0)m(0)_c`$, $`m(t)`$ reaches the value of the upper curves $`m(\mathrm{})`$, while for $`m(0)<m(0)_c`$, $`m(t)`$ decreases to zero. We should note that the basins of attraction remain wide even near $`\alpha _c`$. This can be thought of as an advantage of associative memory.
To predict the correct behavior of the retrieval dynamics, it is necessary for the time correlation of the noise terms to be taken into account, and thus the second-order approximation discussed above is necessary. In the case we consider, this order is also sufficient, but interestingly, it has been reported that the fourth-order approximation is necessary when $`a=1`$ and $`H=0`$ . The question arises why the second-order approximation is not sufficient in this case, while it is sufficient in the case we consider. The key point here is that the order of the last term in Eq.(2.5) is proportional to the square of the activity level $`a^2`$. For $`a<1`$, this factor $`a^2`$ weakens the influence of the time correlation on the recalling process. Consequently, for a sparse coded pattern with $`a<1`$, even the second-order approximation results in reasonable agreement with the numerical results.
### 2.6 Sequence Generator
In this section, we consider the case in which the network retrieves a cyclic sequence of $`P`$ patterns associatively, say $`\xi ^1\xi ^2\mathrm{}\xi ^P\xi ^1\mathrm{}`$. In order to allow for such a process, we employ synaptic connections of the form
$$C_{ij}=\frac{1}{aN}\underset{\nu =1}{\overset{P}{}}\xi _i^{\nu +1}\stackrel{~}{\xi _j^\nu }.$$
(39)
In a manner similar to that in the derivation of Eqs.(29) and (2.5), we obtain the following equations:
$`\sigma ^2(t+1)`$ $`=`$ $`{\displaystyle \frac{\alpha }{2}}Q(m(t),\sigma (t))+\sigma ^2(t)G^2(m(t),\sigma (t))`$ (40)
$`m(t+1)`$ $`=`$ $`f(|m(t)+z(t)|){\displaystyle \frac{m(t)+z(t)}{|m(t)+z(t)|}}_{z(t)}.`$ (41)
Note that, since the target pattern changes from time to time, for the definition of the overlap we adopt $`m(t)m^\mu (t)=|\frac{1}{aN}_{j=1}^N\stackrel{~}{\xi _j^\mu }W_j(t)|`$, where $`\mu `$ is the number of the target pattern at time $`t`$. Here $`G(m(t),\sigma (t))`$ and $`Q(m(t),\sigma (t))`$ are the same as those defined by Eqs.(2.5) and (38), respectively. We should note that in the limit $`N\mathrm{}`$, the last term in Eqs.(29) and (2.5) vanishes in Eq.(40), because the effect of the time correlation can be ignored in the above derivation. Therefore, in a sequence generator, it is expected that our theoretical prediction is almost exact.
Figure5(a) displays a phase diagram obtained from our theoretical analysis. All the lines of theoretical results agree with the numerical ones quite well, as expected. This agreement suggests that higher-order approximations will result in even better agreement and leads us to believe that our theoretical derivation is valid. In (b), it is also shown that, like auto-associative memory, the storage capacity diverges in the limit $`a0`$. This is regarded as expressing the meaning that in the limit of sparse coding (i.e., as $`a0`$), we can embed an infinite length of sequential patterns.
### 2.7 Dynamically Adjusted Threshold
From Fig.4, it appears that the basin of attraction for our model is smaller than those in the binary model and the phase oscillator. In this section, we attempt to increase the size of the basin of attraction by defining the threshold as a dynamical variable that is proportional to the standard deviation of the noise. Therefore, assuming the same condition as in Sec.2.2, we add the equation
$$H(t)=\sqrt{2\mathrm{ln}a}\sigma (t),$$
(42)
where $`H(t)`$ is the threshold at time $`t`$. This choice for the form of $`H(t)`$ is made to insure the relations $`m1`$ and \[firing rate\] $`a`$, which cause the Hamming distance between the state of the neuron and the retrieved pattern to be small . Comparing Fig.6 with Fig.4, it is seen that the basin of attraction can be enlarged by introducing the dynamically adjusted threshold into our model. Therefore, when we introduce the dynamically adjusted threshold, our model has an advantage similar to that of the binary model and the phase oscillator.
## 3 Patterns Generated with Different Firing Rates
To this point, we have assumed that the activity levels are equal for all patterns. However, it is likely that the actual activity level depends on the pattern which the network is retrieving presently, in other words, on the content of the required information. Unfortunately, using a traditional neural network, we encounter difficulties in storing patterns with different activity levels simultaneously. In this section, we demonstrate that, unlike traditional models, the proposed oscillator model can easily store multiple patterns with different activity levels using a simple Hebbian learning rule.
Let us consider a set of complex patterns defined by
$$Prob[|\xi _i^\mu |=1]=\{\begin{array}{cc}a_1\hfill & \text{for }1\mu P_1\hfill \\ a_2\hfill & \text{for }P_1+1\mu P\text{,}\hfill \end{array}$$
(43)
where generally $`a_1a_2`$. Thus, the total number of the patterns with activity level $`a_1`$ is $`P_1`$, while the total number of patterns with the activity level $`a_2`$ is $`P_2=PP_1`$. Using the above patterns, grouped into two different activity levels, we examine whether the network can retrieve patterns having different firing rates. To retrieve such patterns, we use the following modified form of the learning rule Eq.(7):
$$C_{ij}=\frac{1}{a_1N}\underset{\nu =1}{\overset{P_1}{}}\xi _i^\nu \stackrel{~}{\xi _j^\nu }+\frac{1}{a_2N}\underset{\nu =P_1+1}{\overset{P}{}}\xi _i^\nu \stackrel{~}{\xi _j^\nu }.$$
(44)
Note that this is slightly different from the covariance rule adopted in the context of the learning of sparsely coded patterns. Owing to the rotational symmetry of the phase distribution in the patterns, $`C_{ij}`$ can be defined in terms of the patterns $`\xi _i^\nu `$ themselves, rather than the difference between the patterns and the average activity.
With the method described in Sec.2, we can derive both equilibrium and dynamical equations. The results have the same forms as those in Sec.2. However, here the value $`a_1`$ should be used in place of the activity parameter $`a`$ in the previous equations, because the retrieval pattern $`\xi _i^1`$ has an activity level $`a_1`$. Let us define the load parameter as $`\alpha _1=P_1/N`$ with respect to the retrieval pattern. (The usual load parameter is defined as $`\alpha =P/N=(P_1+P_2)/N`$.) The theoretical analysis yields $`\alpha _{1c}=const.\alpha _2`$, where $`\alpha _{1c}`$ is the maximum storage capacity for $`a_1`$, and $`\alpha _2`$ is the storage capacity for $`a_2`$. Figure7 displays $`\alpha _{1c}`$ as a function of $`\alpha _2=P_2/N`$. We can see from Fig.7 that $`\alpha _{1c}=\alpha (a_1)\alpha _2`$, where the constant value $`\alpha (a_1)`$ is given by the equations in (23), (24), (2.4) and (26) for $`a=a_1`$. For the maximum storage capacity $`\alpha _{2c}`$ associated with the activity level $`a_2`$, we have $`\alpha _{2c}=\alpha (a_2)\alpha _1`$.
Let us consider the case $`P_1=P_2`$. Note that in this case, $`\alpha _1=\alpha _2`$ and the total storage capacity $`\alpha `$ has the relation as $`\alpha =2\alpha _1=2\alpha _2`$. We assume that $`a_1=0.1`$, $`a_2=0.2`$ and $`H=0.3`$. Under these conditions, the basin of attraction obtained numerically from the theory is displayed in the left panel of Fig.8. We should note that in the region (b), the patterns with activity $`a_1`$ can be recalled, while the patterns with activity $`a_2`$ cannot. Since we consider the case $`P_1=P_2`$, the vertical lines correspond to half of the usual maximum storage capacities. In the right figures, we display typical behavior of the overlap $`m(t)`$ for the initial condition $`m(0)=0.5`$. The load parameters $`\alpha =0.02,0.05`$ and $`0.08`$ used here correspond to the regions (a), (b) and (c) shown in the left figure, respectively. In the right panel of Fig.8 appear comparisons of the time evolution of the overlap for two different values of $`a`$ in three different regions of $`\alpha `$. The evolutions corresponding to these values of $`a`$ represent the retrieval processes of $`\xi _i^1`$ and $`\xi _i^{P_1+1}`$ associated with the activity levels $`a_1`$ and $`a_2`$, respectively. Note that in the region (b), the patterns with $`a=0.2`$ act only as noise when the network is recalling the pattern with $`a=0.1`$. As a whole, the above findings suggest that the network has a good ability to retrieve patterns with different firing rates. We should remark that the patterns with activity $`a_1`$ can be stored more stably rather than those with activity $`a_2`$ when $`a_1<a_2`$.
## 4 Dilution
In this section, we study the influence of random synaptic dilution on the model’s associative memory capability. For the case of the phase oscillator model, that is, when $`a=1`$ and $`H=0`$, this effect has already been reported . Following the method used in that case to treat random synaptic dilution in our model, we assume that the synaptic efficacies take the form
$$\overline{C}_{ij}=\frac{c_{ij}}{c}C_{ij},$$
(45)
where $`C_{ij}`$ is the standard Hebbian matrix, as defined by Eq.(7), and, the $`c_{ij}`$ are independent random variables, taking the values $`1`$ and $`0`$ with probabilities $`c`$ and $`1c`$, respectively. Note that the dilution parameter $`c`$ represents the ratio of connected synapses. In the limit $`N\mathrm{}`$, Eq.(45) can be regarded as
$$\overline{C}_{ij}=C_{ij}+\eta _{ij},$$
(46)
where the synaptic noise $`\eta _{ij}`$ is a complex Gaussian noise with mean $`0`$ and variance $`\eta ^2/N`$ . The relationship between the dilution parameter $`c`$ and the variance $`\eta `$ can be calculated as
$$\eta ^2=\frac{1c}{c}\alpha .$$
(47)
Under the same assumption as that used in obtaining Eq.(9), we can separate the local field into the signal and two noise parts as follows:
$`h_i(t)`$ $`=`$ $`{\displaystyle \underset{j}{\overset{N}{}}}\overline{C}_{ij}W_j(t)`$ (48)
$`=`$ $`{\displaystyle \underset{j}{\overset{N}{}}}(C_{ij}+\eta _{ij})W_j(t)`$
$`=`$ $`{\displaystyle \underset{j}{\overset{N}{}}}({\displaystyle \frac{1}{aN}}{\displaystyle \underset{\mu =1}{\overset{P}{}}}\xi _i^\mu \stackrel{~}{\xi _j^\mu }+\eta _{ij})W_j(t)`$
$`=`$ $`M^1(t)\xi _i^1+z_i(t)`$
$`=`$ $`M^1(t)\xi _i^1+z_i^c(t)+z_i^s(t),`$ (49)
where $`z_i^c(t)`$ and $`z_i^s(t)`$ are defined as
$`z_i^c(t)`$ $`=`$ $`{\displaystyle \frac{1}{aN}}{\displaystyle \underset{j}{\overset{N}{}}}{\displaystyle \underset{\mu 1}{\overset{P}{}}}\xi _i^\mu \stackrel{~}{\xi _j^\mu }W_j(t)`$ (50)
$`z_i^s(t)`$ $`=`$ $`{\displaystyle \underset{j}{\overset{N}{}}}\eta _{ij}W_j(t).`$ (51)
The new noise term $`z_i^s(t)`$ is caused by synaptic dilution, while the term $`z_i^c(t)`$ is like the noise defined by Eq.(11), representing the crosstalk noise arising from the other embedded patterns. In analogy to our treatment in Sec.2, we assume that $`z_i^c(t)`$ and $`z_i^s(t)`$ are independent complex Gaussian noises with mean $`0`$ and variance $`\sigma _c^2(t)`$ and $`\sigma _s^2(t)`$, respectively. Therefore, $`z_i(t)`$ can be regarded as a complex Gaussian noise with mean $`0`$ and variance $`2\sigma ^2(t)=\sigma _c^2(t)+\sigma _s^2(t)`$. Thus, we can derive dynamical equations in the same way as in Sec.2. The form of the macroscopic order parameter $`m(t+1)`$ is the same as that given by Eq.(13):
$`m(t+1)`$ $`=`$ $`{\displaystyle \frac{1}{aN}}{\displaystyle \underset{j}{}}|\xi _j^1|f(||\xi _j^1|m(t)+z_j(t)|){\displaystyle \frac{|\xi _j^1|m(t)+z_j(t)}{||\xi _j^1|m(t)+z_j(t)|}}`$ (52)
$`=`$ $`f(|m(t)+z(t)|){\displaystyle \frac{m(t)+z(t)}{|m(t)+z(t)|}}_{z(t)}.`$
First, applying the theory of statistical neurodynamics to $`z_i^s(t)`$, we calculate $`\sigma _s^2(t)`$. However, to apply the theory of statistical neurodynamics, it is necessary to take into account the second term in
$$z_i^s(t+1)=\underset{j}{\overset{N}{}}\eta _{ij}f(|h_j^0|)\frac{h_j^0}{|h_j^0|}+W_i(t)\underset{j}{\overset{N}{}}\eta _{ij}\eta _{ji}\left(\frac{f^{}(|h_j^0|)}{2}+\frac{f(|h_j^0|)}{2|h_j^0|}\right),$$
(53)
where $`h_j^0=h_j(t)\eta _{ji}W_i(t)`$. As reported in , for the asymmetrical case, $`\eta _{ij}\eta _{ji}`$, the second term here can be ignored. Then, as seen from Fig.(9), there is little difference between symmetric and asymmetric cases in a sparse coding system. Thus we conjecture that this term can be ignored altogether. Doing so, in the $`N\mathrm{}`$ limit, we obtain
$`\sigma _s^2(t)`$ $`=`$ $`\eta ^2Q(m(t),\sigma (t))`$ (54)
$`=`$ $`{\displaystyle \frac{1c}{c}}\alpha Q(m(t),\sigma (t)),`$
where $`Q(m(t),\sigma (t))`$ is defined by Eq.(38).
For the crosstalk noise $`z_i^c(t)`$, in analogy to Eq.(16), we obtain
$$z_i^c(t+1)=K(m(t),\sigma (t))+z_i^c(t)G(m(t),\sigma (t)),$$
(55)
where $`K(m(t),\sigma (t))`$ and $`G(m(t),\sigma (t))`$ are defined by Eqs.(17) and (2.3), respectively.
### 4.1 Equilibrium State
In the equilibrium state, putting $`z_i^c(t)=z_i^c(t+1)`$ into Eq.(55), we obtain
$`z_i^c`$ $`=`$ $`{\displaystyle \frac{K}{1G}}`$ (56)
$`\sigma _c^2`$ $`=`$ $`{\displaystyle \frac{\alpha Q}{(1G)^2}}.`$ (57)
Therefore, $`\sigma ^2`$, the variance of $`z(t)`$, takes the form
$`\sigma ^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sigma _c^2+{\displaystyle \frac{1}{2}}\sigma _s^2`$ (58)
$`=`$ $`\left({\displaystyle \frac{1}{2(1G)^2}}+{\displaystyle \frac{1c}{2c}}\right)\alpha Q.`$
Consequently, the properties of the network in the equilibrium state can be calculated from Eqs.(23), (58), (2.4) and (26). In Fig.10, we summarize the theoretical results concerning the dependence on the ratio of connected synapses $`c`$. In Figs.10(a) and (b), it is found that the maximum storage capacity is an increasing function of the connectivity $`c`$. Particularly, we can see from (b) that, as the activity level $`a`$ becomes small, the storage capacity decreases almost linearly with the connectivity $`c`$. A similar linear dependence is observed in the case of the diluted Hopfield model. On the other hand, in (c), it is found that even if the ratio of connectivity is small, the maximum storage capacity tends to diverge in the limit $`a0`$. However, it seems that the rate of this divergence decreases as $`c`$ decreases.
### 4.2 Non-equilibrium State
In order to estimate the robustness with respect to synaptic damage, we should also consider the influence of synaptic dilution on the retrieval process, particularly, on the basin of attraction. For this purpose, we can apply the same method as used in the derivation of Eq.(2.5). After some calculations, the resulting recursion equations at second order are given by
$`\sigma ^2(t+1)`$ $`=`$ $`{\displaystyle \frac{\alpha }{2}}Q(m(t),\sigma (t))+\sigma ^2(t)G^2(m(t),\sigma (t))`$ (59)
$`+\alpha G(m(t),\sigma (t))X(t+1,t)`$
$`+\alpha a^2m(t+1)m(t1)G(m(t),\sigma (t))G(m(t1),\sigma (t1))`$
$`+{\displaystyle \frac{1}{2}}(1G^2(m(t),\sigma (t)))\eta ^2Q(m(t),\sigma (t))`$
$`X(t+1,t)`$ $`=`$ $`Re<W_j(t+1)\stackrel{~}{W}_j(t)>_j`$ (60)
$`\rho (t,t1)`$ $`=`$ $`{\displaystyle \frac{\alpha X(t,t1)}{2\sigma (t)\sigma (t1)}}+{\displaystyle \frac{\sigma (t1)}{\sigma (t)}}G(m(t1),\sigma (t1))`$
$`+{\displaystyle \frac{X(t,t1)G(m(t1),\sigma (t1))Q(m(t1),\sigma (t1))}{2\sigma (t)\sigma (t1)}}\eta ^2,`$
where $`\eta `$ is related to the connectivity $`c`$ via Eq.(47). As mentioned in Sec.2, Eq.(4.2) is used to calculate Eq.(60). For initial conditions, we adopt the values used in Sec.2, except in the present case we use $`\sigma ^2(0)=a\alpha /2c`$. In the case of $`a=0.1`$ and $`H=0.5`$, Fig.11 illustrates the theoretical results concerning the basins of attraction for various values of the connectivity $`c`$. We can see that near saturation $`\alpha \alpha _c`$, the basin of attraction remains large even for low connectivity. Therefore, we find that synaptic dilution has little influence on the width of the basin of attraction, even though the storage capacity decreases with the connectivity.
## 5 Conclusion
In this paper, we have presented a simple extended model of oscillator neural networks to allow for the description of the non-firing state. We have studied the model’s associative memory capability for sparsely coded phase patterns, in which some neurons are in the non-firing state and the other neurons encode information in the phase variable representing the timing of neuronal spikes. In particular, applying the theory of statistical neurodynamics, we have evaluated the maximum storage capacity and derived the basin of attraction. We have found the following properties of our model in its basic form:
* The storage capacity diverges as the activity level decreases to zero. It was numerically found that the storage capacity diverges proportionally with $`1/a\mathrm{ln}a`$ in the limit $`a0`$.
* Even just below the maximum storage capacity, the basin of attraction remains large.
We then investigated the model with regard to the size of the basin of attraction. We found that with the model in its basic form, the basin of attraction is smaller than those of the binary model and the phase oscillator. For associative memory, it is desirable that the basin of attraction be large. For this reason, we considered employing a dynamically adjusted threshold, and we found the following:
* The basin of attraction can be enlarged by using the dynamically adjusted threshold.
In view of biology, the neurons may die of age or be injured by accident. Thus, the robustness with respect to synaptic damage is important for real neuronal systems. For this reason, we also investigated our model with regard to robustness, and we found the following:
* It was found that the system is robust with respect to synaptic damage: Even in the case of a high cutting rate, the basin of attraction remains large, and the maximum of the storage capacity diverges in the $`a0`$ limit. For low activity patterns, the maximum storage capacity decreases almost linearly with the ratio of connected synapse.
The above properties are common with the Hopfield model. In addition, we have found that our model possesses a novel feature not seen in the Hopfield model. In realistic situations, the activity level of the firing pattern may generally depend on the content of the information. Using traditional neural network models based only on firing rate coding, however, we encounter difficulties when the network simultaneously stores such patterns with different activity levels. Contrastingly, with our model, we found the following:
* Unlike the Hopfield model, provided that the phase distribution in the embedded patterns is uniform, it was shown that patterns with different activity levels can be memorized simultaneously.
In conclusion, from the above findings it is seen that the oscillator neural network exhibits good performance in the sparse coding situation. Therefore, we believe that these results support the plausibility of temporal coding and we hope that they encourage attempts for more detailed explorations of the nature of temporal coding. As a future work, we wish to consider one of the applications of the oscillator model. In this paper we considered coherent oscillations of only a uniform distribution. It is more interesting, however, to consider systems that are not so strongly constrained. For example, neurons could be organized into ensembles characterized by different phase values. In such a situation, each ensemble would carry out a different function independently, and the neural system as a whole could carry out several tasks simultaneously. Using this kind of synchrony, we may obtain a simple solution to the binding problem.
## 6 Acknowledgement
We express our gratitude to Prof. T. Munakata and Dr. K. Kitano for helpful discussions. This work was supported by the Japanese Grant-in-Aid for Science Research Fund from the Ministry of Education, Science and Culture.
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# Adelic Formulas for Gamma and Beta Functions of One-Class Quadratic Fields: Applications to 4-Particle Scattering String Amplitudes
## 1 INTRODUCTION
Here, it is relevant to note that the ideas and methods developed by the great Russian mathematician and physicist Nikolai Nikolaevich Bogolyubov (see his selected works ) are widely used in the construction of $`p`$-adic mathematical physics. The present work is devoted to the regularization of divergent adelic products for string and superstring 4-particle amplitudes and in many aspects follows the Bogolyubov ideas. In fact, the theory of strings and superstrings has evolved from the theory of dispersion relations in quantum field theory to whose development a fundamental contribution was made by Bogolyubov.
In the last decade, an interest in $`p`$-adic numbers has considerably increased. These numbers were discovered by the German mathematician K. Hensel about 100 years ago. Until recently, they have not found any significant applications although constituted an essential part of number theory and theory of representations \[4–8\]. Unexpected applications of these numbers have been associated with the discovery of a non-Archimedean structure of space–time at supersmall, the so-called Planck, distances (of order $`10^{33}`$ cm) in quantum gravitation and theory of strings. Therefore, one cannot use real numbers with Archimedean structure as space–time coordinates at such small distances; one needs other, non-Archimedean number fields. An attempt to preserve rational numbers as physical observables in appropriate non-Archimedean fields leads us to the unambiguous conclusion that the fields of $`p`$-adic numbers $`_p,`$ where $`p=2,3,5,\mathrm{}`$ are prime numbers, can be chosen as such number fields (the Ostrovskii theorem!) . Thus, we arrive at the problem of constructing $`p`$-adic mathematical physics. By now, a considerable progress has been made in this direction (see , where the motivation, the history of the problem, and extensive bibliography were presented).
The construction of adelic formulas relating $`N`$-particle scattering amplitudes to appropriate $`p`$-adic amplitudes constitutes a separate field of $`p`$-adic mathematical physics. The first such formulas (without regularization) were proposed by Freund and Witten and, independently, by Volovich in 1987 as applied to 4-tachyon string amplitudes. Mathematically, the problem is reduced to the construction of adelic formulas for local gamma and beta functions of algebraic number fields. In , such regularized adelic formulas were constructed in arbitrary fields of algebraic numbers for principal (unramified) quasicharacters and in , for ramified ones. In the latter case, these formulas are valid under the condition that the quasicharacters of the appropriate group of ideles are trivial (i.e., are equal to 1) on the multiplicative group of the original field of algebraic numbers . This condition yields explicit relations between the parameters of local characters, similar to the formulas for a field of rational numbers $``$ . These relations for one-class quadratic fields were obtained in my work (see Section 2). In Section 3, we present regularized adelic formulas for gamma and beta functions as applied to one-class quadratic fields (and to the field of rational numbers). As compared with , here we remove the technical requirements that the ranks of local characters should be identical at all finite points (see Remark 5). In Section 4, we apply the formulas obtained to 4-tachyon tree string amplitudes: to the Veneziano (open strings) and Virasoro (closed strings) amplitudes as well as to the massless 4-particle scattering amplitude of the Ramond–Neveu–Schwarz superstring and massless scattering amplitudes of four charged particles of a heterotic string. We establish certain relations between different superstring amplitudes. We also discuss the problem of other possible applications of adelic formulas for ramified quasicharacters.
## 2 IDELES AND THEIR QUASICHARACTERS FOR <br>ONE-CLASS QUADRATIC FIELDS
First, recall the necessary information from the theory of adeles (and ideles) as applied to the quadratic fields $`(\sqrt{d})`$, where $`dZ`$, $`d0,1`$, is squares free . For this purpose, we first introduce the following notation.
Denote by $`P_d`$ and $`S_d`$ the sets of prime numbers $`p`$ such that $`d_p^{\times 2}`$ or $`d_p^{\times 2}`$, respectively.
Let $`A_d`$ be the adele ring and $`A_d^\times `$ be the idele group of the field $`(\sqrt{d})`$. An arbitrary adele $`XA_d`$ is expressed as
$$X=(X_{\mathrm{}},X_2,X_3,\mathrm{},X_p,\mathrm{}),$$
where the symbol $`X_p`$ determines the following points of the ring $`A_d`$:
$`X_{\mathrm{}}`$ $`=\{\begin{array}{ccc}z,\hfill & & d<0,\hfill \\ (x,x^{})^2,\hfill & & d>0;\hfill \end{array}`$
$`X_p`$ $`=\{\begin{array}{ccc}(x_p,x^{}_p)_p^2,\hfill & & pP_d,\hfill \\ x_p+\sqrt{d}x^{}_p_p(\sqrt{d}),\hfill & & pS_d;\hfill \end{array}`$
in this case, for all $`p`$ starting from a certain value, $`(x_p,x_p^{})Z_p^2`$, $`pP_d`$, and $`X_pZ_p(\sqrt{d})`$, $`pS_d`$. Here, $`Z`$$`Z_p,`$ and $`Z_p(\sqrt{d})`$ are the rings of integers of the fields $``$$`_p`$, and $`_p(\sqrt{d})`$, respectively. The group of ideles $`A_d^\times `$ is defined similarly. Ideles are invertible adeles.
An arbitrary (multiplicative) quasicharacter $`\mathrm{\Theta }`$ of the group $`A_d^\times `$ on the idele $`X`$ is represented as
$$\mathrm{\Theta }(X;\alpha )=\mathrm{\Theta }_{\mathrm{}}(X_{\mathrm{}})\underset{pP_d}{}\theta _p(x_p)\theta _p^{}(x_p^{})|x_p|_p^{i\alpha _p}|x_p^{}|_p^{i\alpha _p^{}}\underset{pS_d}{}\theta _p(X_p)|X_p\overline{X}_p|_p^{i\alpha _p}|X|^\alpha ,\alpha ,$$
(2.1)
where $`\overline{X}_p=x_p\sqrt{d}x_p^{}`$, $`X_p\overline{X}_p=x_p^2dx_{p}^{}{}_{}{}^{2}`$, and
$$\mathrm{\Theta }_{\mathrm{}}(X_{\mathrm{}})=\{\begin{array}{ccc}z^\nu (z\overline{z})^{\nu /2},\hfill & & \nu Z,d<0,\hfill \\ sgn^\nu xsgn^\nu ^{}x^{}|x|^{ia},\hfill & & \nu ,\nu ^{}F_2,a,d>0.\hfill \end{array}$$
(2.2)
Here, $`|X|`$ denotes the normalization of the idele $`X`$ and $`F_2`$ is a field of integers modulo 2.
Suppose that the local characters $`\theta _p`$ and $`\theta _p^{}`$ in (2.1) are normalized by the conditions
$$\theta _p(p)=1,\theta _p^{}(p)=1,pP_d;\theta _p(q)=1,pS_d,$$
(2.3)
where $`q=q_p`$ is the module over the field $`_p(\sqrt{d})`$, $`pS_d`$.
In (2.1), only a finite number of characters $`\theta _p`$ and $`\theta _p^{}`$ are nontrivial (ramified). Denote the ranks of these characters by $`\rho _p=\rho (\theta _p)`$ and $`\rho _p^{}=\rho (\theta _p^{})`$, respectively. Denote by $`F`$ the set of finite points for which the rank of the local character is equal to 0 (unramified points); the set of other finite points is denoted by $`R`$ (ramified points).
The triviality of the quasicharacter $`\mathrm{\Theta }`$ on the principal idele $`X=I_x`$ generated by the number $`x^\times (\sqrt{d})`$ is expressed as
$$\mathrm{\Theta }(I_x;\alpha )\mathrm{\Theta }(x)=1,x^\times (\sqrt{d}).$$
(2.4)
Condition (2.4) implies that, actually, the quasicharacter $`\mathrm{\Theta }`$ is defined over the factor group $`A_d^\times /^\times (\sqrt{d})`$, the idele class group over the field $`(\sqrt{d})`$ .
Denote the norm of the leading ideal $`J`$ of the quasicharacter $`\mathrm{\Theta }`$ by $`N`$ :
$$N=N(J)=\underset{pP_dR}{}p^{\rho _p+\rho _p^{}}\underset{pS_dR}{}q_p^{\rho _p}.$$
(2.5)
A standard additive character $`\chi `$ of the adele ring $`A_d`$ is given by
$$\chi (X)=\chi _{\mathrm{}}(X_{\mathrm{}})\underset{pP_d}{}\chi _p(x_p+x_p^{})\underset{pS_p}{}\chi _p(X_p+\overline{X}_p),XA_d,$$
(2.6)
where
$$\begin{array}{c}\chi _{\mathrm{}}(X_{\mathrm{}})=\{\begin{array}{ccc}\mathrm{exp}[2\pi i(z+\overline{z})],\hfill & & d<0,\hfill \\ \mathrm{exp}[2\pi i(x+x^{})],\hfill & & d>0;\hfill \end{array}\\ \chi _p(x)=\mathrm{exp}(2\pi i\{x\}_p),x_p,\end{array}$$
where $`\{x\}_p`$ is the fractional part of the $`p`$-adic number $`x`$.
Denote by $`r_p`$ the rank of the local character $`\chi _p(X+\overline{X}),`$ $`X_p(\sqrt{d})`$, $`pS_d`$. Only a finite number of these ranks are different from 0.
Denote by $`D`$ the discriminant of the field $`(\sqrt{d})`$:
$$D=\{\begin{array}{ccc}d,\hfill & & d1(mod4),\hfill \\ 4d,\hfill & & d2,3(mod4),\hfill \end{array}|D|=\underset{pS_d}{}q_p^{r_p}.$$
(2.7)
The unities $`ϵ`$ of the quadratic fields $`(\sqrt{d})`$ are $`ϵ=\pm 1`$ and, in addition,
$`ϵ=\pm i\text{ for }d=1,`$
$`ϵ=\pm e^{\pm i\pi /3}=\pm (1\pm \sqrt{3})/2\text{ for }d=3,`$
$`ϵ=\pm \mathrm{\Omega }^\gamma ,\gamma Z,\text{ for }d>0,`$
where $`\mathrm{\Omega }`$ is the primary unity of the field $`(\sqrt{d})`$, $`N(\mathrm{\Omega })=\mathrm{\Omega }\overline{\mathrm{\Omega }}=\pm 1`$, $`\mathrm{\Omega }>0`$.
All one-class quadratic fields $`(\sqrt{d})`$ are well known : there exist nine imaginary fields of this type with $`d=1,2,3,7,11,19,43,67,163`$ and an infinite number of real fields with $`d=p`$, $`p3(mod4)`$; $`d=2p`$, $`p3(mod4)`$; and $`d=pp^{}`$, $`p,p^{}3(mod4)`$ ($`p`$ and $`p^{}`$ are prime numbers).
For one-class fields, the group of classes of divisors is trivial; i.e., all these divisors belong to the ring of integers (to the maximum order), and the prime factor decomposition in this field is unique .
###### Proposition (see ).
All prime divisors $`𝔭𝔭_p`$ and $`\overline{𝔭}\overline{𝔭}_p`$ of one-class quadratic fields $`(\sqrt{d})`$ and their norms $`qq_p`$ are classified under the following cases$`:`$
(a) $`p`$ divides $`D,`$ $`pS_d,`$$`p=ϵ𝔭^2,`$$`q=p;`$ $`(2.8)`$
(b) $`p`$ does not divide $`D,`$ $`pS_d,`$$`p=𝔭,`$$`q=p^2;`$ $`(2.9)`$
$`(\mathrm{c}^{})`$ $`p`$ does not divide $`D,`$ $`pP_d,`$ $`D=4d,`$$`𝔭,\overline{𝔭}=a\pm \sqrt{d}b,`$$`q=p,`$ $`(2.10)`$
where $`a>0`$ and $`b>0`$ are integer solutions to the Diophantine equation
$$p=|x^2dy^2|;$$
(2.11)
$`(\mathrm{c}^{\prime \prime })`$ $`p`$ does not divide $`D,`$ $`pP_d,`$ $`D=d,`$$`𝔭,\overline{𝔭}=a+{\displaystyle \frac{b}{2}}\pm \sqrt{d}{\displaystyle \frac{b}{2}},`$$`q=p,`$ $`(2.12)`$
where $`a>0`$ and $`b0`$ are integer solutions to the Diophantine equation
$$p=\left|x^2+xy+\frac{1d}{4}y^2\right|$$
(2.13)
$`(`$for $`d=7`$ and $`p=2`$, the solution is $`a=0,b=1)`$.
Here, in cases (c) and (c<sup>′′</sup>), the divisors $`𝔭`$ and $`\overline{𝔭}`$ can be chosen so that the conditions
$$|𝔭|_p=1/p,|\overline{𝔭}|_p=1$$
(2.14)
are fulfilled.
###### Remark 1.
For one-class quadratic fields, the Diophantine equations (2.11) and (2.13) are solvable for any prime $`p`$.
The following two theorems give the necessary and sufficient conditions for the triviality of the quasicharacter of the idele group of a field on the multiplicative group of this field.
###### Theorem 1 (see ).
The quasicharacter
$$\mathrm{\Theta }(X;\alpha )=sgn^\nu x_{\mathrm{}}\underset{p=2}{\overset{\mathrm{}}{}}\theta _p(x_p)|x_p|_p^{i\alpha _p}|X|^\alpha ,\nu F_2,\alpha ,$$
(2.15)
$`(\theta _p(p)=1)`$ of the group $`A^\times `$ of ideles $`X=(x_{\mathrm{}},x_2,\mathrm{},x_p,\mathrm{})`$ of the field $``$ is trivial on the group $`^\times `$ if and only if the following conditions hold$`:`$
$$\theta (1)=1,\theta (p)=p^{i\alpha _p},p=2,3,\mathrm{},$$
(2.16)
where $`\theta `$ is a character of the group $`A^\times `$ of the form
$$\theta (X)=sgn^\nu x_{\mathrm{}}\underset{p=2}{\overset{\mathrm{}}{}}\theta _p(x_p).$$
(2.17)
Example of a ramified character of the group $`A^\times /^\times `$:
$$\theta (X)=sgnx_{\mathrm{}}\theta _2(x_2)\underset{p=3}{\overset{\mathrm{}}{}}|x_p|^{i\alpha _p},\theta _2(1)=1,\theta _2(p)=p^{i\alpha _p},p=3,5,\mathrm{}.$$
###### Theorem 2 (see ).
The quasicharacter $`\mathrm{\Theta }`$ of the idele group $`A_d^\times `$ of the one-class quadratic field $`(\sqrt{d})`$ is trivial on the group $`^\times (\sqrt{d})`$ if and only if the following conditions hold$`:`$
$$\theta (1)=1\text{ for all admissible }d;$$
(2.18)
$$\theta (i)=1\text{ for }d=1;$$
(2.19)
$$\theta (e^{i\pi /3})=1\text{ for }d=3;$$
(2.20)
$$\theta (\mathrm{\Omega })=\mathrm{\Omega }^{ia}\text{ for all admissible }d>0;$$
(2.21)
$$p^{i\alpha _p}=\theta (𝔭)|𝔭|^{iah(d)},p^{i\alpha _p^{}}=\theta (\overline{𝔭})|\overline{𝔭}|^{iah(d)},pP_d;$$
(2.22)
$$p^{i\alpha _p}=\theta (𝔭)|𝔭|^{iah(d)},p\text{ divides }D,pS_d;$$
(2.23)
$$p^{2i\alpha _p}=\theta (p)p^{iah(d)},p\text{ does not divide }D,pS_d;$$
(2.24)
here $`\theta `$ denotes a character of the idele group $`A_d^\times ,`$
$$\theta (X)=\theta _{\mathrm{}}(X_{\mathrm{}})\underset{pP_dR}{}\theta _p(x_p)\theta _p^{}(x_p^{})\underset{pS_dR}{}\theta _p(X_p),$$
(2.25)
$$\theta _{\mathrm{}}(X_{\mathrm{}})=\{\begin{array}{ccc}z^\nu (z\overline{z})^{\nu /2},\hfill & & \nu Z,d<0,\hfill \\ sgn^\nu xsgn^\nu ^{}x^{},\hfill & & \nu ,\nu ^{}F_2,d>0\hfill \end{array}$$
(2.28)
$`(h(d)=1`$ for $`d>0,`$ $`h(d)=0`$ for $`d<0)`$.
###### Remark 2.
Equation (2.21) ambiguously determines the real number $`a`$:
$$a=\frac{Arg\theta (\mathrm{\Omega })}{\mathrm{ln}\mathrm{\Omega }}+\frac{2\pi }{\mathrm{ln}\mathrm{\Omega }}Z.$$
Therefore, the real numbers $`\alpha _p`$ and $`\alpha _p^{}`$ are determined by equalities (2.22)–(2.24) ambiguously. However, in what follows, we will need only the exponential functions $`p^{i\alpha _p}`$ and $`p^{i\alpha _p^{}}`$, $`pP_d`$, and $`q^{i\alpha _p}`$, $`pS_d`$, of these numbers; for $`d>0`$, these exponential functions are determined uniquely.
###### Remark 3.
For the principal quasicharacter $`|X|^\alpha `$, $`\alpha `$, conditions (2.18)–(2.24) hold trivially by virtue of the known Artin multiplication formula $`|I_x|=1`$, $`x_p^\times (\sqrt{d})`$.
## 3 ADELIC FORMULAS FOR GAMMA AND BETA <br>FUNCTIONS OF ONE-CLASS QUADRATIC FIELDS <br>AND A FIELD OF RATIONAL NUMBERS
Taking into account equality (2.18), we can represent the regularized adelic formulas of one-class quadratic fields $`(\sqrt{d})`$ for gamma functions of the quasicharacter $`\mathrm{\Theta }(X;\alpha )`$ (defined by (2.1)) of the idele group $`A_d^\times /^\times (\sqrt{d})`$ as
$$\theta (1)=reg\left[\underset{p=2}{\overset{\mathrm{}}{}}\mathrm{\Gamma }(\alpha +i\alpha _p;\theta _p)\underset{pP_dF}{}\mathrm{\Gamma }(\alpha +i\alpha _p^{};\theta _p^{})\right]\{\begin{array}{ccc}\mathrm{\Gamma }_\omega (\alpha ;\nu ),\hfill & & d<0,\hfill \\ \mathrm{\Gamma }_{\mathrm{}}(\alpha +ia;\nu )\mathrm{\Gamma }_{\mathrm{}}(\alpha ;\nu ^{}),\hfill & & d>0,\hfill \end{array}$$
(3.1)
where the character $`\theta (X)`$ is defined by (2.25) and
$$\mathrm{\Gamma }_\omega (\alpha ;\nu )=_{}z^\nu (z\overline{z})^{\nu /2+\alpha 1}\mathrm{exp}[2\pi i(z+\overline{z})]|dzd\overline{z}|$$
$$=i^\nu 2(2\pi )^{2\alpha }\mathrm{\Gamma }\left(\alpha +\frac{\nu }{2}\right)\mathrm{\Gamma }\left(\alpha \frac{\nu }{2}\right)\mathrm{sin}\frac{\pi }{2}(2\alpha \nu ),\nu Z,$$
(3.2)
is the gamma function of quasicharacter $`z^\nu (z\overline{z})^{\nu /2+\alpha }`$ of the field $``$;
$$\mathrm{\Gamma }_{\mathrm{}}(\alpha ;\nu )=_{}sgn^\nu x|x|^{\alpha 1}\mathrm{exp}(2\pi ix)dx=2i^\nu (2\pi )^\alpha \mathrm{\Gamma }(\alpha )\mathrm{cos}\frac{\pi }{2}(\alpha \nu ),\nu F_2,$$
(3.3)
is the gamma function of quasicharacter $`sgn^\nu x|x|^\alpha `$ of the field $``$ ($`\mathrm{\Gamma }(\alpha )`$ is the Euler gamma function);
$$\mathrm{\Gamma }(\alpha ;\theta )=q^{\frac{r}{2}}_𝕂\theta (x)|x|^{\alpha 1}\chi (x)dx=\{\begin{array}{ccc}q^{(\alpha 1/2)r}\mathrm{\Gamma }_q(\alpha ),\hfill & & pF,\hfill \\ \varkappa (\theta )q^{(\alpha 1/2)(r+\rho )},\hfill & & pR;\hfill \end{array}$$
(3.4)
$$\mathrm{\Gamma }_q(\alpha )=G_q(q^\alpha ),\alpha \frac{2\pi i}{\mathrm{ln}q}Z,G_q(x)=\frac{1x/q}{1x^1},$$
(3.4)
is the reduced gamma function of a local $`p`$-field with the module $`q`$; and
$$\varkappa (\theta )=q^{\rho /2}_{|x|=1}\theta (x)\chi (\pi ^{r\rho }x)𝑑x,|\varkappa (\theta )|=1.$$
By Theorem 2 (see Section 2), equalities (2.18)–(2.24) are valid for one-class quadratic fields. Using these equalities, we eliminate the numbers $`\alpha _p`$ and $`\alpha _p^{}`$ from the gamma functions in adelic formulas (3.1). As a result, we obtain the adelic formulas
$$\begin{array}{c}\theta (1)\varkappa (N|D|)^{1/2\alpha }\\ =reg\underset{pF}{}\mathrm{\Gamma }_q(\alpha +i\alpha _p)\underset{pP_dF}{}\mathrm{\Gamma }_p(\alpha +i\alpha _p^{})\{\begin{array}{ccc}\mathrm{\Gamma }_\omega (\alpha ;\nu ),\hfill & & d<0,\hfill \\ \mathrm{\Gamma }_{\mathrm{}}(\alpha +ia,\nu )\mathrm{\Gamma }_{\mathrm{}}(\alpha ,\nu ^{}),\hfill & & d>0,\hfill \end{array}\end{array}$$
(3.5)
where the numbers $`N`$ and $`D`$ are defined by formulas (2.5) and (2.7),
$$\varkappa =\underset{p=2}{\overset{\mathrm{}}{}}\varkappa (\theta _p)q_p^{i\alpha _p(r_p+\rho _p)}\underset{pP_dR}{}\varkappa (\theta _p^{})p^{i\alpha _p^{}\rho _p^{}},|\varkappa |=1,$$
(3.6)
where the quantities $`q_p^{i\alpha _p}`$ and $`p^{i\alpha _p^{}}`$ are determined by equalities (2.22)–(2.24), while $`\mathrm{\Gamma }_q(\alpha +i\alpha _p)`$ and $`\mathrm{\Gamma }_p(\alpha +i\alpha _p^{})`$ are determined by the formulas
$`\mathrm{\Gamma }_p(\alpha +i\alpha _p)`$ $`=G_p[p^\alpha \theta (𝔭)𝔭^{ih(d)a}],`$ (3.7)
$`\mathrm{\Gamma }_p(\alpha +i\alpha _p^{})`$ $`=G_p[p^\alpha \theta (\overline{𝔭})\overline{𝔭}^{ih(d)a}],pP_dF,`$ (3.8)
$`\mathrm{\Gamma }_q(\alpha +i\alpha _p)`$ $`=G_q[q^\alpha \theta (𝔭)p^{ih(d)a}],pS_dF.`$ (3.9)
(In (3.6), we assume that $`\varkappa (\theta _p)=1`$ for $`pF`$.)
Suppose that, in addition to the quasicharacter $`\mathrm{\Theta }(X;\alpha )`$, another quasicharacter $`\mathrm{\Pi }(Y;\beta )`$ of the idele group $`A_d^\times /^\times (\sqrt{d})`$ is defined,
$$\mathrm{\Pi }(Y;\beta )=\pi _{\mathrm{}}(Y_{\mathrm{}})\underset{pP_d}{}\pi _p(y_p)\pi _p^{}(y_p^{})|y_p|_p^{i\beta _p}|y_p^{}|_p^{i\beta _p^{}}\underset{pS_d}{}\pi _p(Y_p)|Y_p\overline{Y}_p|_p^{i\beta _p}|Y|^\beta ,\beta ,$$
(3.10)
where $`b,\beta _p,`$ and $`\beta _p^{}`$ are real numbers and
$$\pi _{\mathrm{}}(Y_{\mathrm{}})=\{\begin{array}{ccc}\zeta ^\mu (\zeta \overline{\zeta })^{\mu /2},\hfill & & \mu Z,d<0,\hfill \\ sgn^\mu ysgn^\mu ^{}y^{}|y|^{ib},\hfill & & \mu ,\mu ^{}F_2,d>0.\hfill \end{array}$$
(3.11)
Let us construct, by formulas (2.25) and (2.26), the character $`\pi (Y)`$ corresponding to the quasicharacter $`\mathrm{\Pi }`$ and then the character $`\sigma (Z)=\overline{\theta }(Z)\overline{\pi }(Z)`$, so that
$$\theta \pi \sigma =1,\sigma _p=\overline{\theta }_p\overline{\pi }_p,p=2,3,\mathrm{},\sigma _p^{}=\overline{\theta }_p^{}\overline{\pi }_p^{},pP_d.$$
(3.12)
Denote by $`R,R^{},R^{\prime \prime }`$ and $`F,F^{},F^{\prime \prime }`$ the sets of ramified and unramified points of the characters $`\theta ,\pi ,`$ and $`\sigma `$, respectively.
The beta functions of local fields are defined as the following products of three gamma functions:
$$\begin{array}{c}\mathrm{B}_\omega (\alpha ,\nu ;\beta ,\mu ;\gamma ,\eta )=\mathrm{\Gamma }_\omega (\alpha ;\nu )\mathrm{\Gamma }_\omega (\beta ;\mu )\mathrm{\Gamma }_\omega (\gamma ;\eta ),\\ \alpha +\beta +\gamma =1,\nu +\mu +\eta =0,\nu ,\mu ,\eta Z,\end{array}$$
(3.13)
for the quasicharacters $`z^\nu (z\overline{z})^{\nu /2+\alpha }`$ and $`\zeta ^\mu (\zeta \overline{\zeta })^{\mu /2+\beta }`$, $`\nu ,\mu Z`$, of the field $``$;
$$\begin{array}{c}\mathrm{B}_{\mathrm{}}(\alpha ,\nu ;\beta ,\mu ;\gamma ,\eta )=\mathrm{\Gamma }_{\mathrm{}}(\alpha ;\nu )\mathrm{\Gamma }_{\mathrm{}}(\beta ;\mu )\mathrm{\Gamma }_{\mathrm{}}(\gamma ;\eta ),\\ \alpha +\beta +\gamma =1,\nu +\mu +\eta =0,\nu ,\mu ,\eta F_2,\end{array}$$
(3.14)
for the quasicharacters $`sgn^\nu x|x|^\alpha `$ and $`sgn^\mu y|y|^\beta `$, $`\nu ,\mu F_2`$, of the field $``$; and
$$\mathrm{B}(\alpha ,\theta ;\beta ,\pi ;\gamma ,\sigma )=\mathrm{\Gamma }(\alpha ;\theta )\mathrm{\Gamma }(\beta ;\pi )\mathrm{\Gamma }(\gamma ;\sigma ),\alpha +\beta +\gamma =1,\theta \pi \sigma =1,$$
(3.15)
for the local quasicharacters $`|x|^\alpha \theta (x)`$ and $`|y|^\beta \pi (y)`$ of the $`p`$-fields.
In particular, for the principal quasicharacters $`|x|_p^\alpha `$, $`|y|_p^\beta `$, $`pP_dF`$, $`|X\overline{X}|_p^\alpha `$, $`|Y\overline{Y}|_p^\beta `$, $`pS_dF`$, of the $`p`$-fields with the module $`q=q_p`$, the beta function is defined as
$$\mathrm{B}_q(\alpha ,\beta ,\gamma )=\mathrm{\Gamma }_q(\alpha )\mathrm{\Gamma }_q(\beta )\mathrm{\Gamma }_q(\gamma ).$$
(3.16)
By virtue of (3.1) and (3.13)–(3.15), the local beta functions defined by the quasicharacters $`\mathrm{\Theta }(X;\alpha )`$ and $`\mathrm{\Pi }(Y;\beta )`$ of the idele group $`A_d^\times /^\times (\sqrt{d})`$ satisfy the following adelic equality (for $`\alpha +\beta +\gamma =1`$):
$$\begin{array}{c}1=reg[\underset{p=2}{\overset{\mathrm{}}{}}\mathrm{B}(\alpha +i\alpha _p,\theta _p;\beta +i\beta _p,\pi _p;\gamma +i\gamma _p,\sigma _p)\\ \times \underset{pP_d(FF^{}F^{\prime \prime })}{}\mathrm{B}(\alpha +i\alpha _p^{\mathrm{}},\theta _p^{\mathrm{}};\beta +i\beta _p^{\mathrm{}},\pi _p^{\mathrm{}};\gamma +i\gamma _p^{\mathrm{}},\sigma _p^{\mathrm{}})]\\ \times \{\begin{array}{c}\mathrm{B}_\omega (\alpha ,\nu ;\beta ,\mu ;\gamma ,\eta ),\nu +\mu +\eta =0,\nu ,\mu ,\eta Z,d<0,\hfill \\ \mathrm{B}_{\mathrm{}}(\alpha +ia,\nu ;\beta +ib,\mu ;\gamma +ic,\eta )\mathrm{B}_{\mathrm{}}(\alpha ,\nu ^{};\beta ,\mu ^{};\gamma ,\eta ^{}),\hfill \end{array}\\ \nu +\mu +\eta =\nu ^{}+\mu ^{}+\eta ^{}=0,\nu ,\mu ,\eta ,\nu ^{},\mu ^{},\eta ^{}F_2,\\ a+b+c=0,a,b,c,d>0;\end{array}$$
(3.17)
here, $`\alpha _p^{\mathrm{}}=\alpha _p^{}`$ and $`\theta _p^{\mathrm{}}=\theta _p^{}`$ if $`pF`$, $`\alpha _p^{\mathrm{}}=\alpha _p`$ and $`\theta _p^{\mathrm{}}=\theta _p`$ if $`pR`$, etc. The real numbers $`\alpha _p,\beta _p,\gamma _p,\alpha _p^{},\beta _p^{},`$ and $`\gamma _p^{}`$ in (3.17) must satisfy the relations
$$\alpha _p+\beta _p+\gamma _p=0,p=2,3,\mathrm{},\alpha _p^{}+\beta _p^{}+\gamma _p^{}=0,pP_d.$$
Using the adelic formula (3.5), we rewrite (3.17) in expanded form in terms of the reduced beta function $`\mathrm{B}_q`$; the real numbers $`\alpha _p,\beta _p,\gamma _p,\alpha _p^{},\beta _p^{},`$ and $`\gamma _p^{}`$ are eliminated from these beta functions by the following relations:
$$\mathrm{B}_p(\alpha +i\alpha _p,\beta +i\beta _p,\gamma +i\gamma _p)$$
$$=G_p[p^\alpha \theta (𝔭)𝔭^{ih(d)a}]G_p[p^\beta \pi (𝔭)𝔭^{ih(d)b}]G_p[p^\gamma \sigma (𝔭)𝔭^{ih(d)c}],pP_dF,$$
(3.18)
$$\mathrm{B}_p(\alpha +i\alpha _p^{},\beta +i\beta _p^{},\gamma +i\gamma _p^{})$$
$$=G_p[p^\alpha \theta (\overline{𝔭})\overline{𝔭}^{ih(d)a}]G_p[p^\beta \pi (\overline{𝔭})\overline{𝔭}^{ih(d)b}]G_p[p^\gamma \sigma (\overline{𝔭})\overline{𝔭}^{ih(d)c}],pP_dF,$$
(3.19)
$$\mathrm{B}_q(\alpha +i\alpha _p,\beta +i\beta _p,\gamma +i\gamma _p)$$
$$=G_q[q^\alpha \theta (𝔭)p^{ih(d)a}]G_q[q^\beta \pi (𝔭)p^{ih(d)b}]G_q[q^\gamma \sigma (𝔭)p^{ih(d)c}],pS_dF.$$
(3.20)
Note that, for $`d>0`$, the following equalities hold in (3.17) in accordance with the formulas (2.21) and (3.12):
$$\mathrm{\Omega }^{ia}=\overline{\theta }(\mathrm{\Omega }),\mathrm{\Omega }^{ib}=\overline{\pi }(\mathrm{\Omega }),\mathrm{\Omega }^{ic}=\overline{\sigma }(\mathrm{\Omega })=\theta (\mathrm{\Omega })\pi (\mathrm{\Omega });$$
(3.21)
therefore, $`a+b+c=0`$ in this case as well.
For the principal quasicharacters $`|X|^\alpha `$ and $`|Y|^\beta `$ ($`\nu =\mu =\eta =0`$, $`\nu ^{}=\mu ^{}=\eta ^{}=0`$, $`a=b=c=0`$, $`\alpha _p=\beta _p=\gamma _p=0`$, $`\alpha _p^{}=\beta _p^{}=\gamma _p^{}=0`$, $`\varkappa =1`$, $`N(J)=1`$, and $`R`$ is an empty set), the adelic formula (3.17) for an arbitrary quadratic field $`(\sqrt{d})`$ is rewritten as
$$\sqrt{|D|}=reg\left[\underset{pP_d}{}\mathrm{B}_p^2(\alpha ,\beta ,\gamma )\underset{pS_d}{}\mathrm{B}_q(\alpha ,\beta ,\gamma )\right]\{\begin{array}{ccc}\mathrm{B}_\omega (\alpha ,\beta ,\gamma ),\hfill & & \alpha +\beta +\gamma =1,d<0,\hfill \\ \mathrm{B}_{\mathrm{}}^2(\alpha ,\beta ,\gamma ),\hfill & & \alpha +\beta +\gamma =1,d>0,\hfill \end{array}$$
(3.22)
where
$`\mathrm{B}_{\mathrm{}}(\alpha ,\beta ,\gamma )`$ $`=\mathrm{B}_{\mathrm{}}(\alpha ,0;\beta ,0;\gamma ,0),`$ (3.23)
$`\mathrm{B}_\omega (\alpha ,\beta ,\gamma )`$ $`=\mathrm{B}_\omega (\alpha ,0;\beta ,0;\gamma ,0).`$ (3.23)
In particular, for the Gauss field $`(\sqrt{1})`$, this formula is represented as
$$\mathrm{B}_\omega (\alpha ,\beta ,\gamma )\mathrm{B}_2(\alpha ,\beta ,\gamma )reg\left[\underset{p1(4)}{}\mathrm{B}_p^2(\alpha ,\beta ,\gamma )\underset{p3(4)}{}\mathrm{B}_{p^2}(\alpha ,\beta ,\gamma )\right]=2,\alpha +\beta +\gamma =1.$$
(3.24)
For the field $``$ of rational numbers, we proceed similarly and in a more simple manner.
By Theorem 1 (see Section 2), the adelic formula for the gamma functions of the field $``$ that are determined by the quasicharacter $`\mathrm{\Theta }(X;\alpha )`$ of the form (2.15) is given by (cf. (3.5))
$$\mathrm{\Gamma }_{\mathrm{}}(\alpha ;\nu )reg\underset{pF}{}\mathrm{\Gamma }_p(\alpha +i\alpha _p)=\theta (1)\underset{pR}{}\varkappa (\theta _p)p^{i\alpha _p\rho _p}N^{1/2\alpha },$$
(3.25)
where
$$\mathrm{\Gamma }_p(\alpha +i\alpha _p)=\frac{1p^{\alpha 1}\theta (p)}{1p^\alpha \overline{\theta }(p)},pF.$$
(3.26)
(The character $`\theta (X)`$ is defined by (2.17).)
For the quasicharacters $`\mathrm{\Theta }(X;\alpha )`$ and $`\mathrm{\Pi }(Y;\beta )`$ of the form (2.15), the adelic formula for the beta functions of the field $``$ is given by
$$\begin{array}{c}\mathrm{B}_{\mathrm{}}(\alpha ,\nu ;\beta ,\mu ;\gamma ,\eta )reg\underset{pFF^{}F^{\prime \prime }}{}\mathrm{B}_p(\alpha +i\alpha _p,\beta +i\beta _p,\gamma +i\gamma _p)=\varkappa (NN^{}N^{\prime \prime })^{1/2\alpha },\\ \alpha +\beta +\gamma =1,\nu +\mu +\eta =0,\nu ,\mu ,\eta F_2,\end{array}$$
(3.27)
where the characters $`\theta `$ and $`\pi `$ are constructed by the quasicharacters $`\mathrm{\Theta }`$ and $`\mathrm{\Pi }`$ by formula (2.17), $`\theta \pi \sigma =1`$; $`\rho _p,\rho _p^{},\rho _p^{\prime \prime }`$ and $`N,N^{},N^{\prime \prime }`$ are the ranks and the norms of the leading ideals of the characters $`\theta ,\pi ,`$ and $`\sigma `$ of the field $``$, respectively; and
$$\varkappa =\underset{pR}{}\varkappa (\theta _p)p^{i\alpha _p\rho _p}\underset{pF^{}}{}\varkappa (\pi _p)p^{i\beta _p\rho _p^{}}\underset{pF^{\prime \prime }}{}\varkappa (\sigma _p)p^{i\gamma _p\rho _p^{\prime \prime }},|\varkappa |=1.$$
(3.28)
By virtue of (2.16), (3.16), and (3.26),
$$\mathrm{B}_p(\alpha +\alpha _p,\beta +\beta _p,\gamma +\gamma _p)=\frac{1p^{\alpha 1}\theta (p)}{1p^\alpha \overline{\theta }(p)}\frac{1p^{\beta 1}\pi (p)}{1p^\beta \overline{\pi }(p)}\frac{1p^{\gamma 1}\sigma (p)}{1p^\gamma \overline{\sigma }(p)},pF,$$
(3.29)
in formula (3.27).
For the principal quasicharacters $`|X|^\alpha `$ and $`|Y|^\beta `$ ($`\nu =\mu =\eta =0`$, $`\alpha _p=\beta _p=\gamma _p=0`$, and $`N=1`$) formula (3.27) takes the form
$$\mathrm{B}_{\mathrm{}}(\alpha ,\beta ,\gamma )reg\underset{p=2}{\overset{\mathrm{}}{}}\mathrm{B}_p(\alpha ,\beta ,\gamma )=1,\alpha +\beta +\gamma =1.$$
(3.30)
To sum up, we present the following theorems.
###### Theorem 3 (see ).
Suppose that the quasicharacters $`\mathrm{\Theta }`$ and $`\mathrm{\Pi }`$ of the form $`(2.15)`$ are given on the idele class group of the field $`,`$ and let $`\theta ,\pi ,`$ and $`\sigma ,`$ $`\theta \pi \sigma =1,`$ be the characters, of the form $`(2.17),`$ of the idele group of the field $``$. Then, the adelic formulas (3.25), (3.27) are valid for local gamma and beta functions.
###### Theorem 4.
Suppose that the quasicharacters $`\mathrm{\Theta }`$ and $`\mathrm{\Pi }`$ of the form $`(2.1)`$ are given on the idele class group of the one-class field $`(\sqrt{d}),`$ and let $`\theta ,\pi ,`$ and $`\sigma ,`$ $`\theta \pi \sigma =1,`$ be the characters, of the form $`(2.25),`$ of the idele group of the field $`(\sqrt{d})`$. Then, the adelic formulas $`(3.5)`$ and $`(3.17)`$ are valid for local gamma and beta functions.
###### Remark 4.
The adelic formulas (3.1) and (3.17) for the gamma function $`\mathrm{\Gamma }_{A_d}(\alpha ;\mathrm{})`$ and the beta function $`\mathrm{B}_{A_d}(\alpha ,\beta ,\gamma ;\mathrm{})`$ of the adele ring $`A_d`$ of the field $`(\sqrt{d})`$ imply that $`\mathrm{\Gamma }_{A_d}(\alpha ;\mathrm{})=\theta (1)`$ and $`\mathrm{B}_{A_d}(\alpha ,\beta ,\gamma ;\mathrm{})=1`$, $`\alpha +\beta +\gamma =1`$. Similar expressions are valid for the field $``$.
###### Remark 5.
In , I established the adelic formulas under the condition that the ranks of the local characters $`\theta _p,\pi _p,`$ and $`\sigma _p`$, $`p=2,3,\mathrm{}`$, and $`\theta _p^{},\pi _p^{},`$ and $`\sigma _p^{}`$, $`pP_d`$, are identical:
$$\rho (\theta _p)=\rho (\pi _p)=\rho (\sigma _p),p=2,3,\mathrm{};\rho (\theta _p^{})=\rho (\pi _p^{})=\rho (\sigma _p^{}),pP_d.$$
In the present work, this requirement is removed at the expense of complicating the adelic formulas.
Example of characters $`\theta ,\pi ,`$ and $`\sigma =\overline{\theta }\overline{\pi }`$ with identical ranks for $`p=5`$: $`x_5=2^k`$, $`aZ_5^\times `$, $`|1a|_5<1`$,
$$\theta _5(x_5)=\pi _5(x_5)=\mathrm{exp}(2\pi ik/5),k=0,1,2,3;\rho (\theta _5)=\rho (\pi _5)=\rho (\overline{\theta }_5)=1.$$
## 4 APPLICATION TO 4-PARTICLE STRING <br>AMPLITUDES
Suppose that $`s,t`$, and $`u`$ are the Mandelstam variables for a 4-particle scattering process with the momenta $`(p_1,p_2,p_3,p_4)`$, $`p_1+p_2+p_3+p_4=0`$, $`p_i^2=m_i^2`$ in an $`n`$-dimensional space $`^n`$ with the Minkowskian metric $`p^2=(p^0)^2(p^1)^2\mathrm{}(p^{n1})^2`$, so that
$$s=(p_1+p_2)^2,t=(p_2+p_3)^2,u=(p_1+p_3)^2,s+t+u=m_i^2.$$
(4.1)
In the general case of 4-particle scattering processes with the Regge trajectory $`\alpha (s)=\alpha +\alpha ^{}s`$, we define the generalized crossing-symmetric Veneziano $`V`$ and Virasoro $`W`$ amplitudes and their $`p`$-adic analogues $`V_p`$ and $`W_p`$ by the following formulas $`(s+t+u=m_i^2)`$ :
$`V(s,t,u)`$ $`=\mathrm{B}_{\mathrm{}}(\alpha \alpha ^{}s,\alpha \alpha ^{}t,\alpha \alpha ^{}u);`$ (4.2)
$`V_p(s,t,u)`$ $`=\mathrm{B}_q(\alpha \alpha ^{}s,\alpha \alpha ^{}t,\alpha \alpha ^{}u);`$ (4.3)
$`W(s,t,u)`$ $`=\mathrm{B}_\omega (\alpha /2\alpha ^{}/2s,\alpha /2\alpha ^{}/2t,\alpha /2\alpha ^{}/2u);`$ (4.4)
$`W_p(s,t,u)`$ $`=\mathrm{B}_q(\alpha /2\alpha ^{}/2s,\alpha /2\alpha ^{}/2t,\alpha /2\alpha ^{}/2u).`$ (4.5)
The slope $`\alpha ^{}`$ and the intercept $`\alpha `$ of the trajectory $`\alpha (s)=\alpha +\alpha ^{}s`$ must satisfy the following condition:
$$3\alpha +\alpha ^{}m_i^2=\{\begin{array}{ccc}1\hfill & \text{for amplitude}\hfill & V,\hfill \\ 2\hfill & \text{for amplitude}\hfill & W.\hfill \end{array}$$
(4.6)
By virtue of (3.21), the amplitudes (4.2)–(4.5) satisfy the following adelic relations (provided that $`s+t+u=m_i^2`$):
$`V^2(s,t,u)reg\left[{\displaystyle \underset{pP_d}{}}V_p^2(s,t,u){\displaystyle \underset{pS_d}{}}V_p(s,t,u)\right]`$ $`=\sqrt{|D|},d>0,`$ (4.7)
$`W(s,t,u)reg\left[{\displaystyle \underset{pP_d}{}}W_p^2(s,t,u){\displaystyle \underset{pS_d}{}}W_p(s,t,u)\right]`$ $`=\sqrt{|D|},d<0.`$ (4.8)
Formulas (4.7), (4.8) are well known . They provide a decomposition of generalized Veneziano $`V(d>0)`$ and Virasoro $`W(d<0)`$ amplitudes into the infinite product of the inverses $`V_p^1`$ and $`W_p^1`$ of the $`p`$-adic amplitudes $`V_p`$ and $`W_p`$, respectively.
Let us present certain important particular cases of the generalized amplitudes introduced. These are, first of all, 4-tachyon crossing-symmetric string amplitudes for tree orientable diagrams—the Veneziano and Virasoro amplitudes \[11–13, 21–25\].
### The Veneziano amplitude for open strings in $`^{26}`$.
There exist only one tree orientable diagram that is conformally equivalent to a unit circle with four punctured points on its boundary. This diagrams corresponds to the crossing-symmetric Veneziano amplitude (for $`\alpha =1`$, $`\alpha ^{}=1/2`$, $`m_i^2=2`$, see formula (4.2))
$$V(s,t,u)=\mathrm{B}_{\mathrm{}}(1s/2,1t/2,1u/2),s+t+u=8.$$
(4.9)
The amplitudes $`V_p`$ are determined similarly by formula (4.3).
### The Virasoro amplitude for closed strings in $`^{26}`$.
There exists only one tree orientable diagram that is conformally equivalent to a unit sphere with four punctured points. This diagrams corresponds to the crossing-symmetric Virasoro amplitude (for $`\alpha =2`$, $`\alpha ^{}=1/4`$, $`m_i^2=8`$, see formula (4.4))
$$W(s,t,u)=\mathrm{B}_\omega (1s/8,1t/8,1u/8),s+t+u=32.$$
(4.10)
The amplitudes $`W_p`$ are determined similarly by formula (4.5).
The Veneziano and Virasoro amplitudes are generalized to the case of ramified quasicharacters by the formulas (under conditions (4.1) and (4.6))
$$\begin{array}{c}V_{\nu \mu \eta }(s,t,u)=\mathrm{B}_{\mathrm{}}(\alpha \alpha ^{}s,\nu ;\alpha \alpha ^{}t,\mu ;\alpha \alpha ^{}u,\eta ),\\ \nu +\mu +\eta =0,\nu ,\mu ,\eta F_2,d>0;\end{array}$$
(4.13)
$$\begin{array}{c}W_{\nu \mu \eta }(s,t,u)=\mathrm{B}_\omega (\alpha \alpha ^{}s,\nu ;\alpha \alpha ^{}t,\mu ;\alpha \alpha ^{}u,\eta ),\\ \nu +\mu +\eta =0,\nu ,\mu ,\eta Z,d<0.\end{array}$$
(4.16)
These amplitudes are crossing symmetric with respect to the permutations of the variables $`(s,\nu )`$, $`(t,\mu )`$, and $`(u,\eta )`$ and satisfy the adelic formulas of Section 3 (for one-class quadratic fields and a field of rational numbers).
### The massless 4-particle amplitude of the Ramond–Neveu–Schwarz superstring in $`^{10}`$.
This amplitude is proportional to
$$A_{\mathrm{}}(s,t,u)=\mathrm{\Gamma }_{\mathrm{}}(s/2;1)\mathrm{\Gamma }_{\mathrm{}}(t/2;1)\mathrm{\Gamma }_{\mathrm{}}(u/2;1),s+t+u=0.$$
(4.17)
The corresponding simple adelic formula for $`\nu =\mu =\eta =\mathrm{\hspace{0.17em}1}`$ and $`\rho _p=\rho (\theta _p)=\rho (\pi _p)=\rho (\sigma _p)`$, $`pR=R^{}=R^{\prime \prime }`$, is rewritten as (see (3.25))
$$A_{\mathrm{}}(s,t,u)reg\underset{pF}{}\mathrm{\Gamma }_p(s/2+i\alpha _p)\mathrm{\Gamma }_p(t/2+i\beta _p)\mathrm{\Gamma }_p(u/2+i\gamma _p)=\varkappa N\sqrt{N},s+t+u=0,$$
(4.18)
where $`\varkappa `$ is defined by (3.28); the functions $`\mathrm{\Gamma }_p(s/2+i\alpha _p),\mathrm{}`$ are calculated by (3.26), and the quantities $`p^{i\alpha _p},\mathrm{}`$, by formulas of the type (2.16).
Another adelic formula for the amplitude $`A_{\mathrm{}}(s,t,u)`$ was proposed in .
### Massless amplitudes for four charged particles of a heterotic superstring in $`^{10}`$.
There are four basic types of such amplitudes; namely ,
$$A_{\mathrm{}}^{(k)}(s,t,u)=\mathrm{B}_{\mathrm{}}(1s/8S/2,1t/8T/2,1u/8U/2),$$
(4.19)
where $`s+t+u=0`$ and $`S+T+U=8`$; here, $`k=1`$ corresponds to the set of indices $`(S=8,T=0`$, and $`U=0);`$ $`k=2`$, to $`(S=6,T=2,`$ and $`U=0);`$ $`k=3`$, to $`(S=4,T=4,`$ and $`U=0);`$ and $`k=4`$, to $`(S=4,T=2,U=2)`$. The other amplitudes are obtained by the permutation of indices $`S,T`$, and $`U`$, $`S+T+U=0`$ (that assume the values $`0,2,4,6,`$ and $`8`$).
The amplitudes $`A_{\mathrm{}}^{(k)}(s,t,u)`$, $`k=1,2,3,4`$, are the beta functions of the field $``$ (which are not crossing symmetric). Therefore, they satisfy any adelic formula from Section 3.
Another type of adelic formulas for these amplitudes is obtained if we represent them as the products of the three gamma functions
$$\begin{array}{c}\mathrm{\Gamma }_{\mathrm{}}(s/8;\nu ),\mathrm{\Gamma }_{\mathrm{}}(t/8;\mu ),\mathrm{\Gamma }_{\mathrm{}}(u/8;\eta ),\\ s+t+u=0,\nu +\mu +\eta =1,\nu ,\mu ,\eta F_2,\end{array}$$
and apply the adelic formula (3.25) to each.
For the calculations, we first apply formulas (3.3):
$`S`$ $`=8,\mathrm{\Gamma }_{\mathrm{}}(1s/8+4)=(16\pi )^3i(16s)(8s)s\mathrm{\Gamma }_{\mathrm{}}(s/8;1),`$ (4.20)
$`S`$ $`=6,\mathrm{\Gamma }_{\mathrm{}}(1s/8+3)=(16\pi )^2(8s)s\mathrm{\Gamma }_{\mathrm{}}(s/8;0),`$ (4.21)
$`S`$ $`=4,\mathrm{\Gamma }_{\mathrm{}}(1s/8+2)=(16\pi )^1is\mathrm{\Gamma }_{\mathrm{}}(s/8;1),`$ (4.22)
$`S`$ $`=2,\mathrm{\Gamma }_{\mathrm{}}(1s/8+1)=\mathrm{\Gamma }_{\mathrm{}}(s/8;0),`$ (4.23)
$`S`$ $`=0,\mathrm{\Gamma }_{\mathrm{}}(1s/8+0)={\displaystyle \frac{16\pi i}{8+s}}\mathrm{\Gamma }_{\mathrm{}}(s/8;1),`$ (4.24)
and then formulas (3.23). As a result, we obtain
$`A_{\mathrm{}}^{(1)}(s,t,u)`$ $`={\displaystyle \frac{i}{16\pi }}{\displaystyle \frac{(16s)(8s)s}{(8+t)(8+u)}}\mathrm{\Gamma }_{\mathrm{}}(s/8;1)\mathrm{\Gamma }_{\mathrm{}}(t/8;1)\mathrm{\Gamma }_{\mathrm{}}(u/8;1),`$ (4.25)
$`A_{\mathrm{}}^{(2)}(s,t,u)`$ $`={\displaystyle \frac{i}{16\pi }}{\displaystyle \frac{(8s)s}{8+u}}\mathrm{\Gamma }_{\mathrm{}}(s/8;0)\mathrm{\Gamma }_{\mathrm{}}(t/8;0)\mathrm{\Gamma }_{\mathrm{}}(u/8;1),`$ (4.26)
$`A_{\mathrm{}}^{(3)}(s,t,u)`$ $`={\displaystyle \frac{i}{16\pi }}{\displaystyle \frac{st}{8+u}}\mathrm{\Gamma }_{\mathrm{}}(s/8;1)\mathrm{\Gamma }_{\mathrm{}}(t/8;1)\mathrm{\Gamma }_{\mathrm{}}(u/8;1),`$ (4.27)
$`A_{\mathrm{}}^{(4)}(s,t,u)`$ $`={\displaystyle \frac{is}{16\pi }}\mathrm{\Gamma }_{\mathrm{}}(s/8;1)\mathrm{\Gamma }_{\mathrm{}}(t/8;0)\mathrm{\Gamma }_{\mathrm{}}(u/8;0).`$ (4.28)
Formulas (4.21)–(4.24) and (4.13) show that the amplitudes $`A_{\mathrm{}}^{(1)}(s,t,u)`$ and $`A_{\mathrm{}}^{(3)}(s,t,u)`$ are proportional to the massless superstring amplitude $`A_{\mathrm{}}(s/4,t/4,u/4).`$
Thus, we established a relation between the scattering amplitudes in the theory of a heterotic string and the Ramond–Neveu–Schwarz superstring.
###### Remark 6.
A function of the form (see (4.13), (4.21)–(4.24))
$$\begin{array}{c}\mathrm{B}_{\mathrm{}}^{}(\alpha ,\nu ;\beta ,\mu ;\gamma ,\eta )=\mathrm{\Gamma }_{\mathrm{}}(\alpha ;\nu )\mathrm{\Gamma }_{\mathrm{}}(\beta ;\mu )\mathrm{\Gamma }_{\mathrm{}}(\gamma ;\eta ),\\ \alpha +\beta +\gamma =0,\nu +\mu +\eta =1,\nu ,\mu ,\eta F_2,\end{array}$$
is not a beta function of the type $`\mathrm{B}_{\mathrm{}}`$ that was considered in Section 3. It as if complements $`\mathrm{B}_{\mathrm{}}`$. We denote this function by $`\mathrm{B}_{\mathrm{}}^{}`$ and call a primed beta function. As we have seen, it describes the superstring amplitudes. The integral representation of this function is given by
$$\mathrm{B}_{\mathrm{}}^{}(\alpha ,\nu ;\beta ,\mu ;\gamma ,\eta )=\pi i\left(sgn^\nu x|x|^{\alpha 1}\right)\left(sgn^\mu x|x|^{\beta 1}\right)\left(sgn^\eta x|x|^{\gamma 1}\right)|_{x=1},$$
where $``$ denotes convolution. We will dwell on this point in a different time and different place.
### The amplitudes $`V_{\nu \mu \eta }`$.
In total, there exist four solutions to the equation $`\nu +\mu +\eta =0`$, $`\nu ,\mu ,\eta F_2`$: 000, 110, 101, and 011. Therefore, by (4.11), there exist four different amplitudes $`V_{\nu \mu \eta }`$. One of them, $`V_{000}=V`$, is the Veneziano amplitude, and the rest three quantities $`V_{110},V_{101},`$ and $`V_{011}`$ form a 3-vector in which each component is expressed in terms of another, for example,
$$V_{101}(s,t,u)=V_{110}(s,u,t)=V_{011}(t,s,u);$$
and all of them are expressed in terms of the Veneziano amplitude $`V`$, for example
$$V_{110}(s,t,u)=\frac{1+\alpha +\alpha ^{}t}{\alpha +\alpha ^{}s}V(s\frac{1}{\alpha ^{}},t+\frac{1}{\alpha ^{}},u),$$
(4.29)
by virtue of the easily verifiable relation
$$\mathrm{B}_{\mathrm{}}(\alpha ,1;\beta ,1;\gamma ,0)=\frac{\beta 1}{\alpha }\mathrm{B}_{\mathrm{}}(\alpha +1,\beta 1,\gamma ),\alpha +\beta +\gamma =1.$$
### The amplitudes $`W_{\nu \mu \eta }`$.
There exist an infinite number of amplitudes $`W_{\nu \mu \eta }`$ (4.12), exactly as many as there are integer solutions to the equation $`\nu +\mu +\eta =0`$, $`\nu ,\mu ,\eta Z`$: $`W_{000}=W`$ is the Virasoro amplitude; the quantities $`W_{\nu \nu \mu },W_{\nu \mu \nu }`$, and $`W_{\eta \nu \nu },`$ $`2\nu +\eta =0`$, $`\nu \eta `$, form a 3-vector; and the quantities $`W_{\nu \mu \eta },W_{\nu \eta \mu },\mathrm{},W_{\eta \mu \nu }`$, $`\nu +\mu +\eta =0`$, $`\nu \mu \eta \nu `$, form a 6-vector. In each of these groups of vectors, each component is expressed in terms of another, for example,
$$W_{\mu \nu \eta }(s,t,u)=W_{\nu \mu \eta }(t,s,u)=\mathrm{}=W_{\eta \mu \nu }(u,s,t).$$
The physical sense of these amplitudes is note yet clear.
Everything that was said in this section concerning one-class quadratic fields may also apply to any one-class fields of algebraic numbers. The appropriate adelic formulas were obtained in \[14–17\].
## ACKNOWLEDGMENTS
I am grateful to I.V. Volovich for fruitful discussions.
This work was supported by the Russian Foundation for Basic Research (project no. 96-01-01008) and the Program “Leading Scientific Schools of the Russian Federation” (project no. 96-15-96131).
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# Contact Systems and Corank One Involutive Subdistributions
## Introduction
Consider $`J^n(^k,^m)`$, the space of $`n`$-jets of smooth maps from $`^k`$ into $`^m`$, and denote by
$$(q^1,\mathrm{},q^k,u_1,\mathrm{},u_m,p_i^\sigma ),\text{ for }1im\text{ and for }1\left|\sigma \right|n\text{,}$$
the *canonical coordinates*, also called natural coordinates, on this space (see e.g. , , and ), where $`q^j`$, for $`1jk`$, represent *independent variables* and $`u_i`$, for $`1im`$, represent *dependent variables*; the vector of non-negative integers $`\sigma =(\sigma _1,\mathrm{},\sigma _k)`$ is a multi-index such that $`\left|\sigma \right|=\sigma _1+\mathrm{}+\sigma _kn`$; and $`p_i^\sigma `$, for $`1im`$, correspond to the partial derivatives $`^{\left|\sigma \right|}u_i/q_\sigma `$. Any smooth map $`\phi `$ from $`^k`$ into $`^m`$ defines a submanifold in $`J^n(^k,^m)`$ by the relations
$$p_i^\sigma =\frac{^{\left|\sigma \right|}\phi _i}{q_\sigma }(q^1,\mathrm{},q^k),$$
for $`1im`$ and $`1\left|\sigma \right|n`$. This submanifold is called the $`n`$*-graph* of $`\phi `$. It turns out that all $`n`$-graphs are integral submanifolds, of dimension $`k`$, of a distribution called the *canonical contact system* on $`J^n(^k,^m)`$ or the *Cartan distribution* on $`J^n(^k,^m)`$. The Pfaffian system that anihilates this distribution, which is also called the *canonical contact system* (see e.g and ), is given in the canonical coordinates of $`J^n(^k,^m)`$ by
$$dp_i^\sigma \underset{j=1}{\overset{k}{}}p_i^{\sigma +1_j}dq^j=0,\text{ for }1im\text{ and for }1\left|\sigma \right|n1\text{,}$$
where $`\sigma +1_j=(\sigma _1,\mathrm{},\sigma _j+1,\mathrm{},\sigma _k)`$.
The above description explains the importance of contact systems in geometric theory of (partial) differential equations and in differential geometry. In the former, a (partial) differential equation is interpreted as a submanifold in $`J^n(^k,^m)`$ and thus it is natural to study the geometry of pairs consisting of a contact system and a submanifold, see e.g. , , and . In the latter, contact systems describe, for instance, diffeomorphisms which preserve the $`n`$-graphs of appplications (for example, $`n`$-graphs of curves in the case $`k=1`$), see e.g. and .
A natural problem which arises is to characterize those distributions which are (locally) equivalent to a canonical contact system. This problem was posed by Pfaff in 1815 and seems still to be open in its full generality although many important particular solutions have been obtained. In the case $`n=1`$, with $`m=1`$ and an arbitrary $`k`$ the final solution has been obtained by Darboux in his famous theorem generalizing earlier results of Pfaff and Frobenius . The case $`n=2`$, $`m=1`$ and $`k=1`$ was solved by Engel in . The case $`n2`$, $`m=1`$ and $`k=1`$ was solved by E. von Weber , Cartan and Goursat (at generic points) and by Libermann , Kumpera and Ruiz , and Murray (at an arbitrary point). The case $`n=1`$, with $`k`$ and $`m`$ arbitrary has been studied and solved by Bryant (see also ).
This paper is devoted to the problem of when a given distribution is locally equivalent to the canonical contact system in the case $`k=1`$, $`n`$ and $`m`$ arbitrary, that is to the canonical contact system for curves. This problem has been studied by Gardner and Shadwick , Murray , and Tilbury and Sastry (as a particular case of the the problem of equivalence to the so-called extended Goursat normal form). Their solutions are based on a result of that assures the equivalence provided that a certain differential form satisfies precise congruence relations. The problem of how to verify the existence of such a form had apparently remained open. This difficulty was solved by Aranda-Bricaire and Pomet , who proposed an algorithm which determines the existence of such a form. Their solution, although being elegant and checkable, uses the formalism of infinite dimensional manifolds and thus goes away from classical results characterizing contact systems.
The aim of our paper is two-fold. Firstly, in Theorem 1.1, we give geometric checkable conditions, based on the classical notion of Engel’s rank, which characterize regular contact systems for curves, i.e., for $`k=1`$ and arbitrary $`n`$ and $`m`$. Secondly, we extend our approach to singular points and we prove, in Theorem 3.2, that any singular contact system can be put to a normal form for which we propose the name extended Kumpera-Ruiz normal form. That form generalizes canonical contact systems for curves in a way analogous to that how Kumpera-Ruiz normal form (see ; compare , , , , and ) generalizes Goursat normal form, that is the canonical contact system on $`J^n(,)`$. When checking conditions of Theorem 1.1 we must determine whether a given distribution possesses a corank one involutive subdistribution. An elegant answer given to this problem by Bryant implies that our conditions become checkable.
The paper is organized as follows. In Section 1 we define the canonical system for curves and we give the first main result of the paper, Theorem 1.1, characterizing distributions that are locally equivalent to a canonical contact system for curves. In Section 2 we discuss the problem of whether a given distribution possesses a corank one involutive subdistribution and we recall Bryant’s solution of this problem (see also Appendices A and B). In Section 3 we introduce extended Kumpera-Ruiz normal forms and give the second main result of the paper, Theorem 3.2, stating that any singular contact system is locally equivalent to an extended Kumpera-Ruiz normal form. Finally, Section 4 contains proofs of all our results.
Acknowledgments: The authors are grateful to Robert Bryant, who kindly sent them a copy of his Ph.D. Thesis and to Jean-Baptiste Pomet for discussion on his results .
## 1 The Canonical Contact System for Curves
A rank$`k`$ *distribution* $`𝒟`$ on a smooth manifold $`M`$ is a map that assigns smoothly to each point $`p`$ in$`M`$ a linear subspace $`𝒟(p)T_pM`$ of dimension $`k`$. In other words, a rank $`k`$ distribution is a smooth rank $`k`$ subbundle of the tangent bundle $`TM`$. Such a field of tangent $`k`$-planes is spanned locally by $`k`$ pointwise linearly independent smooth vector fields $`f_1,\mathrm{},f_k`$ on $`M`$, which will be denoted by $`𝒟=(f_1,\mathrm{},f_k)`$. Two distributions $`𝒟`$ and $`\stackrel{~}{𝒟}`$ defined on two manifolds $`M`$ and $`\stackrel{~}{M}`$, respectively, are *equivalent* if there exists a smooth diffeomorphism $`\phi `$ between $`M`$ and $`\stackrel{~}{M}`$ such that $`(\phi _{}𝒟)(\stackrel{~}{p})=\stackrel{~}{𝒟}(\stackrel{~}{p})`$, for each point$`\stackrel{~}{p}`$ in $`\stackrel{~}{M}`$.
The *derived flag* of a distribution $`𝒟`$ is the sequence of modules of vector fields $`𝒟^{(0)}𝒟^{(1)}\mathrm{}`$ defined inductively by
$$𝒟^{(0)}=𝒟\text{ and }𝒟^{(i+1)}=𝒟^{(i)}+[𝒟^{(i)},𝒟^{(i)}]\text{, \hspace{1em}for }i0\text{.}$$
(1)
The *Lie flag* is the sequence of modules of vector fields $`𝒟_0𝒟_1\mathrm{}`$ defined inductively by
$$𝒟_0=𝒟\text{ and }𝒟_{i+1}=𝒟_i+[𝒟_0,𝒟_i]\text{, \hspace{1em}for }i0\text{.}$$
(2)
In general, the derived and Lie flags are different; though for any point $`p`$ in the underlying manifold the inclusion $`𝒟_i(p)𝒟^{(i)}(p)`$ clearly holds, for $`i0`$.
For a given distribution $`𝒟`$, defined on a manifold $`M`$, we will say that a point $`p`$ of $`M`$ is a *regular point* of $`𝒟`$ if all the elements $`𝒟_i`$ of its Lie flag have constant rank in a small enough neighborhood of $`p`$. A distribution $`𝒟`$ is said to be *completely nonholonomic* if, for each point $`p`$ in $`M`$, there exists an integer $`N(p)`$ such that $`𝒟_{N(p)}(p)=T_pM`$. A distribution $`𝒟`$ is said to be *involutive* if its first derived system satisfies $`𝒟^{(1)}=𝒟^{(0)}`$.
An alternative description of the above defined objects can also be given using the dual language of differential forms. A *Pfaffian system* $``$ of rank $`s`$ on a smooth manifold $`M`$ is a map that assigns smoothly to each point $`p`$ in$`M`$ a linear subspace $`(p)T_p^{}M`$ of dimension $`s`$. In other words, a Pfaffian system of rank $`s`$ is a smooth subbundle of rank $`s`$ of the cotangent bundle $`T^{}M`$. Such a field of cotangent $`s`$-planes is spanned locally by $`s`$ pointwise linearly independent smooth differential $`1`$-forms $`\omega _1,\mathrm{},\omega _s`$ on $`M`$, which will be denoted by $`=(\omega _1,\mathrm{},\omega _s)`$. Two Pfaffian systems $``$ and $`\stackrel{~}{}`$ defined on two manifolds $`M`$ and $`\stackrel{~}{M}`$, respectively, are *equivalent* if there exists a smooth diffeomorphism $`\phi `$ between $`M`$ and $`\stackrel{~}{M}`$ such that $`(p)=(\phi ^{}\stackrel{~}{})(p)`$, for each point$`p`$ in $`M`$.
For a Pfaffian system $``$, we can define its *derived flag* $`^{(0)}^{(1)}\mathrm{}`$ by the relations $`^{(0)}=`$ and $`^{(i+1)}=\{\alpha ^{(i)}:d\alpha 0mod^{(i)}\}`$, for $`i0`$, provided that each element $`^{(i)}`$ of this sequence has constant rank. In this case, it is immediate to see that the derived flag of the distribution $`𝒟=^{}`$ coincides with the sequence of distributions that anihilate the elements of the derived flag of $``$, that is
$$𝒟^{(i)}=(^{(i)})^{},\text{ for }i0.$$
For a given Pfaffian system $``$, we will say that a point $`p`$ of $`M`$ is a *regular point* if $`p`$ is a regular point for the distribution $`𝒟=^{}`$, that is if all elements $`𝒟_i`$ of the Lie flag are of constant rank in a small enough neighborhood of $`p`$.
Consider the space $`J^n(,^m)`$ of jets of order $`n1`$ of functions from $``$ into $`^m`$. This space is diffeomorphic to $`^{(n+1)m+1}`$. The canonical coordinates associated to $``$ (denoted by $`x_0^0`$) and to $`^m`$ (denoted by $`x_1^0,\mathrm{},x_m^0`$) can be used to define the *canonical coordinates* on $`J^n(,^m)`$, which will be denoted by
$$x_0^0,x_1^0,\mathrm{},x_m^0,x_1^1,\mathrm{},x_m^1,\mathrm{},x_1^n,\mathrm{},x_m^n,$$
with obvious indentifications $`x_0^0=q`$ and $`x_i^0=u_i`$, for $`1im`$, and $`x_i^j=p_i^j`$, for $`1im`$ and $`1jn`$ (see the beginning of the Introduction). Observe that any smooth map $`\phi `$ from $``$ into $`^m`$ defines a curve in $`J^n(,^m)`$ by the relations $`x_j^i=\phi _j^{(i)}(x_0^0)`$, for $`0in`$ and $`1jm`$, where $`\phi _j^{(i)}`$ denotes the $`i`$-th derivative with respect to $`x_0^0`$ of the $`j`$-th component of $`\phi `$. This curve is called the $`n`$-*graph* of $`\phi `$. It is clear that not all curves in $`J^n(,^m)`$ are $`n`$-graphs of maps. In order to distinguish the “good” curves from the “bad” ones, we should introduce a set of constraints on the velocities of curves in $`J^n(,^m)`$. In other words, we should endow $`J^n(,^m)`$ with a nonholonomic structure.
The *canonical contact system* on $`J^n(,^m)`$ is the completely nonholonomic distribution spanned by the following family of vector fields:
$$(\begin{array}{c}\frac{}{x_1^n}\hfill \end{array},\mathrm{},\begin{array}{c}\frac{}{x_m^n}\hfill \end{array},\begin{array}{c}\underset{i=0}{\overset{n1}{}}\underset{j=1}{\overset{m}{}}x_j^{i+1}\frac{}{x_j^i}+\frac{}{x_0^0}\hfill \end{array}).$$
(3)
By definition, if a curve in $`J^n(,^m)`$ is the $`n`$-graph of some map then it is an integral curve of the canonical contact system. More precisely, a section $`\sigma :J^n(,^m)`$ is the $`n`$-graph of a curve $`\phi :^n`$ if and only if it is an integral curve of the canonical contact system on $`J^n(,^m)`$ (see e.g. ).
The aim of our paper is to give a complete answer to the question “Which distributions are locally equivalent to the canonical contact system on $`J^n(,^m)`$?”. The following result will be the starting point of our study.
###### Theorem 1.1 (contact systems for curves)
A rank $`m+1`$ distribution $`𝒟`$ on a manifold $`M`$ of dimension $`(n+1)m+1`$ is equivalent, in a small enough neighborhood of any point $`p`$ in $`M`$, to the canonical contact system on $`J^n(,^m)`$ if and only if the two following conditions hold, for $`0in`$.
1. Each element $`𝒟^{(i)}`$ of the derived flag has constant rank $`(i+1)m+1`$ and contains an involutive subdistribution $`_i𝒟^{(i)}`$ that has constant corank one in $`𝒟^{(i)}`$.
2. Each element $`𝒟_i`$ of the Lie flag has constant rank $`(i+1)m+1`$.
This result yields a constructive test for the local equivalence to the canonical contact system for curves, provided that we know how to check whether or not a given distribution admits a corank one involutive subdistribution. We give in the next section a checkable necessary and sufficient condition for the existence of such a distribution. The proof of Theorem 1.1 will be given in Section 4.
## 2 Corank One Involutive Subdistributions
The aim of this Section is to give an answer to the following question: “When does a given constant rank distribution $`𝒟`$ contain an involutive subdistribution $`𝒟`$ that has constant corank one in $`𝒟`$?”. In fact, the answer to this question is an immediate consequence of a result contained in Bryant’s Ph.D. thesis . Links between Bryant’s result and the characterization of the canonical contact system for curves have also been observed by Aranda-Bricaire and Pomet .
A *characteristic vector field* of a distribution $`𝒟`$ is a vector field $`f`$ that belongs to $`𝒟`$ and satisfies $`[f,𝒟]𝒟`$. The *characteristic distribution* of $`𝒟`$, which will be denoted by $`𝒞`$, is the module spanned by all its characteristic vector fields. It follows directly from the Jacobi identity that the characteristic distribution is always involutive. For a constant rank Pfaffian system $``$, the characteristic distribution of $`^{}`$ is often called the *Cartan system* of $``$. We refer the reader to for a definition of the Cartan system given in the language of Pfaffian systems.
The *Engel rank* of a Pfaffian system $``$, at a point $`p`$, is the largest integer $`\rho `$ such that there exists a $`1`$-form $`\omega `$ in $``$ for which we have $`(d\omega )^\rho (p)0mod`$. The Engel rank of a constant rank distribution $`𝒟`$ will be, by definition, the Engel rank of its anihilator$`𝒟^{}`$. Obviously, the Engel rank $`\rho `$ of a distribution equals zero at each point if and only if the distribution is involutive.
We will now give an equivalent definition of the Engel rank in the language of vector fields, in the particular case when $`\rho =1`$, which will be important in the paper. Let $`𝒟`$ be a distribution such that $`𝒟^{(0)}`$ and $`𝒟^{(1)}`$ have constant ranks $`d_0`$ and $`d_1`$, respectively, and denote $`r_0(p)=d_1(p)d_0(p)`$. Assume that $`d_02`$ and $`r_01`$. Take a family of vector fields
$$(f_1,\mathrm{},f_{d_0},g_1,\mathrm{},g_{r_0})$$
such that $`𝒟^{(0)}=(f_1,\mathrm{},f_{d_0})`$ and$`𝒟^{(1)}=(f_1,\mathrm{},f_{d_0},g_1,\mathrm{},g_{r_0})`$. The *structure functions* $`c_{ij}^k`$ associated to those generators, for $`1i<jd_0`$ and $`1kr_0`$, are the smooth functions defined by the following relations:
$$[f_i,f_j]=\underset{k=1}{\overset{r_0}{}}c_{ij}^kg_kmod𝒟^{(0)},\text{ for }1i<jd_0.$$
It is important to point out that the structure functions are not invariantly related to the distribution $`𝒟`$, since they depend on the choice of generators.
Assume that the Engel rank $`\rho `$ of $`𝒟`$ is constant and that $`r_01`$. It is easy to check that $`\rho =1`$ if and only if either $`d_0=2`$, or $`d_0=3`$, or $`d_04`$ and the structure functions satisfy the relations
$$c_{ij}^pc_{kl}^qc_{ik}^pc_{jl}^q+c_{il}^pc_{jk}^q+c_{jk}^pc_{il}^qc_{jl}^pc_{ik}^q+c_{kl}^pc_{ij}^q=0,$$
(4)
for each sextuple $`(i,j,k,l,p,q)`$ of integers such that $`1i<j<k<ld_0`$ and $`1pr_0`$ and $`1qr_0`$.
The following result is a particular case of Bryant’s algebraic Lemma (see also ), which is the cornerstone of Bryant’s characterization of the canonical contact system on $`J^1(^k,^m)`$.
###### Lemma 2.1 (Bryant)
Let $`𝒟`$ be a distribution such that $`𝒟^{(0)}`$ and $`𝒟^{(1)}`$ have constant ranks $`d_0`$ and $`d_1`$, respectively. Assume that $`r_01`$. Then the two following conditions are equivalent:
1. The characteristic distribution $`𝒞`$ of $`𝒟`$ has constant rank $`c_0=d_0r_01`$ and the Engel rank $`\rho `$ of $`𝒟`$ is constant and equals$`1`$;
2. The distribution $`𝒟`$ contains a subdistribution $`𝒟`$ that has constant corank one in $`𝒟`$ and satisfies $`[,]𝒟`$.
Observe that if the first condition is satisfied then we must necessarily have $`r_0d_01`$. The following result is included in the proof of Bryant’s normal form Theorem (see also ). In Appendix A we give an alternative proof of its surprising Item (iii) ; our proof explains the role of the assumption $`r_03`$ by relating it to the Jacobi identity.
###### Lemma 2.2 (Bryant)
Let $`𝒟`$ be a distribution such that $`𝒟^{(0)}`$ and $`𝒟^{(1)}`$ have constant ranks $`d_0`$ and $`d_1`$, respectively. Assume that the distribution $`𝒟`$ contains a subdistribution $`𝒟`$ that has constant corank one in $`𝒟`$ and satisfies $`[,]𝒟`$.
1. If $`r_0=1`$ then the distribution $`𝒟`$ contains an involutive subdistribution $`𝒟`$ that has constant corank one in $`𝒟`$;
2. If $`r_02`$ then $``$ is unique;
3. If $`r_03`$ then $``$ is involutive.
Observe that, in the first item of the above Lemma, the involutive subdistribution $``$ can be different from $``$, which is not necessarily involutive. The following result is a direct consequence of Bryant’s work. In order to avoid the trivial case $`r_0=0`$, for which the existence of a corank one involutive subdistribution is obvious, we will assume that $`r_01`$.
###### Corollary 2.3 (corank one involutive subdistributions)
Let $`𝒟`$ be a distribution such that $`𝒟^{(0)}`$ and $`𝒟^{(1)}`$ have constant ranks $`d_0`$ and $`d_1`$, respectively. Assume that $`r_01`$. Then, the distribution $`𝒟`$ contains an involutive subdistribution $`𝒟`$ that has constant corank one in$`𝒟`$ if and only if the three following conditions hold:
1. The characteristic distribution $`𝒞`$ of $`𝒟`$ has constant rank $`c_0=d_0r_01`$;
2. The Engel rank $`\rho `$ of $`𝒟`$ is constant and equals$`1`$;
3. If $`r_0=2`$ then, additionally, the unique corank one subdistribution $`𝒟`$ such that $`[,]𝒟`$ must be involutive.
We would like to to emphasize that the above conditions are easy to verify, as well as the conditions of Corollary 2.4 below. Indeed, for any distribution, or the corresponding Pfaffian system, we can compute the characteristic distribution $`𝒞`$ and check whether or not the Engel rank equals $`1`$ using, respectively, the formula (12) and the condition (13) of Appendix B. This gives the solution if $`r_02`$. If $`r_0=2`$ we have additionally to check the involutivness of the unique distribution $``$ satisfying $`[,]𝒟`$, whose explicit construction is also given in Appendix B.
Combining Theorem 1.1 and Corollary 2.3 we get the following characterization of the canonical system on $`J^n(,^m)`$, using the language of Pfaffian systems.
###### Corollary 2.4 (contact systems for curves)
Let $``$ be a Pfaffian system of rank $`nm`$, defined on a manifold $`M`$of dimension $`(n+1)m+1`$. If $`m2`$, the Pfaffian system $``$ is locally equivalent, at a given point $`p`$ of $`M`$, to the canonical contact system on $`J^n(,^m)`$ if and only if
1. The rank of each derived system $`^{(i)}`$ is constant and equals $`(ni)m`$, for $`0in`$;
2. The Engel rank of $`^{(i)}`$ is constant and equals $`1`$, for $`0in`$;
3. The rank of each Cartan system $`𝒞(^{(i)})`$ is constant and equals $`(n+1i)m+1`$, for $`0in1`$ ;
4. The point $`p`$ is a regular point for $``$.
In other words, if $`m2`$, the characterization of the canonical contact system on $`J^n(,^m)`$ turns out to be a natural combination of that given for $`J^1(,^m)`$ by Bryant (see and ) and that given for $`J^n(,)`$ by Murray .
## 3 Extended Kumpera-Ruiz Normal Forms
The aim of this Section is to study the class of distributions that satisfy condition (i) of Theorem 1.1 but fail to satisfy condition (ii) of that theorem. The fist condition describes the geometry of the canonical contact system while the second condition characterizes regular points. In this sense, systems that satisfy the former but fail to satisfy the latter can be considered as “singular” contact systems for curves. We will show that any such distribution can be brought to a normal form for which we propose the name extended Kumpera-Ruiz normal form. Those forms generalize the canonical contact system on $`J^1(,^m)`$ in a way analogous to that how Kumpera-Ruiz normal forms (see e.g. , , , and ) generalize the canonical system on $`J^1(,)`$, which is also called Goursat normal form.
Consider the family of vector fields $`\kappa ^1=(\kappa _1^1,\mathrm{},\kappa _m^1,\kappa _0^1)`$ that span the canonical contact system on $`J^1(,^m)`$, where
$`\kappa _1^1`$ $`=\frac{}{x_1^1},\mathrm{},\kappa _m^1=\frac{}{x_m^1}`$
$`\kappa _0^1`$ $`=x_1^1\frac{}{x_1^0}+\mathrm{}+x_m^1\frac{}{x_m^0}+\frac{}{x_0^0},`$
and the the family of vector fields $`\kappa ^2=(\kappa _1^2,\mathrm{},\kappa _m^2,\kappa _0^2)`$ that spans the canonical contact system on $`J^2(,^m)`$, where
$`\kappa _1^2`$ $`=\frac{}{x_1^2},\mathrm{},\kappa _m^2=\frac{}{x_m^2}`$
$`\kappa _0^2`$ $`=x_1^2\frac{}{x_1^1}+\mathrm{}+x_m^2\frac{}{x_m^1}+x_1^1\frac{}{x_1^0}+\mathrm{}+x_m^1\frac{}{x_m^0}+\frac{}{x_0^0}.`$
Loosely speaking, we can write
$`\kappa _1^2`$ $`=\frac{}{x_1^2},\mathrm{},\kappa _m^2=\frac{}{x_m^2}`$
$`\kappa _0^2`$ $`=x_1^2\kappa _1^1+\mathrm{}+x_m^2\kappa _m^1+\kappa _0^1.`$
In order to make this precise we will adopt the following natural notation. Consider an arbitrary vector field $`f`$ given on $`J^{n1}(,^m)`$ by
$$f=\underset{i=0}{\overset{n1}{}}\underset{j=1}{\overset{m}{}}f_j^i(\overline{x}^{n1})\frac{}{x_j^i}+f_0^0(\overline{x}^{n1})\frac{}{x_0^0},$$
where $`\overline{x}^{n1}`$ denotes the coordinates $`x_0^0,x_1^0,\mathrm{},x_m^0,x_1^1,\mathrm{},x_m^1,\mathrm{},x_1^{n1},\mathrm{},x_m^{n1}`$ of $`J^{n1}(,^m)`$. We can *lift* the vector field $`f`$ to a vector field on $`J^n(,^m)`$, which we also denote by $`f`$, by taking
$$f=\underset{i=0}{\overset{n1}{}}\underset{j=1}{\overset{m}{}}f_j^i(\overline{x}^{n1})\frac{}{x_j^i}+f_0^0(\overline{x}^{n1})\frac{}{x_0^0}+0\frac{}{x_1^n}+\mathrm{}+0\frac{}{x_m^n}.$$
That is, we lift $`f`$ by translating it along the directions $`\frac{}{x_1^n},\mathrm{},\frac{}{x_m^n}`$.
###### Notation 3.1 (lifts of vector fields)
From now on, in any expression of the form $`\kappa _0^n=_{i=0}^m\alpha _i(x)\kappa _i^{n1}`$, the vector fields $`\kappa _0^{n1},\mathrm{},\kappa _m^{n1}`$ should be considered as their above defined lifts.
Let $`\kappa ^{n1}=(\kappa _1^{n1},\mathrm{},\kappa _m^{n1},\kappa _0^{n1})`$ denote a family of vector fields defined on $`J^{n1}(,^m)`$. A *regular prolongation*, with a parameter $`c^n`$, of $`\kappa ^{n1}`$, denoted by $`\kappa ^n=R_{c^n}(\kappa ^{n1})`$, is a family of vector fields $`\kappa ^n=(\kappa _1^n,\mathrm{},\kappa _m^n,\kappa _0^n)`$ defined on $`J^n(,^m)`$ by
$$\begin{array}{c}\kappa _1^n=\frac{}{x_1^n},\mathrm{},\kappa _m^n=\frac{}{x_m^n}\hfill \\ \kappa _0^n=(x_1^n+c_1^n)\kappa _1^{n1}+\mathrm{}+(x_m^n+c_m^n)\kappa _m^{n1}+\kappa _0^{n1},\hfill \end{array}$$
(5)
where $`c^n=(c_1^n,\mathrm{},c_m^n)`$ is a vector of$`m`$ real constants. A *singular prolongation*, with a parameter $`c^n`$, of $`\kappa ^{n1}`$, denoted by $`\kappa ^n=S_{c_n}(\kappa ^{n1})`$, is a family of vector fields $`\kappa ^n=(\kappa _1^n,\mathrm{},\kappa _m^n,\kappa _0^n)`$ defined on $`J^n(,^m)`$ by
$$\begin{array}{c}\kappa _1^n=\frac{}{x_1^n},\mathrm{},\kappa _m^n=\frac{}{x_m^n}\hfill \\ \kappa _0^n=(x_1^n+c_1^n)\kappa _1^{n1}+\mathrm{}+(x_{m1}^n+c_{m1}^n)\kappa _{m1}^{n1}+\kappa _m^{n1}+x_m^n\kappa _0^{n1},\hfill \end{array}$$
(6)
where $`c^n=(c_1^n,\mathrm{},c_{m1}^n,0)`$ is a vector of$`m`$ real constants, the last one being zero.
A family of vector fields $`\kappa ^n`$ on $`J^n(,^m)`$, for $`n1`$, will be called an *extended Kumpera-Ruiz normal form* if $`\kappa ^n=\sigma _n\mathrm{}\sigma _2(\kappa ^1)`$, where for each $`2in`$ the map $`\sigma _i`$ equals either $`R_{c^i}`$ or $`S_{c^i}`$, for some vector parameters$`c^i`$. In other words, a Kumpera-Ruiz normal form is a family of vector fields obtained by successive prolongations from the family of vector fields that spans the canonical contact system on $`J^1(,^m)`$.
The above defined prolongations and prolongations-based definition of extended Kumpera-Ruiz normal forms generalizes for contact systems analogous operations introduced by the authors for Goursat structures.
Let $`x:M^{(n+1)m+1}J^n(,^m)`$ be a local coordinate system on a manifold $`M`$, in a neighborhood of a given point $`p`$ in $`M`$. We will say that an extended Kumpera-Ruiz normal form on $`J^n(,^m)`$, defined in $`x`$-coordinates, is *centered* at $`p`$ if we have $`x(p)=0`$. For example, on $`J^2(,^2)`$, we have the two following extended Kumpera-Ruiz normal forms
$`(\begin{array}{c}\frac{}{x_1^2}\end{array},\begin{array}{c}\frac{}{x_2^2}\end{array},\begin{array}{c}x_1^2\frac{}{x_1^1}+x_2^2\frac{}{x_2^1}+x_1^1\frac{}{x_1^0}+x_2^1\frac{}{x_2^0}+\frac{}{x_0^0}\end{array})`$
$`(\begin{array}{c}\frac{}{x_1^2}\end{array},\begin{array}{c}\frac{}{x_2^2}\end{array},\begin{array}{c}x_1^2\frac{}{x_1^1}+\frac{}{x_2^1}+x_2^2\left(x_1^1\frac{}{x_1^0}+x_2^1\frac{}{x_2^0}+\frac{}{x_0^0}\right)\end{array}),`$
defined by $`R_{(0,0)}(\kappa ^1)`$ and $`S_{(0,0)}(\kappa ^1)`$, respectively. These two normal forms are obviously centered at zero.
The following theorem is the second main contribution of the paper. It asserts that extended Kumpera-Ruiz normal forms serve as local normal forms for all singular contact systems for curves, that is for all distributions that satisfy condition (i) of Theorem 1.1 but fail to fulfill the regularity condition (ii) of that theorem.
###### Theorem 3.2 (extended Kumpera-Ruiz normal forms)
A distribution $`𝒟`$ of rank $`m+1`$ on a manifold$`M`$ of dimension $`(n+1)m+1`$ is equivalent, in a small enough neighborhood of any point$`p`$ in$`M`$, to a distribution spanned by an extended Kumpera-Ruiz normal form, centered at$`p`$ and defined on a suitably chosen neighborhood of zero, if and only if each element$`𝒟^{(i)}`$ of its derived flag has constant rank $`(i+1)m+1`$ and contains an involutive subdistribution$`_i𝒟^{(i)}`$ that has constant corank one in$`𝒟^{(i)}`$, for $`0in`$.
## 4 Proof of the Two Main Theorems
In this Section we provide a proof of Theorem 3.2. We start with several Lemmas — which will be used in the proof — that describe the geometry of incidence between characteristic distributions $`𝒞_i`$ and involutive corank one subdistributions $`_i`$ of $`𝒟^{(i)}`$. Then we prove Theorem 3.2. Finally, we conclude Theorem 1.1 as a corollary of Theorem 3.2.
For $`i0`$, we will denote by $`𝒞_i`$ the characteristic distribution of $`𝒟^{(i)}`$. It follows directly from the Jacobi identity that $`𝒞_i𝒞_{i+1}`$. Define $`d_i(p)=dim𝒟^{(i)}(p)`$ and $`c_i(p)=dim𝒞_i(p)`$. Moreover, denote $`r_i(p)=d_{i+1}(p)d_i(p)`$.
Though the following result is a direct consequence of Bryant’s algebraic Lemma, we will give its proof as a warm up exercice. Indeed, the method used to prove the inclusion $`𝒞_0_0`$ is also used in the proofs of inclusions $`_0_1`$ and $`_0𝒞_1`$, which will be considered later.
###### Lemma 4.1 ($`𝒞_0_0`$)
Let $`𝒟`$ be a distribution such that $`𝒟^{(0)}`$ and $`𝒟^{(1)}`$ have constant ranks $`d_0`$ and $`d_1d_0+1`$, respectively. If the distribution $`𝒟`$ contains an involutive subdistribution $`_0𝒟^{(0)}`$ that has constant corank one in $`𝒟^{(0)}`$ then the ranks of $`𝒟^{(0)}`$ and $`𝒟^{(1)}`$ satisfy $`r_0d_01`$. Moreover:
1. The characteristic distribution $`𝒞_0`$ satisfies $`𝒞_0_0`$;
2. The rank of $`𝒞_0`$ is constant and equal to $`d_0r_01`$.
Proof: Assume that $`𝒟`$ contains an involutive subdistribution $`_0𝒟^{(0)}`$ of constant corank one. The relation $`r_0d_01`$ is obvious. In order to prove by contradiction that $`𝒞_0_0`$, assume that for some point $`p`$ the vector space $`𝒞_0(p)`$ is not contained in $`_0(p)`$. Then, in a small enough neighborhood of$`p`$, we can assume that the distribution $`𝒟^{(0)}`$ is a direct sum $`𝒟^{(0)}=(h)_0`$, where $`h`$ is a vector field that belongs to $`𝒞_0`$ but that does not belong to $`_0`$. Since $`[_0,_0]𝒟^{(0)}`$ we have $`𝒟^{(1)}=𝒟^{(0)}+[h,_0]`$. But since $`[h,𝒟^{(0)}]𝒟^{(0)}`$ we have $`𝒟^{(1)}=𝒟^{(0)}`$, which is impossible because $`r_01`$.
Now, let us compute the rank of $`𝒞_0`$. Since the corank of $`_0`$ in $`𝒟^{(0)}`$ equals $`1`$, we have a local decomposition $`𝒟^{(0)}=(f_{d_0})_0`$, where $`f_{d_0}`$ is an arbitrary vector field that belongs to $`𝒟^{(0)}`$ but that does not belong to $`_0`$. Since $`𝒟^{(1)}=𝒟^{(0)}+[f_{d_0},_0]`$ we can find, locally, a family of vector fields
$$(f_1,\mathrm{},f_{r_0},f_{r_0+1},\mathrm{},f_{d_01})$$
that span $`_0`$ and satisfies
$$𝒟^{(1)}=𝒟^{(0)}([f_{d_0},f_1],\mathrm{},[f_{d_0},f_{r_0}]).$$
It follows that, for $`r_0+1id_01`$, we can find some smooth functions $`\alpha _{ij}`$ such that $`[f_{d_0},f_i]=_{j=1}^{r_0}\alpha _{ij}[f_{d_0},f_j]mod𝒟^{(0)}`$. On the one hand, we have $`dim𝒞_0(p)d_0r_01`$, at each point $`p`$, because the vector fields $`h_i=f_i_{j=1}^{r_0}\alpha _{ij}f_j`$, for $`r_0+1id_01`$, satisfy $`[h_i,𝒟^{(0)}]𝒟^{(0)}`$ and are pointwise linearly independent. But, on the other hand, we have $`dim𝒞_0(p)d_0r_01`$, at any point $`p`$, because otherwise we would have $`dim𝒟^{(1)}(p)d_11`$. Hence $`dim𝒞_0(p)=d_0r_01`$, for each point $`p`$ in the underlying manifold. $`\mathrm{}`$
###### Lemma 4.2 ($`_0_1`$)
Let $`𝒟`$ be a distribution such that $`𝒟^{(0)}`$, $`𝒟^{(1)}`$, and $`𝒟^{(2)}`$ have constant ranks $`d_0`$, $`d_1d_0+2`$, and $`d_2d_1+2`$, respectively. Assume that each distribution $`𝒟^{(i)}`$, for $`i=0`$ and$`1`$, contains an involutive subdistribution $`_i𝒟^{(i)}`$ that has constant corank one in $`𝒟^{(i)}`$. Then $`_0_1`$.
Proof: Assume that there exists a point $`p`$ such that the vector space$`_0(p)`$ is not contained in $`_1(p)`$. We will show that this assumption leads to $`d_2d_1+1`$. To start with, observe that on the one hand we have $`dim𝒟^{(0)}(q)_1(q)d_01`$, for any point $`q`$ in a small enough neighborhood of$`p`$, because $`_0(p)`$ is not contained in$`_1(p)`$; but that on the other hand we have
$`dim𝒟^{(1)}(q)`$ $`dim\left(𝒟^{(0)}(q)_1(q)\right)`$
$`=dim𝒟^{(0)}(q)+dim_1(q)dim\left(𝒟^{(0)}(q)_1(q)\right),`$
for any point $`q`$, which implies that $`dim\left(𝒟^{(0)}(q)_1(q)\right)d_01`$. Therefore, locally around $`p`$, we have $`dim𝒟^{(0)}(q)_1(q)=d_01`$. An analogous argument can be applied to the intersection $`_0_1`$ in order to show that $`dim\left(_0(q)_1(q)\right)=d_02`$.
The above relations between the ranks of $`𝒟^{(0)}`$, $`𝒟^{(0)}_1`$, and $`_0_1`$ imply that we can find a local basis $`(f_1,\mathrm{},f_{d_0})`$ of $`𝒟^{(0)}`$ such that $`_0=(f_1,\mathrm{},f_{d_01})`$ and $`_0_1=(f_1,\mathrm{},f_{d_02})`$ and $`𝒟^{(0)}_1=(f_1,\mathrm{},f_{d_02},f_{d_0})`$. Moreover, we can assume that $`𝒞_0=(f_1,\mathrm{},f_{c_0})`$, where $`c_0<d_01`$. Indeed, by Lemma 4.1, we have $`𝒞_0_0`$ and $`𝒞_0𝒞_1_1`$, which obviously implies $`𝒞_0_0_1`$.
Since the vector field $`f_{d_0}`$ does not belong to $`_0`$, we have $`𝒟^{(1)}=𝒟^{(0)}+[f_{d_0},_0]`$. Therefore, the distribution $`𝒟^{(1)}`$ admits the following local decomposition
$$𝒟^{(1)}=𝒟^{(0)}([f_{d_0},f_{c_0+1}],\mathrm{},[f_{d_0},f_{d_01}]).$$
Denote $`g_i=[f_{d_0},f_i]`$, for $`c_0+1id_01`$. Since the vector field $`f_{d_01}`$ does not belong to $`_1`$, we have $`𝒟^{(2)}=𝒟^{(1)}+[f_{d_01},𝒟^{(1)}]`$. Therefore, all vector fields in $`𝒟^{(2)}`$ are linear combinations of those belonging to $`𝒟^{(1)}`$ and of the vector fields $`[f_{d_01},g_i]`$, for $`c_0+1id_01`$.
Now, observe that $`_0_1𝒞_1`$. Indeed, the distribution $`𝒟^{(1)}`$ admits a local decomposition $`𝒟^{(1)}=(f_{d_01})_1`$. But $`[_0_1,_1]_1`$ and $`[_0_1,f_{d_01}]_0`$. Hence $`[_0_1,𝒟^{(1)}]𝒟^{(1)}`$.
We claim that each vector field $`[f_{d_01},g_i]`$, for $`c_0+1id_02`$, belongs to $`𝒟^{(1)}`$. Indeed, the Jacobi identity gives
$$[f_{d_01},g_i]+[f_{d_0},[f_i,f_{d_01}]]+[f_i,[f_{d_01},f_{d_0}]]=0.$$
On the one hand, the vector field $`[f_i,[f_{d_01},f_{d_0}]]`$ belongs to $`𝒟^{(1)}`$ because the vector field $`f_i`$ belongs to $`_0_1`$, which is contained in $`𝒞_1`$. But on the other hand, the vector field $`[f_{d_0},[f_i,f_{d_01}]]`$also belongs to $`𝒟^{(1)}`$ because $`[f_i,f_{d_01}]`$ belongs to $`_0`$. It follows that the vector field $`[f_{d_01},g_i]`$ belongs to $`𝒟^{(1)}`$. Hence $`𝒟^{(2)}=𝒟^{(1)}+[f_{d_01},g_{d_01}]`$, which obviously implies $`d_2d_1+1`$. It follows that we must have $`_0_1`$. $`\mathrm{}`$
###### Lemma 4.3 ($`_0𝒞_1`$)
Let $`𝒟`$ be a distribution such that $`𝒟^{(0)}`$, $`𝒟^{(1)}`$, and $`𝒟^{(2)}`$ have constant ranks $`d_0`$, $`d_1d_0+2`$, and $`d_2d_1+2`$, respectively. Assume that each distribution $`𝒟^{(i)}`$, for $`i=0`$ and$`1`$, contains an involutive subdistribution $`_i𝒟^{(i)}`$ that has constant corank one in $`𝒟^{(i)}`$. Then $`_0𝒞_1`$.
Proof: Take local generators $`(f_1,\mathrm{},f_{d_01},f_{d_0})`$ of $`𝒟^{(0)}`$ such that
$$_0=(f_1,\mathrm{},f_{d_01}).$$
Since by Lemma 4.2 we have $`_0_1`$, each vector field $`f_i`$, for $`1id_01`$, belongs to $`_1`$. Now observe that, since the corank of $`_1`$ in $`𝒟^{(1)}`$ equals $`1`$, we can take $`r_0`$ vector fields $`g_1,\mathrm{},g_{r_0}`$ in $`𝒟^{(1)}`$ that are linearly independent $`mod𝒟^{(0)}`$ and such that each vector field $`g_i`$, for $`1ir_01`$ belongs to $`_1`$. It follows that there exist two smooth functions $`\alpha `$ and $`\beta `$ such that
$$_1=(f_1,\mathrm{},f_{d_01},g_1,\mathrm{},g_{r_01},\alpha f_{d_0}+\beta g_{r_0}).$$
We want to prove that $`_0𝒞_1`$, that is that $`[f_i,𝒟^{(1)}]𝒟^{(1)}`$, for $`1id_01`$. By the definition of $`𝒟^{(1)}`$ we have $`[f_i,𝒟^{(0)}]𝒟^{(1)}`$, for $`1id_01`$, and by the involutivity of $`_1`$ we have $`[f_i,g_j]𝒟^{(1)}`$, for $`1jr_01`$. Therefore, what remains to prove is that $`[_0,g_{r_0}]𝒟^{(1)}`$.
We will prove that $`[_0,g_{r_0}]𝒟^{(1)}`$ by contradiction. Assume that, for some $`1id_01`$, there exists a point $`p`$ such that $`[f_i,g_{r_0}](p)𝒟^{(1)}(p)`$. This implies that $`[f_i,g_{r_0}](q)𝒟^{(1)}(q)`$, for each point $`q`$ in a small neighborhood $`U`$ of$`p`$. But, since $`_1`$ is involutive, the vector field $`[f_i,\alpha f_{d_0}+\beta g_{r_0}]`$ belongs to $`𝒟^{(1)}`$, which clearly implies that $`\beta [f_i,g_{r_0}]`$ also belongs to $`𝒟^{(1)}`$. Therefore, we must have $`\beta (q)=0`$, for each point $`q`$ in $`U`$. It follows that, in a small enough neighborhood of $`p`$, we have $`𝒟^{(0)}_1`$, which implies that $`𝒟^{(1)}_1`$ because $`_1`$ is involutive. Since $`_1`$ has corank one in $`𝒟^{(1)}`$, this is impossible. Therefore $`[_0,g_{r_0}]𝒟^{(1)}`$, which implies that $`_0𝒞_1`$. $`\mathrm{}`$
###### Lemma 4.4 (canonical distribution)
Let $`𝒟`$ be a distribution such that $`𝒟^{(0)}`$, $`𝒟^{(1)}`$, and $`𝒟^{(2)}`$ have constant ranks $`d_0=m+1`$, $`d_1=2m+1`$, and $`d_2=3m+1`$, respectively. If $`m2`$ then assume, additionally, that each distribution $`𝒟^{(i)}`$, for $`i=0`$ and $`1`$, contains an involutive subdistribution $`_i𝒟^{(i)}`$ that has constant corank one in $`𝒟^{(i)}`$. Under these assumptions, we have $`_0=𝒞_1`$, that is the distribution $`𝒟^{(0)}`$ contains a unique involutive subdistribution $`_0`$ that has constant corank one in $`𝒟^{(0)}`$ and satisfies $`[_0,𝒟^{(1)}]𝒟^{(1)}`$.
Proof: For $`m=1`$ the result is well known (see e.g. , , , and ; see also ). Now, if $`m2`$ then Item (ii) of Lemma 4.1 (applied to $`𝒞_1`$) implies that $`dim𝒞_1(p)=2(2m+1)(3m+1)1=m`$ and Lemma 4.3 that $`_0(p)𝒞_1(p)`$, for each point $`p`$ in the underlying manifold. But $`dim_0(p)=m`$. Thus $`_0(p)=𝒞_1(p)`$, which implies that $`_0`$ is uniquely characterized by $`[_0,𝒟^{(1)}]𝒟^{(1)}`$. $`\mathrm{}`$
The following result is a natural generalization of a theorem used by E. von Weber in his study of Goursat structures. In fact, the main idea we will use in our proof of Theorem 3.2 is quite close to Weber’s original idea. A good introduction to the work of E. von Weber is Cartan’s paper . In our own paper , the two main results of are given in a more modern language.
###### Lemma 4.5 (extended Weber normal form)
Let $`𝒟`$ be a distribution defined on a manifold$`M`$ of dimension $`(n+1)m+1`$. Assume that$`𝒟^{(0)}`$ and$`𝒟^{(1)}`$ have constant ranks $`m+1`$ and $`2m+1`$, respectively, and that $`𝒟^{(0)}`$ contains an involutive subdistribution $`_0𝒟^{(0)}`$ that has constant corank one in $`𝒟^{(0)}`$ and satisfies $`_0𝒞_1`$. Then, in a small enough neighborhood of any point$`p`$ in$`M`$, the distribution$`𝒟`$ is equivalent to a distribution spanned on $`J^n(,^m)`$ by a family of vector fields that has the following form:
$$(\begin{array}{c}\frac{}{y_1^n}\end{array},\mathrm{},\begin{array}{c}\frac{}{y_m^n}\end{array},\begin{array}{c}y_1^n\zeta _1^{n1}+\mathrm{}+y_m^n\zeta _m^{n1}+\zeta _0^{n1}\end{array}),$$
where $`_0=(\frac{}{y_1^n},\mathrm{},\frac{}{y_m^n})`$ and the vector fields $`\zeta _1^{n1},\mathrm{},\zeta _m^{n1},\zeta _0^{n1}`$ are lifts of vector fields on $`J^{n1}(,^m)`$, that is,
$$\zeta _i^{n1}=\zeta _i^{n1}(y_0^0,y_1^0,\mathrm{},y_m^0,\mathrm{},y_1^{n1},\mathrm{},y_m^{n1})\text{, for }0im\text{.}$$
Moreover, the set of local coordinates $`(y_0^0,y_1^0,\mathrm{},y_m^0,\mathrm{},y_1^n,\mathrm{},y_m^n)`$, from $`M`$ into $`J^n(,^m)`$ can be taken to be centered at$`p`$.
Proof: It follows directly from Frobenius’ theorem, applied to the distribution $`_0`$, that the distribution$`𝒟`$ is locally equivalent to a distribution spanned on$`^{(n+1)m+1}`$ by a family of vector fields that has the following form:
$$(\begin{array}{c}\frac{}{z_1^n}\end{array},\mathrm{},\begin{array}{c}\frac{}{z_m^n}\end{array},\begin{array}{c}\underset{i=0}{\overset{n1}{}}\underset{j=1}{\overset{m}{}}\alpha _j^i(z)\frac{}{z_j^i}+\frac{}{z_0^0}\hfill \end{array}),$$
where $`_0=(\frac{}{z_1^n},\mathrm{},\frac{}{z_m^n})`$ and the local coordinates $`z_0^0,z_1^0,\mathrm{},z_m^0,\mathrm{},z_1^n,\mathrm{},z_m^n`$ are centered at$`p`$.
Since $`dim𝒟^{(1)}(p)=2m+1`$ we can assume, after a permutation of the $`z`$-coordinates, if necessary, that the real $`m\times m`$ matrix
$$T=\left(\frac{\alpha _i^{n1}}{z_j^n}\right),\text{ for }1im\text{ and }1jm\text{,}$$
has full rank $`m`$, in a small enough neighborhood of zero. We can assume, moreover, that $`\alpha _i^{n1}(0)=0`$, for $`1im`$. Otherwise, replace the coordinate $`z_i^{n1}`$ by $`z_i^{n1}z_0^0\alpha _i^{n1}(0)`$. Now, we can define a new set of centered local coordinates
$$(y_0^0,y_1^0,\mathrm{},y_m^0,\mathrm{},y_1^n,\mathrm{},y_m^n)=\psi (z_0^0,z_1^0,\mathrm{},z_m^0,\mathrm{},z_1^n,\mathrm{},z_m^n)$$
by taking $`y_i^n=\alpha _i^{n1}(z)`$, for $`1im`$, and by taking $`y_j^i=z_j^i`$, as the remaining coordinates. Since the matrix $`T`$ has rank $`m`$, this change of coordinates is indeed a local diffeomorphism. Hence, the distribution$`𝒟`$ is locally equivalent to a distribution spanned, on a small enough neighborhood of zero, by a pair of vector fields that has the following form:
$$(\begin{array}{c}\frac{}{y_1^n}\end{array},\mathrm{},\begin{array}{c}\frac{}{y_m^n}\end{array},\begin{array}{c}\underset{j=1}{\overset{m}{}}y_j^n\frac{}{y_j^{n1}}+\underset{i=0}{\overset{n2}{}}\underset{j=1}{\overset{m}{}}\beta _j^i(y)\frac{}{y_j^i}+\frac{}{y_0^0}\hfill \end{array}).$$
Since $`_0𝒞_1`$ we have $`[_0,𝒟^{(1)}]𝒟^{(1)}`$. But this inclusion clearly implies that, for $`1in2`$ and $`1jm`$, we have $`^2\beta _j^i/y_k^ny_l^n0`$, for $`1km`$ and $`1lm`$. It follows that all functions $`\beta _j^i`$, for $`1in2`$ and $`1jm`$, are affine with respect to the variables $`y_1^n,\mathrm{},y_m^n`$, that is
$$\beta _j^i(y)=\underset{k=1}{\overset{m}{}}a_{jk}^i(\overline{y}^{n1})y_k^n+a_{j0}^i(\overline{y}^{n1}),$$
where $`\overline{y}^{n1}`$ denotes the coordinates $`y_0^0,y_1^0,\mathrm{},y_m^0,\mathrm{},y_1^{n1},\mathrm{},y_m^{n1}`$. Now, define
$$\zeta _k^{n1}=\frac{}{y_j^{n1}}+\underset{i=1}{\overset{n2}{}}\underset{j=1}{\overset{m}{}}a_{jk}^i(\overline{y}^{n1})\frac{}{y_j^i},\text{ for }1km,$$
and
$$\zeta _0^{n1}=\underset{i=1}{\overset{n2}{}}\underset{j=1}{\overset{m}{}}a_{j0}^i(\overline{y}^{n1})\frac{}{y_j^i}+\frac{}{y_0^0}.$$
This definition shows that $`𝒟`$ is locally equivalent to
$$(\begin{array}{c}\frac{}{y_1^n}\end{array},\mathrm{},\begin{array}{c}\frac{}{y_m^n}\end{array},\begin{array}{c}y_1^n\zeta _1^{n1}+\mathrm{}+y_m^n\zeta _m^{n1}+\zeta _0^{n1}\end{array}),$$
where the vector fields $`\zeta _1^{n1},\mathrm{},\zeta _m^{n1},\zeta _0^{n1}`$ are lifts of vector fields on $`J^{n1}(,^m)`$. It follows directly from our construction that $`_0=(\frac{}{y_1^n},\mathrm{},\frac{}{y_m^n})`$. $`\mathrm{}`$
Proof of Theorem 3.2: We will proceed by induction on the integer $`n1`$. For $`n=1`$, the Theorem is a direct consequence of Lemma 4.5. Thus, assume that the Theorem is true for $`n11`$ and consider a rank $`m+1`$ distribution $`𝒟`$, defined on a manifold $`M`$ of dimension $`(n+1)m+1`$, such that each element $`𝒟^{(i)}`$ of its derived flag has constant rank $`(i+1)m+1`$ and contains an involutive subdistribution $`_i𝒟^{(i)}`$ that has constant corank one in $`𝒟^{(i)}`$, for $`0in`$. Let $`p`$ be an arbitrary point in$`M`$.
By Lemma 4.4, the involutive distribution $`_0`$, which has corank one in $`𝒟^{(0)}`$, satisfies $`_0𝒞_1`$. We can thus apply Lemma 4.5, which states that the distribution $`𝒟`$ is equivalent, in a small enough neighborhood of $`p`$, to a distribution spanned on $`J^n(,^m)`$ by a family of vector fields $`(\zeta _1^n,\mathrm{},\zeta _m^n,\zeta _0^n)`$ that has the following form:
$`\zeta _1^n`$ $`=\frac{}{y_1^n},\mathrm{},\zeta _m^n=\frac{}{y_m^n}`$
$`\zeta _0^n`$ $`=y_1^n\zeta _1^{n1}+\mathrm{}+y_m^n\zeta _m^{n1}+\zeta _0^{n1},`$
where the vector fields $`\zeta _1^{n1},\mathrm{},\zeta _m^{n1},\zeta _0^{n1}`$ are lifts of vector fields on $`J^{n1}(,^m)`$. In the rest of the proof we will assume that $`𝒟=(\zeta _1^n,\mathrm{},\zeta _m^n,\zeta _0^n)`$. Note that the $`y`$-coordinates are centered at zero.
The aim of the proof will be to construct a local change of coordinates
$$(x_0^0,x_1^0,\mathrm{},x_m^0,\mathrm{},x_1^n,\mathrm{},x_m^n)=\varphi ^n(y_0^0,y_1^0,\mathrm{},y_m^0,\mathrm{},y_1^n,\mathrm{},y_m^n),$$
a Kumpera-Ruiz normal form $`(\kappa _1^n,\mathrm{},\kappa _m^n,\kappa _0^n)`$ on $`J^n(,^m)`$, and a smooth map $`\mu ^n:J^n(,^m)GL(m+1,)`$, given by an $`(m+1)\times (m+1)`$ matrix $`(\mu _{ij}^n(y))`$, such that
$$\varphi _{}^n(\zeta _i^n)=\underset{j=1}{\overset{m}{}}(\mu _{ij}^n\psi ^n)\kappa _j^n,\text{ for }1im,$$
(7)
and
$$\varphi _{}^n(\zeta _0^n)=\underset{j=0}{\overset{m}{}}(\mu _{0j}^n\psi ^n)\kappa _j^n,$$
(8)
where $`\psi ^n=(\varphi ^n)^1`$ denotes the inverse of the local diffeomorphism $`\varphi ^n`$. Moreover, we will ask the $`x`$-coordinates to be centered at zero, that is $`\varphi ^n(0)=0`$. Observe that we take $`\mu _{i0}^n=0`$, for $`1im`$, in order to transform the canonical distribution of our distribution $`𝒟`$, which is given by $`_0=(\zeta _1^n,\mathrm{},\zeta _m^n)`$ (see Lemma 4.4), into the canonical distribution of the Kumpera-Ruiz normal form, which is given by $`(\kappa _1^n,\mathrm{},\kappa _m^n)`$.
Recall that the vector fields $`\zeta _1^{n1},\mathrm{},\zeta _m^{n1},\zeta _0^{n1}`$ are lifts of vector fields on $`J^{n1}(,^m)`$. If we take $`=(\zeta _1^{n1},\mathrm{},\zeta _m^{n1},\zeta _0^{n1})`$ then we will obtain a decomposition $`𝒟^{(1)}=_0`$ . Since $`𝒟^{(1)}`$ contains an involutive subdistribution $`_1`$ that has constant corank one in $`𝒟^{(1)}`$, it follows directly from the relation $`_0_1`$ (see Lemma 4.2) that $``$ contains an involutive subdistribution that has constant corank one in $``$. In fact, it is easy to prove (using the relations $`𝒞_i𝒞_{i+1}`$ and $`_i=𝒞_i`$) that if $`𝒟`$ satisfies the conditions of Theorem 3.2 on $`J^n(,^m)`$ then $``$ satisfies the conditions of this Theorem on $`J^{n1}(,^m)`$. Now, recall that we have assumed that the Theorem is true on $`J^{n1}(,^m)`$. The distribution $``$ is thus locally equivalent to a distribution spanned by a Kumpera-Ruiz normal form on $`J^{n1}(,^m)`$, centered at zero. It follows that there exists a local diffeomorphism
$$(x_0^0,x_1^0,\mathrm{},x_m^0,\mathrm{},x_1^{n1},\mathrm{},x_m^{n1})=\varphi ^{n1}(y_0^0,y_1^0,\mathrm{},y_m^0,\mathrm{},y_1^{n1},\mathrm{},y_m^{n1}),$$
a Kumpera-Ruiz normal form $`(\kappa _1^{n1},\mathrm{},\kappa _m^{n1},\kappa _0^{n1})`$ on $`J^{n1}(,^m)`$, and a smooth map $`\mu ^{n1}:J^{n1}(,^m)GL(m+1,)`$ such that
$$\varphi _{}^{n1}(\zeta _i^{n1})=\underset{j=0}{\overset{m}{}}(\mu _{ij}^{n1}\psi ^{n1})\kappa _j^{n1}\text{,\hspace{1em}for }0im\text{,}$$
where $`\psi ^{n1}=(\varphi ^{n1})^1`$ denotes the inverse of the local diffeomorphism $`\varphi ^{n1}`$. Note that we have $`\varphi ^{n1}(0)=0`$.
The following Lemma can be easily proved by a direct computation.
###### Lemma 4.6 (triangular tangent maps)
Let $`\varphi ^n=(\varphi ^{n1},\varphi _1^n,\mathrm{},\varphi _m^n)^\mathrm{T}`$ be a diffeomorphism of $`J^n(,^m)`$ such that its first $`nm+1`$ components, which are given by $`\varphi ^{n1}`$, depend on the first $`nm+1`$ coordinates only. Moreover, let $`f`$ be a vector field on $`J^n(,^m)`$ of the form $`f=\alpha f^{n1}+f_n`$, where $`\alpha `$ is a smooth function on $`J^n(,^m)`$, the vector field $`f^{n1}`$ is the lift of a vector field on $`J^{n1}(,^m)`$, and the only non-zero components of $`f_n`$ are those that multiply $`\frac{}{y_1^n},\mathrm{},\frac{}{y_m^n}`$. Then, we have
$$\varphi _{}^n(f)=(\alpha \psi ^n)\varphi _{}^{n1}(f^{n1})+\underset{i=1}{\overset{m}{}}\left((\mathrm{L}_f\varphi _i^n)\psi ^n\right)\frac{}{x_i^n}.$$
(9)
Note that the vector field $`\varphi _{}^{n1}(f^{n1})`$ is lifted to $`J^n(,^m)`$ and that the coordinates $`x_1^n,\mathrm{},x_m^n`$ are those given by $`\varphi _1^n,\mathrm{},\varphi _m^n`$, respectively.
*Regular case:* If $`\mu _{00}^{n1}(0)0`$ then we can complete $`\varphi ^{n1}`$ to a zero preserving diffeomorphism $`\varphi ^n`$ of $`J^n(,^m)`$ by taking $`\varphi ^n=(\varphi ^{n1},\varphi _1^n,\mathrm{},\varphi _m^n)^\mathrm{T}`$, where
$$\varphi _j^n(y)=\frac{_{i=1}^m\mu _{ij}^{n1}y_i^n+\mu _{0j}^{n1}}{_{i=1}^m\mu _{i0}^{n1}y_i^n+\mu _{00}^{n1}}\frac{\mu _{0j}^{n1}(0)}{\mu _{00}^{n1}(0)}\text{,\hspace{1em}for }1jm\text{.}$$
It is easy to check, using Lemma 4.7 below, that $`\varphi ^n`$ is a local diffeomorphism (because $`\mu ^{n1}`$ is invertible). In this case, we define $`c_i^n=(\mu _{0i}^{n1}/\mu _{00}^{n1})(0)`$ for $`1im`$, $`\mu _{ij}^n=\mathrm{L}_{\zeta _i^{n1}}\varphi _j^n`$ for $`0im`$ and $`1jm`$, and $`\mu _{00}^n=_{i=1}^m\mu _{i0}^{n1}y_i^n+\mu _{00}^{n1}`$. Moreover, the Kumpera-Ruiz normal form $`(\kappa _1^n,\mathrm{},\kappa _m^n,\kappa _0^n)`$ is defined to be the regular prolongation, with parameter $`c^n=(c_1^n,\mathrm{},c_m^n)`$, of $`(\kappa _1^{n1},\mathrm{},\kappa _m^{n1},\kappa _0^{n1})`$. Let us check that, in this case, relation (8) holds. By relation (9) we have:
$`\varphi _{}^n(\zeta _0^n)`$ $`=`$ $`\underset{i=1}{\overset{m}{}}(y_i^n\psi ^n)\varphi _{}^{n1}(\zeta _i^{n1})+\varphi _{}^{n1}(\zeta _0^{n1})+\underset{i=1}{\overset{m}{}}\left((\mathrm{L}_{\zeta _0^n}\varphi _i^n)\psi ^n\right)\frac{}{x_i^n}`$
$`=`$ $`\underset{i=1}{\overset{m}{}}(y_i^n\psi ^n)\left(\underset{j=0}{\overset{m}{}}(\mu _{ij}^{n1}\psi ^{n1})\kappa _j^{n1}\right)+\underset{j=0}{\overset{m}{}}(\mu _{0j}^{n1}\psi ^{n1})\kappa _j^{n1}`$
$`+\underset{i=1}{\overset{m}{}}\left(\mu _{0i}^n\psi ^n\right)\kappa _i^n`$
$`=`$ $`\underset{j=0}{\overset{m}{}}\left(\left(\underset{i=1}{\overset{m}{}}\mu _{ij}^{n1}y_i^n+\mu _{0j}^{n1}\right)\psi ^n\right)\kappa _j^{n1}+\underset{i=1}{\overset{m}{}}\left(\mu _{0i}^n\psi ^n\right)\kappa _i^n`$
$`=`$ $`\left(\left(\underset{i=1}{\overset{m}{}}\mu _{i0}^{n1}y_i^n+\mu _{00}^{n1}\right)\psi ^n\right)`$
$`\times \left(\underset{j=1}{\overset{m}{}}\left({\displaystyle \frac{_{i=1}^m\mu _{ij}^{n1}y_i^n+\mu _{0j}^{n1}}{_{i=1}^m\mu _{i0}^{n1}y_i^n+\mu _{00}^{n1}}}\psi ^n\right)\kappa _j^{n1}+\kappa _0^{n1}\right)+\underset{i=1}{\overset{m}{}}\left(\mu _{0i}^n\psi ^n\right)\kappa _i^n`$
$`=`$ $`\left(\mu _{00}^n\psi ^n\right)\left(\underset{j=1}{\overset{m}{}}(x_j^n+c_j^n)\kappa _j^{n1}+\kappa _0^{n1}\right)+\underset{i=1}{\overset{m}{}}\left(\mu _{0i}^n\psi ^n\right)\kappa _i^n`$
$`=`$ $`\underset{i=0}{\overset{m}{}}\left(\mu _{0i}^n\psi ^n\right)\kappa _i^n.`$
Moreover, for $`1im`$, we have
$$\varphi _{}^n(\zeta _i^n)=\underset{j=1}{\overset{m}{}}\left((\mathrm{L}_{\zeta _i^n}\varphi _j^n)\psi ^n\right)\kappa _j^n=\underset{j=1}{\overset{m}{}}\left(\mu _{ij}^n\psi ^n\right)\kappa _j^n.$$
It follows that both (7) and (8) hold.
*Singular case:* Suppose now that $`\mu _{00}^{n1}(0)=0`$. Since the matrix $`\mu ^{n1}`$ is invertible in a small enough neighborhood of zero, we can assume that there exists an integer $`1im`$ such that $`\mu _{0i}^{n1}(0)0`$. After a permutation of the coordinates $`y_1^n,\mathrm{},y_m^n`$, if necessary, we can assume that $`\mu _{0m}^{n1}(0)0`$. Now, like in the regular case, we can complete $`\varphi ^{n1}`$ to a zero preserving diffeomorphism$`\varphi ^n`$ of $`J^n(,^m)`$ by taking $`\varphi ^n=(\varphi ^{n1},\varphi _1^n,\mathrm{},\varphi _m^n)^\mathrm{T}`$, where
$$\varphi _j^n(y)=\frac{_{i=1}^m\mu _{ij}^{n1}y_i^n+\mu _{0j}^{n1}}{_{i=1}^m\mu _{im}^{n1}y_i^n+\mu _{0m}^{n1}}\frac{\mu _{0j}^{n1}(0)}{\mu _{0m}^{n1}(0)}\text{,\hspace{1em}for }1jm1\text{,}$$
and
$$\varphi _m^n(y)=\frac{_{i=1}^m\mu _{i0}^{n1}y_i^n+\mu _{00}^{n1}}{_{i=1}^m\mu _{im}^{n1}y_i^n+\mu _{0m}^{n1}}.$$
In this case, we define $`c_i^n=(\mu _{0i}^{n1}/\mu _{0m}^{n1})(0)`$ for $`1im1`$. Observe that we can take $`c_m^n=0`$ because $`\mu _{00}^{n1}(0)=0`$. We take $`\mu _{ij}^n=\mathrm{L}_{\zeta _i^{n1}}\varphi _j^n`$, for $`0im`$ and $`1jm`$, and $`\mu _{00}^n=_{i=1}^m\mu _{im}^{n1}y_i^n+\mu _{0m}^{n1}`$. Moreover, the Kumpera-Ruiz normal form $`(\kappa _1^n,\mathrm{},\kappa _m^n,\kappa _0^n)`$ is defined to be the singular prolongation, with parameter $`c^n=(c_1^n,\mathrm{},c_{m1}^n,0)`$, of $`(\kappa _1^{n1},\mathrm{},\kappa _m^{n1},\kappa _0^{n1})`$. Let us check that relation (8) holds. We have:
$`\varphi _{}^n(\zeta _0^n)`$ $`=`$ $`\underset{i=1}{\overset{m}{}}(y_i^n\psi ^n)\varphi _{}^{n1}(\zeta _i^{n1})+\varphi _{}^{n1}(\zeta _0^{n1})+\underset{i=1}{\overset{m}{}}\left((\mathrm{L}_{\zeta _0^n}\varphi _i^n)\psi ^n\right)\frac{}{x_i^n}`$
$`=`$ $`\underset{j=0}{\overset{m}{}}\left(\left(\underset{i=1}{\overset{m}{}}\mu _{ij}^{n1}y_i^n+\mu _{0j}^{n1}\right)\psi ^n\right)\kappa _j^{n1}+\underset{i=1}{\overset{m}{}}\left(\mu _{0i}^n\psi ^n\right)\kappa _i^n`$
$`=`$ $`\left(\underset{i=1}{\overset{m}{}}\left(\mu _{im}^{n1}y_i^n+\mu _{0m}^{n1}\right)\psi ^n\right)`$
$`\times (\underset{j=1}{\overset{m1}{}}({\displaystyle \frac{_{i=1}^m\mu _{ij}^{n1}y_i^n+\mu _{0j}^{n1}}{_{i=1}^m\mu _{im}^{n1}y_i^n+\mu _{0m}^{n1}}}\psi ^n)\kappa _j^{n1}+\kappa _m^{n1}`$
$`+({\displaystyle \frac{_{i=1}^m\mu _{i0}^{n1}y_i^n+\mu _{00}^{n1}}{_{i=1}^m\mu _{im}^{n1}y_i^n+\mu _{0m}^{n1}}}\psi ^n)\kappa _0^{n1})+_{i=1}^m(\mu _{0i}^n\psi ^n)\kappa _i^n`$
$`=`$ $`\left(\mu _{00}^n\psi ^n\right)\left(\underset{j=1}{\overset{m1}{}}(x_j^n+c_j^n)\kappa _j^{n1}+\kappa _m^{n1}+x_m^n\kappa _0^{n1}\right)+\underset{i=1}{\overset{m}{}}\left(\mu _{0i}^n\psi ^n\right)\kappa _i^n`$
$`=`$ $`\underset{i=0}{\overset{m}{}}\left(\mu _{0i}^n\psi ^n\right)\kappa _i^n.`$
Moreover, like in the regular case, we have
$$\varphi _{}^n(\zeta _i^n)=\underset{j=1}{\overset{m}{}}\left((\mathrm{L}_{\zeta _i^n}\varphi _j^n)\psi ^n\right)\kappa _j^n=\underset{j=1}{\overset{m}{}}\left(\mu _{ij}^n\psi ^n\right)\kappa _j^n,$$
for $`1im`$. It follows that relations (7) and (8) hold in both cases.
We have thus proved that the conditions of Theorem 3.2 are sufficient for converting a distribution into extended Kumpera-Ruiz normal form. It is straightforward to check that these conditions are also necessary. $`\mathrm{}`$
###### Lemma 4.7 (Möbius transformations)
Consider a real $`n\times n`$ matrix $`M`$ that has the following form:
$$M=\left(\begin{array}{cc}A& b\\ c& d\end{array}\right),$$
where $`c`$ is a row vector and $`b`$ a column vector, both of dimension $`n1`$, the real constant $`d`$ is non-zero, and $`A`$ is a real $`(n1)\times (n1)`$ matrix. The linear fractional transformation$`\phi `$, from$`^{n1}`$ into$`^{n1}`$, defined in a small enough neighborhood of zero by $`\phi (x)=\left(Ax+b\right)/\left(cx+d\right)`$ is a local diffeomorphism if and only if the matrix $`M`$ is invertible.
Proof: We have $`\phi _{}(0)=\left(Adbc\right)/d^2`$ and thus $`det\phi _{}(0)=(1/d^2)det(Adbc)`$. But $`detM=(1/d^{n2})det(Adbc)`$. Hence $`det\phi _{}(0)0`$ if and only if $`detM0`$. $`\mathrm{}`$
Proof of Theorem 1.1: Let $`𝒟`$ be a distribution of rank $`m+1`$, defined on a manifold $`M`$ of dimension $`(n+1)m+1`$, that satisfies the conditions of Theorem 1.1. In particular, the distribution $`𝒟`$ satisfies the conditions of Theorem 3.2. Therefore there exists a Kumpera-Ruiz normal form $`(\kappa _1^n,\mathrm{},\kappa _m^n,\kappa _0^n)`$ on $`J^n(,^m)`$, defined on a small enough neighborhood of zero, that is equivalent to the distribution $`𝒟`$ considered on a small enough neighborhood of any point $`p`$ in $`M`$. We will thus assume that $`𝒟=(\kappa _1^n,\mathrm{},\kappa _m^n,\kappa _0^n)`$.
Now, if we exclude Engel’s case ($`m=1`$ and $`n=2`$), for which Theorem 3.2 is well known to be true, it is straightforward to check that if the sequence of prolongations that defines our Kumpera-Ruiz normal form contains a singular prolongation then there exists some integer $`2in`$ such that the Lie flag of $`𝒟`$ satisfies
$$dim𝒟_i(0)<(i+1)m+1.$$
It thus follows that $`(\kappa _1^n,\mathrm{},\kappa _m^n,\kappa _0^n)`$ has necessarily been obtained by a sequence of regular prolongations from the canonical contact system on $`J^1(,^m)`$. That is:
$$𝒟=(\begin{array}{c}\frac{}{x_1^n}\hfill \end{array},\mathrm{},\begin{array}{c}\frac{}{x_m^n}\hfill \end{array},\begin{array}{c}\underset{i=0}{\overset{n1}{}}\underset{j=1}{\overset{m}{}}\left(x_j^{i+1}+c_j^{i+1}\right)\frac{}{x_j^i}+\frac{}{x_0^0}\hfill \end{array}).$$
(10)
What remains to prove now is that the distribution (10), which is defined on a small enough neighborhood of zero, is locally equivalent to the canonical contact system on $`J^n(,^m)`$, also considered on a small enough neighborhood of zero. In other words, we have to normalize all constants $`c_j^{i+1}`$ by making them equal to zero. To this end, observe that the Lie algebra
$$𝔤=\underset{}{span}\{\kappa _1^n,\mathrm{},\kappa _m^n,\kappa _1^{n1},\mathrm{},\kappa _m^{n1},\mathrm{},\kappa _1^1,\mathrm{},\kappa _m^1,\kappa _0^n\},$$
of dimension $`(n+1)m+1`$, generated by the vector fields $`(\kappa _1^n,\mathrm{},\kappa _m^n,\kappa _0^n)`$ defining (10), has the same structure constants independently of the values of the parameters $`c_j^{i+1}`$. Indeed, the only non-zero Lie brackets are those given by the relations $`[\kappa _j^i,\kappa _0^n]=\kappa _j^{i1}`$, for $`1in`$ and $`1jm`$. By Cartan’s theorem on equivalence of frames, our distribution is locally equivalent to the canonical contact system at zero (see e.g. for a modern account on Cartan’s equivalence method). $`\mathrm{}`$
## Appendix A Appendix
###### Lemma A.1 (Bryant)
Let $`𝒟`$ be a distribution such that $`𝒟^{(0)}`$ and $`𝒟^{(1)}`$ have constant ranks $`d_0`$ and $`d_1`$, respectively. Put $`r_0=d_1d_0`$. Assume that the distribution $`𝒟`$ contains a subdistribution $`𝒟`$ that has constant corank one in $`𝒟`$ and satisfies $`[,]𝒟`$. If $`r_03`$ then $``$ is involutive.
Proof: Assume that $`𝒟`$ contains a subdistribution $`𝒟`$ that has constant corank one in $`𝒟`$ and satisfies $`[,]𝒟`$. By Lemma 4.1, the rank of the characteristic distribution $`𝒞_0`$ of $`𝒟^{(0)}`$ is constant and $`𝒞_0`$. Therefore, there exists a local basis $`(f_1,\mathrm{},f_{d_0})`$ of $`𝒟`$ such that
$$𝒞_0=(f_1,\mathrm{},f_{c_0})\text{ and }=(f_1,\mathrm{},f_{d_01}),$$
Since $``$ satisfies $`[,]𝒟`$ we have $`𝒟^{(1)}=𝒟^{(0)}+[f_{d_0},]`$ or, more precisely,
$$𝒟^{(1)}=𝒟^{(0)}([f_{d_0},f_{c_0+1}],\mathrm{},[f_{d_0},f_{d_01}]).$$
(11)
Note that the relation $`d_1d_0+3`$ implies that $`card\{c_0+1,\mathrm{},d_01\}3`$.
We want to prove that $`[,]`$. Let us first prove that $`[f_i,f_j]`$, for $`c_0+1id_01`$ and $`c_0+1jd_01`$. Consider an arbitrary triple $`f_i`$, $`f_j`$, and $`f_k`$ of vector fields such that the indices $`i`$, $`j`$, and $`k`$ are pairwise different and contained in $`\{c_0+1,\mathrm{},d_01\}`$ (it is important to stress that such a triple exists because $`d_1d_0+3`$). It follows from relation (11) that the three vector fields $`[f_{d_0},f_i]`$, $`[f_{d_0},f_j]`$, and $`[f_{d_0},f_k]`$ are linearly independent $`mod𝒟^{(0)}`$. Moreover, since $`[,]𝒟^{(0)}`$, there exist three smooth functions $`a`$, $`b`$, and $`c`$ such that $`[f_i,f_j]=af_{d_0}mod`$, $`[f_j,f_k]=bf_{d_0}mod`$, and $`[f_k,f_i]=cf_{d_0}mod`$. The Jacobi identity gives:
$$[f_i,[f_j,f_k]]+[f_j,[f_k,f_i]]+[f_k,[f_i,f_j]]=0,$$
which implies that $`[f_i,bf_{d_0}]+[f_j,cf_{d_0}]+[f_k,af_{d_0}]`$ belongs to $`𝒟^{(0)}`$, and thus that $`b[f_i,f_{d_0}]+c[f_j,f_{d_0}]+a[f_k,f_{d_0}]`$ belongs to $`𝒟^{(0)}`$. The latter relation implies that $`a`$, $`b`$, and $`c`$ are identically zero because $`[f_{d_0},f_i]`$, $`[f_{d_0},f_j]`$, and $`[f_{d_0},f_k]`$ are linearly independent $`mod𝒟^{(0)}`$. It follows that we have $`[f_i,f_j]`$, for $`c_0+1id_01`$ and $`c_0+1jd_01`$.
Since $`[𝒞_0,𝒞_0]=𝒞_0`$, what remain to prove is that $`[𝒞_0,]=`$. The proof follows again from the Jacobi identity, applied to any triple $`f_i`$, $`f_j`$, and $`f_k`$ of pairwise linearly independent vector fields such that $`f_i`$ belongs to $`𝒞_0`$ and both $`f_j`$ and $`f_k`$ belong to $``$ but do not belong to $`𝒞_0`$. $`\mathrm{}`$
## Appendix B Appendix
In this appendix we will provide, following Bryant , a way to check the conditions of Corollaries 2.3 and 2.4. Indeed, we will show how to verify whether or not the Engel rank equals $`1`$, and how to construct explicitly the characteristic distribution of $`𝒟`$ and — when it exists — the unique corank one subdistribution $`𝒟`$ satisfying $`[,]𝒟`$.
Consider a distribution $`𝒟`$ of constant rank $`d_0`$, defined on a manifold of dimension $`N`$. Let $`\omega _1,\mathrm{},\omega _{s_0}`$, where $`s_0=Nd_0`$, be differential 1-forms locally spanning $`𝒟^{}`$, the annihilator of $`𝒟`$, which we denote by
$$𝒟^{}=(\omega _1,\mathrm{},\omega _{s_0}).$$
We will denote by $``$ the Pffafian system generated by $`\omega _1,\mathrm{},\omega _{s_0}`$.
For any form $`\omega 𝒟^{}`$, we put
$$𝒲(\omega )=\{f𝒟:f\mathrm{}d\omega 𝒟^{}\}.$$
Clearly, the characteristic distribution $`𝒞`$ of $`𝒟`$ is given by
$$𝒞=\underset{i=1}{\overset{s_0}{}}𝒲(\omega _i).$$
(12)
Now assume that $`𝒟^{(1)}`$ is of constant rank $`d_1>d_0`$, that is $`r_01`$, or, equivalently, that the first derived system $`^{(1)}`$ is of constant rank smaller than $`s_0`$. By a direct calculation we can check (see e.g. ) that the Engel rank of the distribution $`𝒟`$, or of the corresponding Pfaffian system $``$, equals $`1`$ at $`p`$ if and only if
$$(d\omega _id\omega _j)(p)=0mod,$$
(13)
for any $`1ijs_0`$.
Now let us choose a family of differential 1-forms $`\omega _1,\mathrm{},\omega _{r_0},\omega _{r_0+1},\mathrm{},\omega _{s_0}`$ such that $`(𝒟^{(0)})^{}=(\omega _1,\mathrm{},\omega _{s_0})`$ and $`(𝒟^{(1)})^{}=(\omega _{r_0+1},\mathrm{},\omega _{s_0})`$. Independently of the value of $`r_02`$, the unique distribution $``$ satisfying $`[,]𝒟`$ is given, as shown by Bryant , by
$$=\underset{i=1}{\overset{r_0}{}}𝒲(\omega _i).$$
(14)
In fact, Bryant has also proved that it is enough to take in the above sum only two terms corresponding to any $`1i<jr_0`$. In order to verify, in the case $`r_0=2`$, the conditions of Corollary 2.3 we have additionally to check the involutivity of this explicitly calculable distribution.
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# Steady State of microemulsions in shear flow
## Abstract
Steady-state properties of microemulsions in shear flow are studied in the context of a Ginzburg-Landau free-energy approach. Explicit expressions are given for the structure factor and the time correlation function at the one loop level of approximation. Our results predict a four-peak pattern for the structure factor, implying the simultaneous presence of interfaces aligned with two different orientations. Due to the peculiar interface structure a non-monotonous relaxation of the time correlator is also found.
Self-assembling amphiphilic systems are binary or ternary mixtures with a surfactant forming interfaces between the other fluid components. These systems show a very rich phase behavior with ordered and structured disordered phases at different temperatures and relative concentrations . Particularly interesting for the applications is the microemulsion phase where coherent domains of oil and water on scales between 100 and 1000 $`\AA `$ form disordered isotropic intertwined structures on larger scales .
While a satisfactory comprehension of the equilibrium behavior of amphiphilic systems is now available the dynamical features are far less understood. In particular, in microemulsions, the existence of correlated mesoscopic structures makes the effects of imposed flows very different than in simple fluids. Interfaces between oil and water are elongated in the flow direction and an excess stress appears in the system. This can give a peculiar dynamics with unusual relaxation properties and a non-Newtonian rheological behavior . Moreover, from a statistical mechanics point of view, complex fluids in external flow are an interesting example of an out of equilibrium driven system .
In this Letter we study the behavior of microemulsions in shear flow with a continuum free-energy approach similar to that used by Fredrickson and Larson for the study of rheology of block copolymers and by Onuki and Kawasaki and others for the critical properties. A Ginzburg-Landau model for describing the rheological behavior of ternary mixtures was considered by Pätzold and Dawson . However the results of are based on an evaluation of the structure factor obtained numerically or analytically only in the limit of vanishing shear rate $`\gamma `$. Here we solve exactly the model for any value of $`\gamma `$ showing that in an intermediate shear range the structure factor is characterized by four pronounced peaks in the plane of the shear and flow directions. This new phenomenon can be interpreted in terms of a complex spatial pattern where interfaces with two different orientations coexist; their relative abundance is tuned by $`\gamma `$. We also consider the behavior of the two-time correlator. Due to the coupling between the flow and the interface structure, this function has a remarkable non monotonic behavior characterized by the presence of maxima preceding a fast asymptotic decay.
The Ginzburg-Landau free-energy generally used to describe the equilibrium properties of ternary mixtures is
$`\{\phi \}={\displaystyle }d^3x\{{\displaystyle \frac{a}{2}}\phi ^2+{\displaystyle \frac{b}{4!}}\phi ^4`$ $`+`$ $`{\displaystyle \frac{1}{2}}(g_0+g_2\phi ^2)\phi ^2`$ (1)
$`+`$ $`{\displaystyle \frac{c}{2}}(^2)^2\phi \}`$ (2)
where the scalar field $`\phi `$ represents the concentration difference between oil and water components. The amount of surfactant present in the system is related to the value of $`g_0`$ which can be negative, favoring the appearing of interfaces. The term proportional to $`c`$ makes stable the free energy at large momenta and weights the curvature of interfaces. The quadratic part of the free-energy in the disordered phase with $`a>0`$ and $`g_0^2<4ac`$ gives a suitable description of the microemulsion phase; in particular the space correlation function
$$G(r)\frac{e^{r/\xi }}{r}\mathrm{sin}\frac{2\pi r}{d}$$
(3)
well fits the experimental data . The characteristic length $`d=2\pi (1/2\sqrt{a/c}g_0/4c)^{1/2}`$ represents the size of coherent oil or water regions which are correlated up to a distance $`\xi =(1/2\sqrt{a/c}+g_0/4c)^{1/2}`$. At negative $`g_0`$, the structure factor obtained by Fourier transforming Eq. (2) has a peak at a finite value of the momentum given by $`k_M=\sqrt{|g_0|/2c}`$. Finally, the non-linear terms in eq.(1) describe the possible effects due to mode coupling and become more important in approaching the phase boundary.
The kinetic behavior of the mixture, neglecting hydrodynamical effects , is described by the convection-diffusion equation
$$\frac{\phi }{t}+\stackrel{}{}(\phi \stackrel{}{v})=\mathrm{\Gamma }^2\frac{\delta }{\delta \phi }+\eta $$
(4)
where the velocity field is given by
$$\stackrel{}{v}=\gamma y\stackrel{}{e}_x$$
(5)
$`\gamma `$ being the shear rate and $`\stackrel{}{e}_x`$ the unit vector in the flow direction. The stochastic term $`\eta `$ is a gaussian white noise, representing thermal fluctuations, with zero mean and correlation $`\eta (\stackrel{}{r},t)\eta (\stackrel{}{r}^{},t^{})=2T\mathrm{\Gamma }^2\delta (\stackrel{}{r}\stackrel{}{r}^{})\delta (tt^{})`$, as required by the fluctuation-dissipation theorem, where $`\mathrm{}`$ means the ensemble average. Here $`\mathrm{\Gamma }`$ is a mobility coefficient and $`T`$ the temperature of the heat bath in contact with the fluid.
The cubic terms appearing in the functional derivative of Eq. (3) will be treated in the one-loop approximation which describes appropriately disordered phases. The resulting equation for the Fourier components $`\phi (\stackrel{}{k},t)`$ is given by
$`{\displaystyle \frac{\phi (\stackrel{}{k},t)}{t}}`$ $``$ $`\gamma k_x{\displaystyle \frac{\phi (\stackrel{}{k},t)}{k_y}}=\mathrm{\Gamma }k^2[X(k)+g_2k^2S_0(t)`$ (6)
$`+`$ $`({\displaystyle \frac{b}{2}}S_0(t)+g_2S_2(t))]\phi (\stackrel{}{k},t)+\eta (\stackrel{}{k},t)`$ (7)
where $`X(k)=ck^4+g_0k^2+a`$, $`S_0(t)`$ and $`S_2(t)`$ are defined self-consistently by the relations
$$S_0(t)=_{|\stackrel{}{k}|<\mathrm{\Lambda }}\frac{d\stackrel{}{k}}{(2\pi )^D}C(\stackrel{}{k},t)$$
(8)
$$S_2(t)=_{|\stackrel{}{k}|<\mathrm{\Lambda }}\frac{d\stackrel{}{k}}{(2\pi )^D}k^2C(\stackrel{}{k},t)$$
(9)
and $`\mathrm{\Lambda }`$ is a phenomenological ultraviolet cut-off. Starting from Eq.(5) a straightforward calculation give the dynamical equations for the various correlation functions . We first focus on the properties of the structure factor $`\phi (\stackrel{}{k},t)\phi (\stackrel{}{k},t)`$. Using the method of characteristics the $`t\mathrm{}`$ stationary expression of the structure factor $`C(\stackrel{}{k})`$ can be calculated as
$$C(\stackrel{}{k})=_0^{\mathrm{}}𝑑z2\mathrm{\Gamma }T\stackrel{}{K}^2(z)e^{_0^z𝑑s2\mathrm{\Gamma }\stackrel{}{K}^2(s)[a^R+g^R\stackrel{}{K}^2(s)+c\stackrel{}{K}^4(s)]}$$
(10)
where $`\stackrel{}{K}(s)=(k_x,k_y+\gamma k_xs,k_z)`$, $`\stackrel{}{K}^4(\stackrel{}{K}^2)^2`$, $`a^R=a+b/2S_0^{\mathrm{}}+g_2S_2^{\mathrm{}}`$, $`g^R=g_0+g_2S_0^{\mathrm{}}`$, and $`S_0^{\mathrm{}}`$ and $`S_2^{\mathrm{}}`$ are the limit for $`t\mathrm{}`$ of $`S_0(t)`$ and $`S_2(t)`$.
Since $`\xi `$ and $`d`$ are the physical quantities to be compared with experimental data, the model parameters are tuned with the following procedure: We first calculate the one-loop renormalization of the bare coefficients $`a`$ and $`g`$ without shear and express the correlation lengths $`\xi `$ and $`d`$ in terms of the renormalized parameters. Then, for a given choice of $`\xi `$ and $`d`$, we invert the above mentioned relations and compute the values of $`a,g_0`$ to be used in the case with shear. In the following we show results for the case $`\xi =7,d=5`$. The other parameters are set to the values $`c=2`$, $`b=g_2=T=1`$. Similar results have been obtained for different choices of parameters with $`d/\xi `$ in the range between 0.5 and 2.5.
The effects of the flow on $`C(\stackrel{}{k})`$ are shown in Fig. 1, where the projections of the three-dimensional structure factor on the planes $`k_y=0`$ and $`k_z=0`$ are plotted at different $`\gamma `$. At $`\gamma =0`$ the structure factor is isotropic and its shape on each cartesian plane is that of a circular volcano with radius $`k_M`$. When shear is applied, the projection on the $`k_x=0`$ plane gives the same results of the $`\gamma =0`$ case because the velocity field is along the $`x`$ direction. On the $`k_z=0`$ plane, instead, this pattern is progressively distorted as the shear rate is increased. For small $`\gamma `$ (see the case $`\gamma =0.25`$ in the left column of Fig. 1) the edge of the volcano becomes elliptical and slightly depressed along the $`k_x=k_y`$ direction so that two weak symmetric maxima located at $`\theta =\mathrm{arctan}(k_y/k_x)=\pi /4,3\pi /4`$ are developed. This is confirmed by a small $`\gamma `$ expansion of Eq. (10). As $`\gamma `$ is increased each of these two maxima is splitted into two peaks, as shown in Fig. 1 at $`\gamma =2`$ and $`\gamma =4`$. By increasing the shear rate, the maxima located at $`k_x0`$ become comparatively more important while the other peaks rotate clockwise and decrease their amplitude linearly in $`\gamma `$ until they disappear. Indeed, in the limit $`\gamma \mathrm{}`$, since terms proportional to powers of $`\gamma k_x`$ damp the exponential term on the r.h.s. of Eq. (10), only the maxima of $`C(\stackrel{}{k})`$ with $`k_x=0`$ survive. Fig. 1 shows that this is already observable for $`\gamma =8`$. Then, by letting $`k_x=0`$ in Eq. (10) we find the peaks located at $`k_y=\pm k_M`$. The description is completed with the projection of the structure factors on the plane $`k_y=0`$ shown in the right column of Fig. 1. Here two peaks at $`k_x=0,k_z=\pm k_M`$, are observed to become sharper and sharper as the shear is increased.
A peak of $`C(k)`$ is generally interpreted as the signature of a characteristic length in the system which is proportional to the inverse of its position. In this case, since the system is not isotropic, to each maximum one associates three lengths, one for each space direction. Due to the symmetry $`\stackrel{}{k}\stackrel{}{k}`$ only the peaks not related by reflection around the origin can be considered. For large shear the existence of a single couple of maxima at $`k_x=0`$ signals that interfaces are aligned along the flow. In the transverse directions the characteristic lengths are the same as without shear. For intermediate values of $`\gamma `$ the additional peaks at $`(\stackrel{~}{k}_x,\stackrel{~}{k}_y,\stackrel{~}{k}_z)`$ reveal the presence of interfaces oriented with an angle $`\alpha =\mathrm{arctan}(\stackrel{~}{k}_x/\stackrel{~}{k}_y)`$ with respect to the flow, besides those aligned along the $`x`$ direction. As $`\gamma `$ is increased these features are progressively tilted in the direction of the flow while their relative abundance diminishes, as suggested by the behavior of the maxima with $`k_x0`$ previously discussed. The existence of a four-fold peaked structure factor has been also reported in scattering experiments on segregating mixtures . For small $`\gamma `$ it is interesting to note that, although the microemulsion is almost isotropic, the depression of $`C(\stackrel{}{k})`$ along $`k_x=k_y`$ indicates a slight predominance of interfaces directed against the flow at $`\alpha 3\pi /4`$.
Stretching of domains requires work against surface tension and results in an increase $`\mathrm{\Delta }\eta `$ of the viscosity . The excess viscosity is generally defined as $`\mathrm{\Delta }\eta =\sigma _{xy}/\gamma `$ with the shear stress given by $`\sigma _{xy}=\frac{d\stackrel{}{k}}{(2\pi )^D}k_xk_y(g^R+2ck^2)C(\stackrel{}{k})`$ . In fig. 2 $`\sigma _{xy}`$ is plotted as a function of $`\gamma `$. Shear thinning is observed. For small $`\gamma `$, $`\sigma _{xy}`$ grows linearly with $`\gamma `$ as also found in . This behaviour is observed in correspondence of a structure factor similar to that shown in Fig. 1 at $`\gamma =0.25`$. When $`C(\stackrel{}{k})`$ develops four maxima at $`k_z=0`$, $`\sigma _{xy}`$ keeps increasing with a much smaller effective exponent consistent with the value $`0.13`$. For very large $`\gamma `$, when the peaks at $`k_x=0`$ alone survive, $`\sigma _{xy}`$ decreases to zero because the excess stress vanishes when the interfaces are aligned with the flow. A similar behavior is shown by the first normal stress $`N_1=\frac{d\stackrel{}{k}}{(2\pi )^D}(k_y^2k_x^2)(g^R+2ck^2)C(\stackrel{}{k})`$ .
Next we consider the dynamical properties of the microemulsion phase. In the steady state the two-time correlation function $`𝒟(\stackrel{}{k},\stackrel{}{k}^{},t)=\phi (\stackrel{}{k},t)\phi (\stackrel{}{k}^{},0)=D(\stackrel{}{k},t)\delta \left(\stackrel{}{K}(t)+k^{}\right)`$ can be computed through
$$D(\stackrel{}{k},t)=C\left(\stackrel{}{K}(t)\right)e^{_0^t\mathrm{\Gamma }\stackrel{}{K}^2(s)[a^R+g^R\stackrel{}{K}^2(s)+c\stackrel{}{K}^4(s)]𝑑s}$$
(11)
$`D`$ is the product of two terms: the structure factor evaluated at the translated momentum $`\stackrel{}{K}(t)`$ and an exponential term. For a given $`\stackrel{}{k}=\stackrel{}{K}(0)`$, $`\stackrel{}{K}(t)`$ moves along the $`k_y`$ direction upward or downward depending on the sign of $`k_x`$, as represented schematically in the lower inset of Fig. 3. Therefore, when $`\stackrel{}{K}(t)`$ crosses the edge of the structure factor, $`C\left(\stackrel{}{K}(t)\right)`$ shows a pronounced maximum. This feature is peculiar to microemulsions where the presence of interfaces gives the patterns of structure factor discussed above and cannot be observed in simple fluids. The exponential term in Eq. (11) contains a polynomial of $`t`$ and behaves as $`\mathrm{exp}(\mathrm{\Gamma }\gamma ^6k_x^6ct^7/7)`$ for long times. However, due to the presence of negative coefficients, the decay can be much slower at earlier times, as it can be observed in the upper inset of Fig. 3. In particular one can show that an inflection point is developed at $`\mathrm{tan}\theta =k_x/k_y=\gamma t`$ in the sectors with $`k_xk_y<0`$. Other inflection points come out, even in the other sectors, due to the negative value of $`g^R`$. Such a slow decay can preserve the observation of the maximum of $`C\left(\stackrel{}{K}(t)\right)`$ in the full time correlator (11). In Fig. 3 the behaviour of $`D`$ is shown for three typical values of $`\stackrel{}{k}`$. In case A, $`\stackrel{}{K}(t)`$ does not intersect the volcano of the structure factor so that $`D`$ decays monotonously. When the edge of $`C(\stackrel{}{k})`$ is crossed once or twice, as in cases B and C respectively, $`D`$ is characterized by a corresponding number of peaks. This rich relaxation behavior has never been described before to our knowledge and could be searched for in experiments .
In conclusion, we have presented explicit expressions for the steady-state structure factor and the time correlation function of a Ginzburg-Landau model describing mixtures of oil, water and surfactant in the microemulsion phase under the action of a shear flow. The structure factor exhibits a four peak pattern which implies a non-uniform orientation of interfaces with two preferred directions. The time correlator, reminiscent of the interface structure, shows maxima superimposed on a global decay. All these predictions can be tested experimentally and show that together with the equilibrium behavior also the dynamics of complex fluid is characterized by a very rich phenomenology.
We thank Kenneth Dawson for helpful discussions. F.C. is grateful to M.Cirillo and R. Del Sole for hospitality in the University of Rome. F.C. acknowledges support by the TMR network contract ERBFMRXCT980183 and by MURST(PRIN 97). G.G. acknowledges support by MURST (PRIN 97).
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# Untwisting of a cholesteric elastomer by a mechanical field
## Abstract
A mechanical strain field applied to a monodomain cholesteric elastomer will unwind the helical director distribution. There is an analogy with the classical problem of an electric field applied to a cholesteric liquid crystal, but with important differences. Frank elasticity is of minor importance unless the gel is very weak. The interplay is between director anchoring to the rubber elastic matrix and the external mechanical field. Stretching perpendicular to the helix axis induces the uniform unwound state via the elimination of sharp, pinned twist walls above a critical strain. Unwinding through conical director states occurs when the elastomer is stretched along the helical axis.
PACS numbers: 61.30.-v, 61.41.+e, 78.20.Ek
Monodomain cholesteric elastomers are formed by crosslinking mesogenic chiral polymers in the cholesteric state with a properly formed helical director twist. The subsequent retention of the helical state as an elastic equilibrium is a natural consequence of topological imprinting of textures in the crosslinked network, seen in a number of other elastomers with liquid crystalline order and other microstructure. Recently an interesting aspect of chiral imprinting was established by crosslinking nematic polymers in a chiral state purely induced by a chiral solvent ; on removal of the solvent, the network of chemically achiral nematic chains remains macroscopically cholesteric. Such an imprinting has been envisaged long time ago on phenomenological grounds. It is now important to consider the mechanical possibilities of such solids with a helical microstructure, expecting new transitions and instabilities characteristic of liquid crystalline elastomers. Additionally there are obvious device applications of such materials, which combine all the optical properties of twisted nematic liquids with the remarkable mechanical characteristics of rubbers. There is some experimental evidence that such effects are indeed observable and our hope is that this theoretical work will stimulate more studies in this field.
Consider a monodomain cholesteric elastomer with an ideal helically twisted director $`𝐧_0(z)`$ in the $`xy`$ plane, initially making angle $`\varphi _0=q_0z`$ with the $`x`$ axis, Fig. 1. We shall examine two specific cases of imposed uniaxial extension: (i) the transverse deformation $`\lambda _{\mathrm{xx}}=\lambda `$, in the plane including $`𝐧_0`$, and (ii) the longitudinal deformation along the helix axis $`\lambda _{\mathrm{zz}}=\lambda `$.
The symmetry obvious from Fig. 1 requires that in the case (i) the director remains in the $`xy`$ plane, characterised by the azimuthal angle $`\varphi (z)`$, while in the case (ii) one may expect a conical texture with $`𝐧(z)`$ inclined towards the stretching axis $`z`$ and, therefore, described by two angles $`\theta `$ and $`\varphi `$ (cf. Fig. 3 below). In ordinary liquid cholesterics subjected to, e.g., a magnetic field $`H_z`$, such conical states are not generally seen, preempted by the $`90^\mathrm{o}`$-switching of the helix axis and then untwisting in the “transverse” geometry . We shall see that in elastomers, due to the chiral imprinting, this regime is not possible and the conical director configurations should occur.
An elastic material with a microstructure represented by an independently mobile director orientation is analogous to the Cosserat medium. In the limit of linear elasticity the relative rotation coupling between the director rotation $`𝝎=[𝐧\times \delta 𝐧]`$ and the antisymmetric part of strain, $`\mathrm{\Omega }_\mathrm{i}=ϵ_{\mathrm{ijk}}\epsilon _{\mathrm{jk}}`$,
$$\frac{1}{2}D_1[𝐧\times (𝛀\mathbf{}𝝎)]^2+D_2𝐧\underset{¯}{\underset{¯}{\epsilon }}^{(s)}[𝐧\times (𝛀\mathbf{}𝝎)],$$
(1)
has been first written down phenomenologically by de Gennes , $`\underset{¯}{\underset{¯}{\epsilon }}^{(s)}`$ being the symmetric part of the small strain defined as $`\underset{¯}{\underset{¯}{\epsilon }}=\underset{¯}{\underset{¯}{\lambda }}\underset{¯}{\underset{¯}{\delta }}`$. This symmetry-based expression is only valid for small deformations, having only linear and quadratic terms in the local relative rotation.
The microscopic statistical-mechanical theory of nematic rubber elasticity, e.g. , obtains a generalisation of the classical rubber-elastic energy density in the form of a complete frame-independent expression
$$F=\frac{1}{2}\mu \mathrm{𝖳𝗋}\left(\underset{¯}{\underset{¯}{\mathrm{}}}_0\underset{¯}{\underset{¯}{\lambda }}^T\underset{¯}{\underset{¯}{\mathrm{}}}^1\underset{¯}{\underset{¯}{\lambda }}\right),$$
(2)
plus the constraint of material incompressibility, expressed by the condition $`\mathrm{𝖣𝖾𝗍}\left(\underset{¯}{\underset{¯}{\lambda }}\right)=1`$ on the strain tensor. Apart from the strain tensor, the other entries in the Eq. (2) are $`\underset{¯}{\underset{¯}{\mathrm{}}}_0=\mathrm{}_{}\underset{¯}{\underset{¯}{\delta }}+(\mathrm{}_{}\mathrm{}_{})𝐧_0𝐧_0`$ and $`\underset{¯}{\underset{¯}{\mathrm{}}}^1=(1/\mathrm{}_{})\underset{¯}{\underset{¯}{\delta }}+(1/\mathrm{}_{}1/\mathrm{}_{})𝐧𝐧`$, the reduced shape and inverse shape tensors characterising the Gaussian distribution of nematic polymer chains before and after the distortion $`\underset{¯}{\underset{¯}{\lambda }}`$. The rubber shear modulus $`\mu =n_sk_BT`$ (with $`n_s`$ the number density of network strands, proportional to the crosslink density) is that characteristic of the underlying isotropic rubber and sets the energy scale of distortions. The free energy density Eq. (2) is known to be valid up to large strains and correctly predicts the opto-mechanical responses and the soft elasticity of nematic elastomers. The free energy $`F`$ is a function only of the chain anisotropy $`r=\mathrm{}_{}/\mathrm{}_{}`$, the ratio of the effective step lengths parallel and perpendicular to the director. It is an independently measured parameter accessible from neutron scattering or from spontaneous mechanical distortions on going from the nematic to isotropic phase. Unless there is a large nematic order change induced by $`\underset{¯}{\underset{¯}{\lambda }}`$, the shape $`\underset{¯}{\underset{¯}{\mathrm{}}}`$ is essentially just a rotated version of $`\underset{¯}{\underset{¯}{\mathrm{}}}_0`$, a uniaxial ellipsoid with the long axis (at $`r>1`$) oriented along $`𝐧`$ instead of $`𝐧_0`$.
Embedded in the general expression Eq. (2) is the penalty for local director deviations from the orientation $`𝐧_0`$ imprinted into the network at formation. When no elastic strains are allowed, this elastic energy reduces to
$$F\frac{3}{2}\mu +\frac{1}{2}\mu \frac{(r1)^2}{r}\mathrm{sin}^2\mathrm{\Theta }$$
(3)
where $`\mathrm{\Theta }`$ is the local angle between $`𝐧`$ and $`𝐧_0`$. The elastic penalty for such a deviation, appropriately proportional to the square of chain anisotropy, is the coefficient $`D_1`$ of the de Gennes’ phenomenological expression at small deformations, Eq.(1). This has to be compared with the Frank elastic penalty for director curvature deformations, $`\frac{1}{2}K(𝐧)^2`$. The length scale $`\xi \sqrt{K/\mu }`$ at which the two energy contributions are comparable is usually small: $`\xi 10^8`$m for a typical $`K10^{11}\text{J/m},\mu 10^5\text{J/m}^3`$ and not too small anisotropy, $`r`$. This is rather less than the cholesteric pitch $`p`$, which is a characteristic scale in our problem. Therefore, the anchoring of the director $`𝐧`$ to the rubbery matrix, described by Eq. (2), tends to dominate over Frank effects.
We shall assume that a cholesteric elastomer is locally like a nematic in its elastic response: rubber elasticity is determined on the scale of network crosslink separations (a few nanometers), whereas cholesteric pitches are $`10^3`$ times longer. We can at once see why the chiral structure is stable and how mechanical fields can destabilise it. With no elastic strain, the free energy penalty is $`\frac{1}{2}D_1(\varphi \varphi _0)^2`$ for rotating the director away from its original helical texture $`\varphi _0=q_0z`$. On the other hand, if strains are applied, the rubber can lower its elastic energy Eq. (2) by rotating the director $`𝐧`$ towards the axis of principal extension. This general principle of adjusting the microstructure to minimise the elastic energy is seen in its ultimate form in the effect of soft elasticity , when a stretched nematic rubber may reduce its effective modulus (the slope of a stress-strain curve) to zero by optimising the director rotation and associated shear strains.
Distortions in a cholesteric elastomer cannot be soft, because of elastic compatibility constraints in matching different director and shear modes along the helix. If the director at a position $`z`$ rotates towards the $`x`$-axis, it is known that the elongation $`\lambda _{\mathrm{xx}}`$, contraction $`\lambda _{\mathrm{yy}}`$, and shear $`\lambda _{\mathrm{xy}}`$, are precisely determined by the initial orientation, $`\varphi _0`$, and the rotation from it, if the process is to be soft. The next slab of material, at $`z+dz`$, has the initial orientation $`\varphi _0+q_0dz`$ and a different set of soft strains $`\underset{¯}{\underset{¯}{\lambda }}`$ must arise. Material points at $`y`$ translate to $`\lambda _{\mathrm{yy}}(z)y`$ and $`\lambda _{\mathrm{yy}}(z+dz)y`$ in the two neighbouring slabs along the helix, that is they differ by a relative displacement $`(\lambda _{\mathrm{yy}}/z)dzy`$. There is thus a generated shear $`\lambda _{\mathrm{yz}}=(\lambda _{\mathrm{yy}}/z)y`$ that diverges as the linear $`y`$-dimension of the sample. We accordingly conclude that the transverse contractions are uniform. Such deformations, e.g. $`\lambda _{\mathrm{zz}}`$ and $`\lambda _{\mathrm{yy}}`$, are equal in the first approximation, in spite of the apparent anisotropy along the pitch axis $`z`$ (each deformation generates the same energy $`k_BT`$ per chain independently of its alignment with respect to the director) and are thus both equal to $`1/\sqrt{\lambda }`$.
(i) Transverse elongation $`\lambda _{\mathrm{xx}}=\lambda `$. We consider the strain tensor in the following form:
$$\underset{¯}{\underset{¯}{\lambda }}=\left(\begin{array}{ccc}\lambda & 0& 0\\ 0& 1/\sqrt{\lambda }& 0\\ 0& 0& 1/\sqrt{\lambda }\end{array}\right).$$
(4)
Although one expects the director rotation in the azimuthal plane (cf. Fig. 1), there are no associated shear strains. Such shears, $`\lambda _{\mathrm{xy}}(z)`$ and $`\lambda _{\mathrm{yx}}(z)`$, would both lead to elastic compatibility problems (e.g. a generated shear $`\lambda _{\mathrm{xz}}y`$, the sample dimension) and we assume they are suppressed. The shears $`\lambda _{\mathrm{xz}}(z)`$ and $`\lambda _{\mathrm{xz}}(z)`$ are not subject to compatibility requirements. However, they should not appear on symmetry grounds, which is easily confirmed by direct minimisation. Now $`𝐧_0=\{\mathrm{cos}\varphi _0,\mathrm{sin}\varphi _0,\mathrm{\hspace{0.17em}0}\}`$ and the rotated director after deformation is $`𝐧=\{\mathrm{cos}\varphi ,\mathrm{sin}\varphi ,\mathrm{\hspace{0.17em}0}\}`$. Note that the helix is $`\varphi _0=q_0z`$ in the initial undistorted material. After deformation, because of the affine contraction $`\lambda _{\mathrm{zz}}=1/\sqrt{\lambda }`$, the material frame shrinks: $`zz/\sqrt{\lambda }`$. This has an effect of affine contraction of the helical pitch so that $`\stackrel{~}{q}=\sqrt{\lambda }q_0`$ in all expressions below. With the $`\underset{¯}{\underset{¯}{\mathrm{}}}_0`$ and $`\underset{¯}{\underset{¯}{\mathrm{}}}`$ implied by these $`𝐧_0`$ and $`𝐧`$, the free energy density Eq. (2) yields
$`F_{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mu (\lambda ^2+{\displaystyle \frac{2}{\lambda }}+{\displaystyle \frac{r1}{r}}[\lambda ^2(rc_0^2s^2c^2s_0^2)`$ (5)
$`+{\displaystyle \frac{1}{\lambda }}(rc^2s_0^2s^2c_0^2)2\sqrt{\lambda }(r1)s_0c_0sc]),`$ (6)
where $`c_0`$ and $`s_0`$ are shorthand for $`\mathrm{cos}\varphi _0`$ and $`\mathrm{sin}\varphi _0`$; analogously, $`c`$ and $`s`$ stand for $`\mathrm{cos}\varphi `$ and $`\mathrm{sin}\varphi `$. The appearance of terms linear and quadratic in $`\varphi `$ (or rather $`\mathrm{sin}\varphi `$ because all values of the azimuthal angle will be found along the cholesteric helix) indicate that rotations can always lower the energy for $`\lambda 1`$. Minimization of $`F_{}`$ with respect to $`\varphi `$ results in the expression for the local director angle $`\varphi (z)`$ at a given imposed extension $`\lambda `$, depending on the phase of cholesteric helix:
$$\mathrm{tan}2\varphi =\frac{2\lambda ^{3/2}(r1)\mathrm{sin}\mathrm{\hspace{0.17em}2}\stackrel{~}{q}z}{(r1)(\lambda ^3+1)\mathrm{cos}\mathrm{\hspace{0.17em}2}\stackrel{~}{q}z+(r+1)(\lambda ^31)}$$
(7)
Initially, all directors at $`0<\stackrel{~}{q}z<\pi /2`$ are induced to rotate “backward” towards $`\varphi =0`$, and all directors at $`\pi /2<\stackrel{~}{q}z<0`$ rotate “forward” towards $`\varphi =\pi `$, as the imposed deformation $`\lambda `$ increases, see Fig. 2. Although $`\varphi =0`$ and $`\pi `$ describe equivalent directors, the twist wall between these two states becomes more and more sharp. Due to the helix imprinting, the orientations $`\varphi =0`$ at $`\stackrel{~}{q}z=0`$ and $`\varphi =\pi `$ at $`\stackrel{~}{q}z=\pi `$ are pinned, as is the middle-point of the twist wall at $`\stackrel{~}{q}z=\pi /2`$. As a result, no change of the helical pitch can occur. This is in contrast with cholesteric liquid crystals, where in a classical problem of helix unwinding by electric or magnetic field one finds an increase in cholesteric pitch along with the coarsening of the helix.
Examining the Eq.(7) one finds that the denominator changes sign and remains negative in the region of the twist wall, centered at $`\stackrel{~}{q}z=\pi /2`$ between the values $`\varphi =\pi /4`$ and $`3\pi /4`$. The width of such a wall is
$$w\frac{2}{q_0}\sqrt{\frac{r\lambda ^3}{(r1)(\lambda ^3+1)}}.$$
As the increasing applied strain reaches a critical value $`\lambda _c=r^{1/3}`$, the wall width $`w0`$ and the discontinuous transition occurs. The director in the mid-point of the wall breaks away from the pinning and jumps from $`\varphi =\pi /2`$ to $`\varphi =0`$, along the strain axis, thus removing the topologically constrained twist wall. From this point there is no barrier for director rotation towards the final uniform orientation with $`\varphi =0`$, as the last two curves in Fig. 2 indicate. The bifurcation occurs between the family of $`\varphi (z)`$ curves at $`\lambda <\lambda _c`$ and those at $`\lambda >\lambda _c`$. They are separated by the critical curve $`\varphi _c`$ given by $`\mathrm{tan}2\varphi _c=2\left(\frac{\sqrt{r}}{r+1}\right)\mathrm{tan}\stackrel{~}{q}z`$.
A discontinuous director jump at a critical strain has been predicted, and indeed observed in nematic elastomers stretched at $`90^\mathrm{o}`$ to their initial director $`𝐧_0`$ . It was also found that even small deviations from the exact perpendicular geometry, or the possibility of soft shears in stripe domains lead to a continuous director rotation. However, in the present cholesteric case the shears are prohibited by compatibility constraints and one always finds an exact phase angle $`\varphi =\pi /2`$ along the helix – where the center of narrowing twist wall becomes pinned from both sides. It is this point that experiences a discontinuous jump in order to allow the system to proceed to its final uniform equilibrium state.
(ii) Stretching along the pitch axis $`\lambda _{\mathrm{zz}}=\lambda `$. We now take (cf. Fig. 1)
$$\underset{¯}{\underset{¯}{\lambda }}=\left(\begin{array}{ccc}1/\sqrt{\lambda }& 0& \lambda _{\mathrm{xz}}\\ 0& 1/\sqrt{\lambda }& \lambda _{\mathrm{yz}}\\ 0& 0& \lambda \end{array}\right)$$
(8)
No compatibility problem with shears $`\lambda _{\mathrm{xz}}(z)`$ and $`\lambda _{\mathrm{yz}}(z)`$ arises from their variation with $`z`$ along the helical pitch. By contrast, their conjugate strains $`\lambda _{\mathrm{zx}}`$ and $`\lambda _{\mathrm{zy}}`$, which would also have to vary with $`z`$, would lead to a serious compatibility mismatch, e.g. $`\lambda _{\mathrm{zx}}/z=\lambda _{\mathrm{zz}}/x`$. We therefore assume $`\lambda _{\mathrm{zx}}`$ and $`\lambda _{\mathrm{zy}}`$ are suppressed even though in other settings these are the generators of soft elastic response.
In this geometry one expects the director rotation out of the azimuthal $`xy`$ plane, see Fig. 3(a). The initial director is, as before, $`𝐧_0=\{\mathrm{cos}q_0z,\mathrm{sin}q_0z,\mathrm{\hspace{0.17em}0}\}`$, while after deformation the rotated director is $`𝐧=\{\mathrm{cos}\theta \mathrm{cos}\stackrel{~}{q}z,\mathrm{cos}\theta \mathrm{sin}\stackrel{~}{q}z,\mathrm{sin}\theta \}`$. As in the case (i), all physical dimensions in the deformed sample are scaled by the affine strain. In particular, here $`z\lambda z`$, resulting in the corresponding expansion of the cholesteric pitch: $`\stackrel{~}{q}=q_0/\lambda `$. With the $`\underset{¯}{\underset{¯}{\mathrm{}}}_0`$ and $`\underset{¯}{\underset{¯}{\mathrm{}}}`$ defined by the axes $`𝐧_0`$ and $`𝐧`$, the free energy density Eq. (2) now becomes a function of three variables: the director tilt angle $`\theta `$ and the two shear strains $`\lambda _{\mathrm{xz}}(z)`$ and $`\lambda _{\mathrm{yz}}(z)`$ (we continue to neglect the effects of director gradients and Frank elasticity). Algebraic minimisation over these components of strain tensor is not difficult and results in
$$\left(\begin{array}{c}\lambda _{\mathrm{xz}}\\ \lambda _{\mathrm{yz}}\end{array}\right)=\lambda \frac{(r1)\mathrm{sin}\mathrm{\hspace{0.17em}2}\theta }{r+1(r1)\mathrm{cos}\mathrm{\hspace{0.17em}2}\theta }\left(\begin{array}{c}\mathrm{cos}\stackrel{~}{q}z\\ \mathrm{sin}\stackrel{~}{q}z\end{array}\right),$$
(9)
in phase with the azimuthal angle along the helical pitch. Eq. (9) describes small distortions in the $`xy`$ plane, perpendicular to the helix axis, rotating following the initial orientation $`𝐧_0`$. On substitution of these optimal shears back into the free energy density one obtains
$$F_{}=\frac{1}{2}\mu \left(\frac{\lambda ^2}{1+(r1)\mathrm{sin}^2\theta }+\frac{2+(r1)\mathrm{sin}^2\theta }{\lambda }\right)$$
(10)
$`F_{}`$ expands at small tilt angle $`\theta `$ as
$$F_{}\frac{1}{2}\mu (\lambda ^2+2/\lambda )\frac{1}{2}\mu \theta ^2(r1)(\lambda ^21/\lambda ),$$
(11)
that is the director starts to rotate down to define a cone of semiangle $`\pi /2\theta `$ immediately as the strain $`\lambda >1`$ is imposed. The equilibrium director tilt is obtained by minimisation of the full free energy density $`F_{}(\theta )`$:
$$\mathrm{sin}^2\theta =\frac{\lambda ^{3/2}1}{r1};\theta =\mathrm{arcsin}\sqrt{\frac{\lambda ^{3/2}1}{r1}}.$$
(12)
The director rotation starts and ends in a characteristically singular fashion Fig. 3(b) (reminiscent of the universal opto-mechanical response seen in nematic elastomers ). The rotation is complete with the director aligned along the extension axis ($`\theta =\pi /2`$) at $`\lambda =r^{2/3}`$ which, for some elastomers, can be a very large extension.
In contrast to conventional cholesteric liquid crystals, we have altogether ignored effects of Frank elastic energy. The most compelling evidence for this is the very stability of the imprinted helical state in the face of the Frank penalty $`\frac{1}{2}K_2q_0^2`$. The argument for this relies upon the great difference in characteristic length scales, the elastomer penetration depth, more accurately expressed as $`\xi \frac{1}{r1}\sqrt{K/\mu }`$ \[cf. Eq. (3)\], and the director modulation wavelength estimated by the helical pitch $`p=\pi /q_0\xi `$. There are two possibilities to alter this inequality – by increasing the penetration depth $`\xi `$ (either by making a weaker gel, or a less anisotropic one), or by locally increasing the director gradient (for instance, in the ever narrowing twist wall, Fig. 2).
One can estimate how weak a gel must be for the Frank elasticity to intervene in our analysis. When $`\xi p`$, for example with a pitch $`p4\times 10^7`$m, then a rubber modulus of only $`\mu 60\text{J/m}^3`$ is required (assuming $`[r1]1`$). Nematic elastomers typically have $`\mu 10^310^5\text{J/m}^3`$ and their cholesteric analogues would clearly find Frank-elastic effects minor. However, an elastomer with a reasonable $`\mu 10^3\text{J/m}^3`$ would feel the director gradients when its polymer chain anisotropy becomes as low as $`r=\mathrm{}_{}/\mathrm{}_{}1.25`$. Such a value is easily reached in side-chain liquid crystal polymers, especially near the clearing point .
Another interesting test of the role of Frank elasticity is in the twist wall described in our case (i). The width of the wall decreases to zero, and the Frank energy density grows, being maximal at the centre of the wall. There it is $`\frac{1}{2}K_2q_0^2\lambda ^3(r1)^2/(r\lambda ^3)`$ and diverges at the critical strain $`\lambda _c`$. Therefore, the local analysis of Eqs. (5)-(7) is only valid outside the region of strain $`\mathrm{\Delta }\lambda \left(\frac{r^21}{3r^{2/3}}\right)(q_0\xi )^2`$ around $`\lambda _c`$. In a typical hard nematic rubber this is a narrow region of $`\mathrm{\Delta }\lambda 0.01`$, but in a weak gel with low chain anisotropy it may become more substantial. Moreover, the finite width of the twist wall, demanded by the Frank gradient energy, raises the question of topological mechanism for eliminating the twist stored in the cholesteric helix, perhaps by a disclination loop expansion.
To summarise, we have predicted a qualitatively new response of an elastomer with chiral cholesteric microstructure to applied fields that is different from classical cholesteric liquids. Likewise, the chiral imprinting and its modification by elastic fields is a new effect in rubbers and solids. One could envisage tuning these effects by the use of solvents (with or without chiral power) and by other fields affecting the director, for instance electric.
YM is grateful to St John’s College, Cambridge for a research fellowship and RBM acknowledges support by the NSF, through grant DMR-9974388.
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# Dynamical Chiral Symmetry Breaking on the Light Front. II. The Nambu–Jona-Lasinio Model
## I Introduction
Our expectation that light-front (LF) formalism enables us to relate QCD directly to the constituent quark model at field-theoretic level , seriously requires a full understanding of dynamical chiral symmetry breaking (D$`\chi `$SB) on the LF. Central to this issue is, first of all, a well-known problem of how to reconcile a LF “trivial” vacuum with a chirally broken vacuum having a nonzero fermion condensate. The secondary problem is to determine the property of “LF chiral transformation” which is defined differently from the usual one. The most surprising fact of the LF chiral transformation is that it is an exact symmetry even for a massive free fermion .
In the present paper, we discuss this issue within the Nambu–Jona-Lasinio (NJL) model which is a typical example of D$`\chi `$SB. Previously we consider the same problem from different point of view . Our interest was in describing D$`\chi `$SB of the NJL model, but we actually took a roundabout way in order to apply an idea which works well for spontaneous symmetry breaking of a scalar model, to a fermionic theory. We know that the longitudinal zero modes of scalar fields are responsible for describing spontaneous symmetry breaking on the LF. Indeed, it is achieved by solving the “zero-mode constraints” (i.e., constraint equations for the longitudinal zero modes) nonperturbatively . The zero-mode constraint appears in the discretized light-cone quantization (DLCQ) approach , where we set periodic boundary conditions for scalars in the longitudinal direction with finite extension. Of course the NJL model has no scalar fields as fundamental degrees of freedom, but we overcame the situation by considering the chiral Yukawa model. This model shows D$`\chi `$SB in large $`N`$ limit ($`N`$ is the number of fermions) and goes to the NJL model in infinitely heavy mass limit of scalar and pseudoscalar bosons. We showed that the zero-mode constraint of the scalar field correctly produces a gap equation for a chiral condensate and calculated masses of the scalar and pseudoscalar bosons from poles of their propagators. Therefore, in Ref. , we succeeded in describing indirectly the chiral symmetry breaking of the NJL model on the LF.
Since the very essence of the previous analysis was the existence of scalar fields, one may ask a question: How can one formulate D$`\chi `$SB without scalars? In order to answer the question, we treat the NJL model without introducing auxiliary field. An important key was already shown in Ref. . It was found that the “fermionic constraint” plays the same role as that of the zero-mode constraint. Splitting the fermion field as $`\mathrm{\Psi }=\psi _++\psi _{},\psi _\pm =\mathrm{\Lambda }_\pm \mathrm{\Psi }`$ by using projectors $`\mathrm{\Lambda }_\pm =\gamma ^{}\gamma ^\pm /2`$, we easily find that the “bad” component $`\psi _{}`$ is a dependent variable and subject to a constraint equation called the “fermionic constraint.” In the LF NJL model, the fermionic constraint is very complicated and it is difficult to solve it as an operator equation. However, we will see that to solve this equation is a crucial step for describing the broken phase and will find a close parallel between the fermionic constraint for D$`\chi `$SB and the zero-mode constraint for spontaneous symmetry breaking in scalar models. Although such special importance of the fermionic constraint might be restricted only to the LF NJL model and therefore most of the analysis might be model dependent, but what we are eventually interested in is the physics consequences of the chiral symmetry breaking. And of course we cannot reach the chiral symmetry breaking in QCD unless we understand the simpler and typical example of this phenomenon. Therefore the importance of our analysis is evident.
Let us comment on other attempts for the NJL model on the LF. First of all, Heinzl et al. treated the model within the mean-field approximation and insisted delicacy of the infrared cutoff to obtain a chiral condensate. The meaning and necessity of such cutoff scheme was clarified in Ref. . As mentioned above, an observation that a gap equation for a chiral condensate emerges from the fermionic constraint was first pointed out by one of the authors . The light-cone (LC) wavefunction of a pionic state was calculated through the LC projection of the Bethe-Salpeter amplitude which was derived in the equal-time quantization . Bentz et al. introduced the auxiliary fields to fermion bilinears and solved the constraint equations for them by $`1/N`$ expansion . They obtained “effective” Lagrangian for the broken phase and discussed the structure function of the pionic state. With all these studies, however, there still remains many unknowns concerning basic problems. Especially, we still do not understand well the LF chiral transformation itself. To what extent is it different from the usual chiral symmetry? How is the chiral symmetry breaking realized on the LF? These fundamental problems will be resolved in the present paper.
The paper is organized as follows. The rest of this section is devoted to introduction of the NJL model and our notation. In Sec. II, we discuss the complexity of the fermionic constraint in great detail. We explicitly solve the fermionic constraint in classical treatment and investigate properties of LF chiral transformation. In Section III, we solve the fermionic constraint in quantum theory by the $`1/N`$ expansion. Here we introduce the boson expansion method in order to solve the bilocal fermionic constraint with systematic $`1/N`$ expansion. We see emergence of the gap equation for the chiral condensate from the fermionic constraint. We obtain the Hamiltonian with respect to the (bilocal) bosons which is introduced by the boson expansion method. In Section IV, some physics consequences of the chiral symmetry breaking are discussed. First of all, we see how the chiral symmetry breaking is realized in the LF formalism. We discuss unusual chiral transformation of fields and nonconservation of the light-front chiral charge. Secondly, we construct the bound-state equation for mesonic states and solve it for scalar and pseudoscalar mesons. Thirdly, we derive the partially conserving axial current (PCAC) relation. Summary and conclusion are given in the last section. Miscellaneous topics with detailed calculation are presented in Appendices.
Before ending this section, let us fix our model and notation. Since the primary purpose of our paper is to study basic properties of the LF chiral symmetry, we consider only one flavor case for simplicity. Thus the model we discuss is
$$=\overline{\mathrm{\Psi }}(i\text{/}m_0)\mathrm{\Psi }+\frac{g^2}{2}\left[(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2+(\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi })^2\right].$$
(1)
Here $`\mathrm{\Psi }_\alpha ^a(x)(a=1,\mathrm{},N)`$ is a four component spinor with “color” internal symmetry $`U(N)`$, which has been introduced so that we can use the $`1/N`$ expansion as a nonperturbative technique. We always work with a nonzero bare mass $`m_00`$. The primary reason is that the Hamiltonian with a massless fermion is plagued with a troublesome situation in 1+1 dimensions: As we will see, if we set $`m_0=0`$ from the beginning, the canonical LF Hamiltonian $`P^{}`$ of the Gross-Neveu model vanishes at all. Even in 3+1 dimensions, we will see that absence of the bare mass term causes an inconsistency of the results. The secondary reason is to avoid massless particles which can move in parallel with $`x^+=`$constant surface. The difficulty of describing massless particles is intimately connected with the fact that on the LF, the (massless) Nambu-Goldstone boson becomes physically meaningful only when we first include explicit breaking term and then take the vanishing limit of it . The same situation was observed in the chiral Yukawa model .
In practical calculation, it is convenient to introduce the two-component representation for the gamma matrices so that the projectors $`\mathrm{\Lambda }_\pm `$ are expressed as
$$\mathrm{\Lambda }_+=\left(\begin{array}{cc}\mathrm{𝟏}& 0\\ 0& 0\end{array}\right),\mathrm{\Lambda }_{}=\left(\begin{array}{cc}0& 0\\ 0& \mathrm{𝟏}\end{array}\right).$$
(2)
Then, the projected fermions have only upper or lower components:
$$\psi _+=\mathrm{\Lambda }_+\mathrm{\Psi }2^{\frac{1}{4}}\left(\begin{array}{c}\psi \\ 0\end{array}\right),\psi _{}=\mathrm{\Lambda }_{}\mathrm{\Psi }2^{\frac{1}{4}}\left(\begin{array}{c}0\\ \chi \end{array}\right),$$
(3)
where we defined two-component spinors $`\psi `$ and $`\chi `$. Among various representations which satisfy Eq. $`(\text{2})`$, we choose a representation having a similar structure to the chiral representation in 1+1 dimensions, i.e., $`\gamma ^0=\sigma ^1,\gamma ^1=i\sigma ^2,\gamma _5=\gamma ^0\gamma ^1=\sigma ^3`$. Explicitly, they are
$$\gamma ^0=\left(\begin{array}{cc}0& \mathrm{𝟏}\\ \mathrm{𝟏}& 0\end{array}\right),\gamma ^3=\left(\begin{array}{cc}0& \mathrm{𝟏}\\ \mathrm{𝟏}& 0\end{array}\right),\gamma ^i=\left(\begin{array}{cc}i\sigma ^i& 0\\ 0& i\sigma ^i\end{array}\right),$$
(4)
for $`i=1,2`$ and
$$\gamma _5=\left(\begin{array}{cc}\sigma ^3& 0\\ 0& \sigma ^3\end{array}\right).$$
(5)
Results of the chiral Gross-Neveu model in 1+1 dimensions can be easily obtained if we make a replacement for the Pauli matrices $`\sigma ^31`$ and $`\sigma ^i0`$, and regard $`\psi `$ and $`\chi `$ as one component spinors. The explicit form of the Lagrangian in the two-component representation is given in Appendix A.
In the previous work , we made the longitudinal direction finite in order to carefully treat the longitudinal zero modes of the scalar fields. However, in the present analysis, we work in an infinite longitudinal space. There is no need of introducing finite $`x^{}`$. When we take the inverse of $`_{}`$, we need to specify the boundary conditions. We here follow the conventional antiperiodic boundary conditionFor a scalar field, antiperiodic boundary condition in infinite longitudinal space leads to inconsistency . However, fermionic fields are free from such troubles. which is standard for free fermions:
$`\psi _a(x^{}=\mathrm{},x_{}^i)`$ $`=`$ $`\psi _a(x^{}=\mathrm{},x_{}^i),`$ (6)
$`\chi _a(x^{}=\mathrm{},x_{}^i)`$ $`=`$ $`\chi _a(x^{}=\mathrm{},x_{}^i).`$ (7)
So our results must always have a smooth free field limit.
## II Complexity of the fermionic constraint and the chiral symmetry
It is highly complicated structure of the fermionic constraint that makes the analysis of the LF NJL model difficult. However, we cannot know anything about LF chiral symmetry unless we confront with this complexity. Therefore in this section, we investigate the fermionic constraint in great detail. First of all, we classically solve the fermionic constraint. Using the explicit solution, we then discuss properties of the LF chiral transformation. Especially we show that LF chiral transformation is no longer an exact symmetry when $`m_00`$. Finally we consider the implication of the fermionic constraint in quantum theory.
### A Fermionic constraint and its classical solution
The fermionic constraint is immediately obtained as the Euler-Lagrange equation for $`\chi `$:
$$i_{}\chi _a=\frac{1}{\sqrt{2}}\left(\sigma ^i_i+m_0\right)\psi _a\frac{g^2}{2}\left\{\psi _a\left(\psi _b^{}\chi _b+\chi _b^{}\psi _b\right)+\sigma ^3\psi _a\left(\psi _b^{}\sigma ^3\chi _b\chi _b^{}\sigma ^3\psi _b\right)\right\},$$
(1)
where summation over color and spinor indices are implied. If we want to solve this equation as an operator equation in quantum theory, we need a commutation relation between $`\chi `$ with $`\psi `$ which must be given by the Dirac bracket. Since the anticommutator $`\{\chi ,\psi \}`$ is very complicated, it seems almost hopeless to find an exact quantum solution of it. However, in a classical theory where we treat all the variables just as Grassmann numbers, the equation becomes tractable and it is not difficult to solve it. Indeed, the exact solution with antiperiodic boundary condition is given by (see Appendix B for more details)
$$\left(\begin{array}{c}\chi _{1a}(x)\\ \chi _{2a}^{}(x)\end{array}\right)=\frac{1}{\sqrt{2}}_{\mathrm{}}^{\mathrm{}}𝑑y^{}G_{ab}(x^{},y^{},x_{})\left(\begin{array}{c}m_0\psi _{1b}(y^{})_z\psi _{2b}(y^{})\\ _z\psi _{1b}^{}(y^{})+m_0\psi _{2b}^{}(y^{})\end{array}\right),$$
(2)
where $`_z=_1i_2`$ and the “Green function” $`G_{ab}(x^{},y^{},x_{})`$ is
$`G_{ab}(x^{},y^{},x_{})=G^{(0)}(x)\left[{\displaystyle \frac{1}{2i}}ϵ(x^{}y^{})+C\right]G^{(0)1}(y),`$ (3)
$`G^{(0)}(x)=\mathrm{P}e^{ig^2_{\mathrm{}}^x^{}𝒜(y^{})𝑑y^{}},`$ (4)
$`𝒜_{ijab}=\left(\begin{array}{cc}\psi _{1a}\psi _{1b}^{}& \psi _{1a}\psi _{2b}\\ \psi _{2a}^{}\psi _{1b}^{}& \psi _{2a}^{}\psi _{2b}\end{array}\right).`$ (5)
The integral constant $`C`$ is determined so that the solution satisfies the antiperiodic boundary condition. In Eq. (4), P stands for the path-ordered product. Note that we easily derive the solution for the chiral Gross-Neveu model by extracting the 1-1 component of $`𝒜`$ and neglecting $`_z`$. The result is equivalent to the solution of the Thirring model obtained by Domokos . And also, if we take the free fermion limit $`g^20`$, we of course recover the free solution due to $`G^{(0)}(x)1`$ and $`G(x^{},y^{})ϵ(x^{}y^{})/2i`$.
### B Chiral symmetry on the light front
Since the bad component $`\chi `$ is a constrained variable in the LF formalism, we impose the chiral transformation only on the good component $`\psi _+e^{i\theta \gamma _5}\psi _+`$ or in the two-component representation \[see Eq. (5)\]
$$\psi e^{i\theta \sigma ^3}\psi .$$
(6)
Now we have completely solved the fermionic constraint for $`\chi `$, we can explicitly demonstrate its transformation property under the LF chiral transformation. However, before discussing the NJL model, it will be instructive to remind you of the LF chiral symmetry in the free massive fermion.
As we mentioned before, the massive free fermion is chiral invariant under the transformation (6). Let us see this fact directly in the Lagrangian even though it is a little lengthy. It is convenient to separate the solution of the fermionic constraint $`\chi =(\sqrt{2}i_{})^1\left(\sigma ^i_i+m_0\right)\psi `$ into mass-independent and dependent parts $`\chi =\chi ^{(0)}+\chi ^{(m)}`$ as
$$\chi ^{(0)}=\frac{1}{\sqrt{2}}\sigma ^i_i\frac{1}{i_{}}\psi ,\chi ^{(m)}=\frac{m_0}{\sqrt{2}}\frac{1}{i_{}}\psi .$$
Note that there is a relation between $`\chi ^{(0)}`$ and $`\chi ^{(m)}`$:
$$\sigma ^i_i\chi ^{(m)}+m_0\chi ^{(0)}=0.$$
(7)
As a result of the LF chiral transformation (6), we find
$`\chi ^{(0)}e^{i\theta \sigma ^3}\chi ^{(0)},`$ (8)
$`\chi ^{(m)}e^{i\theta \sigma ^3}\chi ^{(m)}.`$ (9)
The free fermion Lagrangian is compactly expressed as $`_{\mathrm{free}}=\psi ^{}\omega _{\mathrm{EOM}}+\chi ^{}\omega _{\mathrm{FC}},`$ where $`\omega _{\mathrm{EOM}}=i_+\psi \frac{1}{\sqrt{2}}(\sigma ^i_i+m_0)\chi =0`$ is the equation of motion for $`\psi `$ and $`\omega _{\mathrm{FC}}=i_{}\chi \frac{1}{\sqrt{2}}(\sigma ^i_i+m_0)\psi =0`$ is the fermionic constraint. The second term is zero and is invariant under the LF chiral transformation. Now substituting $`\chi =\chi ^{(0)}+\chi ^{(m)}`$ into the Lagrangian, the first term is decomposed into apparently invariant and (seemingly) non-invariant terms
$$\psi ^{}\omega _{\mathrm{EOM}}=\psi ^{}\left[i_+\psi \frac{1}{\sqrt{2}}\left(\sigma ^i_i\chi ^{(0)}+m_0\chi ^{(m)}\right)\right]+\psi ^{}\left[\frac{1}{\sqrt{2}}\left(\sigma ^i_i\chi ^{(m)}+m_0\chi ^{(0)}\right)\right].$$
The first term consists of the $`m_0`$-independent term and quadratically dependent term $`𝒪(m_0^2)`$, while the second term linearly depends on $`m_0`$. The $`𝒪(m_0)`$ term changes under the chiral transformation, but due to the relation (7), it eventually vanishes and therefore the Lagrangian is invariant even if there is a mass term. As a result, we have a conserved Noether current
$`j_{5\mathrm{F}\mathrm{r}\mathrm{e}\mathrm{e}}^\mu =\overline{\mathrm{\Psi }}\gamma ^\mu \gamma _5\mathrm{\Psi }m_0\overline{\mathrm{\Psi }}\gamma ^\mu \gamma _5{\displaystyle \frac{1}{i_{}}}\gamma ^+\psi _+,`$ (10)
$`_\mu j_{5\mathrm{F}\mathrm{r}\mathrm{e}\mathrm{e}}^\mu =0,`$ (11)
which of course reduces to the usual current in the massless limit.
Now let us consider the NJL model. Decomposition of $`\chi `$ is straightforward:
$`\left(\begin{array}{c}\chi _{1a}^{(0)}\\ \chi _{2a}^{(0)}\end{array}\right)={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑y^{}G_{ab}(x^{},y^{},x_{})\left(\begin{array}{c}_z\psi _{2b}\\ _z\psi _{1b}^{}\end{array}\right),`$ (12)
$`\left(\begin{array}{c}\chi _{1a}^{(m)}\\ \chi _{2a}^{(m)}\end{array}\right)={\displaystyle \frac{m_0}{\sqrt{2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑y^{}G_{ab}(x^{},y^{},x_{})\left(\begin{array}{c}\psi _{1b}\\ \psi _{2b}^{}\end{array}\right).`$ (13)
Since the matrix $`𝒜`$, and thus $`G_{ab}(x,y)`$ is invariant under the transformation (6), it is easy to find that $`\chi ^{(0)}`$ and $`\chi ^{(m)}`$ transform as Eqs. (8) and (9). Therefore, if $`m_0=0`$, the LF chiral transformation (6) is equivalent to the usual chiral transformation. The chiral current and the chiral charge are given by
$`j_5^\mu =\overline{\mathrm{\Psi }}\gamma ^\mu \gamma _5\mathrm{\Psi },`$ (14)
$`Q_5^{\mathrm{LF}}={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x^{}d^2x_{}j_5^+(x)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x^{}d^2x_{}\psi ^{}\sigma ^3\psi .`$ (15)
How about the massive case? As we explicitly showed above, the mass term does not prevent chiral symmetry in the free fermion case. We must bare in mind such possibility even in the NJL model. Thus it is worth while to check whether the massive NJL model is invariant under the LF chiral transformation. To see this, it is convenient to treat the Hermite Lagrangian
$$_{\mathrm{Hermite}}=\frac{1}{2}i\psi ^{}\underset{+}{\overset{}{}}\psi \frac{1}{2\sqrt{2}}\left[\left(\psi ^{}\sigma ^i_i\chi +_i\chi ^{}\sigma ^i\psi \right)+m_0\left(\psi ^{}\chi +\chi ^{}\psi \right)\right].$$
(16)
Note that this is equivalent to the free Lagrangian except that $`\chi `$ is a solution of Eq. (1). Now the apparently non-invariant term is a term linearly depending on $`m_0`$:
$$\frac{1}{2\sqrt{2}}\psi ^{}\left(\sigma ^i_i\chi ^{(m)}+m_0\chi ^{(0)}\right)+\mathrm{H}.\mathrm{c}..$$
(17)
In the massive free fermion case, we had the same term but it eventually vanished due to Eq. (7). However, in the NJL model, it is evident from Eqs. (12) and (13) such relation does not hold because $`G`$ depends on $`x_{}`$. Therefore we have verified that the massive NJL model is not invariant under the LF chiral transformation. If and only if $`m_0=0`$, the LF chiral transformation is the symmetry of the NJL model and equivalent to the usual chiral transformation. This is of course not a surprising result but must be checked explicitly. Anyway, we do not stick to this problem anymore.
Irrespective of whether we have a mass term or not, we always use the definition for the chiral current Eq. (14) which was derived for the massless fermion. In the massless case, it is of course a conserved current $`_\mu j_5^\mu =0`$, while in the massive case, a usual relation holds
$$_\mu j_5^\mu =2m_0\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi },$$
(18)
which is derived by using the Euler-Lagrange equation for the massive fermion. Equation (18) is used when we discuss nonconservation of the chiral charge and the PCAC relation in Sec. IV.
### C Implication of the fermionic constraint
So far we treated the fermionic constraint in classical theory and obtained the exact solution (2). However, this solution does not give a nonzero condensate and the resulting Hamiltonian does not describe the broken phase. The situation is very similar to the previous analysis of the chiral Yukawa model . The chiral Yukawa model in the DLCQ approach has three constraint equations. We solved them in classical theory but we could not find any way to describe the broken phase with the classical solutions. What we finally found is that it is very important to treat the constraint equations, especially the zero-mode constraint for a scalar field, nonperturbatively in quantum theory. This fact is true of our present case. To obtain a nonzero condensate, we must treat the fermionic constraint as an operator equation and solve it with some nonperturbative method.
To strengthen this, let us briefly overview the procedure in the previous work . In the chiral Yukawa model, we have three dependent variables: Two of them are the longitudinal zero-modes of scalar and pseudoscalar fields $`\sigma _0(x_{})=(1/2L)_L^L𝑑x^{}\sigma (x),\pi _0(x_{})=(1/2L)_L^L𝑑x^{}\pi (x)`$ where $`L`$ is an extension of the longitudinal direction $`x^{}[L,L]`$, and the rest is the bad component of a fermion $`\psi _{}(x)`$. So there are three constraints:
$$\left(\frac{\mu ^2}{\lambda }_{}^2\right)\left(\begin{array}{c}\sigma _0\\ \pi _0\end{array}\right)=\frac{\mu ^2}{N}\frac{1}{\sqrt{2}}_L^L\frac{dx^{}}{2L}\left[\psi _+^a\left(\begin{array}{c}1\\ i\gamma _5\end{array}\right)\gamma ^{}\psi _{}^a+\psi _{}^a\left(\begin{array}{c}1\\ i\gamma _5\end{array}\right)\gamma ^+\psi _+^a\right],$$
(19)
$$2i_{}\psi _{}^a=\left(i\gamma _{}^i_i+m_0+\sigma i\pi \gamma _5\right)\gamma ^+\psi _+^a,$$
(20)
where $`\lambda =g^2N`$ in the present notation and $`\mu `$ is a dimensionless parameter which controls the scalar and pseudoscalar masses. In the infinitely heavy mass limit, $`\mu \mathrm{}`$, we recover the NJL model. The procedure of Ref. is as follows: First, we formally solved the fermionic constraint (20) and substitute the solution into the zero-mode constraints (19). Second, we solved the zero-mode constraints by $`1/N`$ expansion with a fixed operator ordering and found that the leading order of the scalar zero-mode constraint can be identified with the gap equation. Selecting a nontrivial solution of the gap equation, we again substitute it back to the fermionic constraint. Then we obtain the final expression for the bad component $`\psi _{}`$ in terms of independent variables. Thus we solved three coupled equations step by step. On the other hand, we have only one constraint equation. The procedure in the chiral Yukawa model suggests that we will have to do almost the same procedure at once when we solve the fermionic constraint. Note that just the same as in the chiral Yukawa model, a perturbative solution cannot reach the broken phase even in quantum theory. Therefore, we naturally expect that solving the fermionic constraint (1) in quantum theory using some nonperturbative method is necessary for describing the chiral symmetry breaking .
## III Solving the fermionic constraint by $`1/N`$ expansion in quantum theory
As we discussed above, it is important to solve the fermionic constraint (1) in quantum theory by some nonperturbative method. Here we solve it with a fixed operator ordering by using the $`1/N`$ expansion. For systematic $`1/N`$ expansion, we introduce the bilocal formulation. We rewrite the fermionic constraint in terms of bilocal fields and expand it following the Holstein-Primakoff type expansion of the boson expansion method. We always work with fixed $`x^+`$.
### A Quantization and the operator ordering
To solve the constraint in quantum theory, we must first perform the Dirac quantization for constraint systems In Ref. , the authors solved the constraint equations for auxiliary fields before canonical quantization was specified and gave a c-number to the scalar auxiliary field in leading order of $`1/N`$. Nevertheless, the condensation in the NJL model is a quantum phenomenon and thus this procedure is not justified.. After tedious but straightforward calculation of the Dirac brackets, we find a familiar relation for the good component $`\psi _\alpha (\alpha =1,2)`$
$$\{\psi _\alpha ^a(x),\psi _\beta ^b(y)\}_{x^+=y^+}=\delta _{\alpha \beta }\delta _{ab}\delta (x^{}y^{})\delta ^{(2)}(x_{}y_{}),$$
(1)
and so on. We introduce the simplest mode expansion at $`x^+=0`$ as in Ref. :
$$\psi _\alpha ^a(x)=_{\mathrm{}}^{\mathrm{}}\frac{d^2k_{}}{2\pi }_0^{\mathrm{}}\frac{dk^+}{\sqrt{2\pi }}\left[b_\alpha ^a(𝒌)e^{i𝒌𝒙}+d_\alpha ^a(𝒌)e^{i𝒌𝒙}\right],$$
(2)
where $`𝒌𝒙k^+x^{}+k_{}^ix_{}^i`$. The vacuum is defined by the annihilation operators as
$$b_\alpha ^a(𝒌)|0=d_\alpha ^a(𝒌)|0=0.$$
(3)
When we deal with the quantum fermionic constraint, we have to specify the operator ordering. In many publications discussing the zero-mode constraints, people often choose the Weyl ordering with respect to both constrained and unconstrained variables. However, in a previous paper , we discussed that the ideal situation was to find a “consistent” operator ordering. For example, let us consider an anticommutator $`\{\chi ,\psi \}`$ in the NJL model. It can be evaluated in two different ways: (i) by using the solution $`\chi _{\mathrm{sol}}=\chi (\psi )`$ of the fermionic constraint and the standard quantization rule (1), and (ii) by calculating the Dirac bracket for $`\{\chi ,\psi \}`$. For the case (i), we select a specific operator ordering for the fermionic constraint, and the result depends on the ordering. For the case (ii), we must also determine the ordering in the r.h.s. of the Dirac bracket $`\{\chi ,\psi \}_\mathrm{D}=\mathrm{}`$. These two results must be equivalent to each other. We have two ambiguity of the operator ordering: those of the constraint equation in (i) and the right hand side of the Dirac bracket in (ii). “Consistent operator ordering” should be imposed so that these two quantities be identical. In other words, we determine the operator ordering of the r.h.s. in the Dirac brackets so that it coincides with the direct evaluation. In the chiral Yukawa model, we could not check that the ordering we adopted was consistent or not. Again in the NJL model, this is a very difficult task and we choose a specific operator ordering defined by Eq. (1). However, the chiral Gross-Neveu model in 1+1 dimensions allows us to check the consistency of this operator ordering. This is briefly shown in Appendix C.
### B Boson expansion method as $`1/N`$ expansion of bilocal operators
How can we solve the “operator equation” Eq. (1) by the $`1/N`$ expansion? It is generally difficult to count the order $`𝒪(N^n)`$ of an operator instead of its matrix element. What is worse, it is physically hard to justify the $`1/N`$ expansion of the fermionic field itself. However, as was discussed in Ref. , there is a powerful method to this problem. We can perform a systematic $`1/N`$ expansion of operators if we introduce the bilocal operators and use the boson expansion method. The boson expansion method is one of the traditional techniques in nonrelativistic many-body problems . Originally this was invented for describing bosonic excitations in non-bosonic systems such as collective excitation in nuclei or spin systems.
Let us rewrite the fermionic constraint (1) in terms of bilocal operators. We introduce the following “color” singlet bilocal operators at equal light-front time
$`_{\alpha \beta }(𝒙,𝒚)={\displaystyle \underset{a=1}{\overset{N}{}}}\psi _\alpha ^a(x^+,𝒙)\psi _\beta ^a(x^+,𝒚),`$ (4)
$`𝒯_{\alpha \beta }(𝒙,𝒚)={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{a=1}{\overset{N}{}}}\left(\psi _\alpha ^a(x^+,𝒙)\chi _\beta ^a(x^+,𝒚)+\chi _\beta ^a(x^+,𝒚)\psi _\alpha ^a(x^+,𝒙)\right),`$ (5)
$`𝒰_{\alpha \beta }(𝒙,𝒚)={\displaystyle \frac{i}{\sqrt{2}}}{\displaystyle \underset{a=1}{\overset{N}{}}}\left(\psi _\alpha ^a(x^+,𝒙)\chi _\beta ^a(x^+,𝒚)\chi _\beta ^a(x^+,𝒚)\psi _\alpha ^a(x^+,𝒙)\right).`$ (6)
We define the Fourier transformation of them as
$$_{\alpha \beta }(𝒑,𝒒)=_{\mathrm{}}^{\mathrm{}}\frac{d^3𝒙}{(2\pi )^{3/2}}_{\mathrm{}}^{\mathrm{}}\frac{d^3𝒚}{(2\pi )^{3/2}}_{\alpha \beta }(𝒙,𝒚)e^{i𝒑𝒙i𝒒𝒚},$$
and so on. Note that this definition allows the longitudinal momenta to take negative values. Using these bilocal operators, the fermionic constraint (1) is equivalently rewritten as
$`i{\displaystyle \frac{}{y^{}}}𝒯_{\alpha \beta }(𝒙,𝒚)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{_i^y\left(\sigma _{\beta \gamma }^i_{\alpha \gamma }(𝒙,𝒚)\sigma _{\gamma \beta }^i_{\gamma \alpha }(𝒚,𝒙)\right)+m_0\left(_{\alpha \beta }(𝒙,𝒚)_{\beta \alpha }(𝒚,𝒙)\right)\right\}`$ (9)
$`{\displaystyle \frac{g^2}{2}}\{_{\alpha \gamma }(𝒙,𝒚)(\delta _{\gamma \beta }𝒯(𝒚,𝒚)+i\sigma _{\gamma \beta }^3𝒰(𝒚,𝒚))`$
$`(\delta _{\beta \gamma }𝒯(𝒚,𝒚)i\sigma _{\beta \gamma }^3𝒰(𝒚,𝒚))_{\gamma \alpha }(𝒚,𝒙)\},`$
and
$`i^2{\displaystyle \frac{}{y^{}}}𝒰_{\alpha \beta }(𝒙,𝒚)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{_i^y\left(\sigma _{\beta \gamma }^i_{\alpha \gamma }(𝒙,𝒚)+\sigma _{\gamma \beta }^i_{\gamma \alpha }(𝒚,𝒙)\right)+m_0\left(_{\alpha \beta }(𝒙,𝒚)+_{\beta \alpha }(𝒚,𝒙)\right)\right\}`$ (12)
$`{\displaystyle \frac{g^2}{2}}\{_{\alpha \gamma }(𝒙,𝒚)(\delta _{\gamma \beta }𝒯(𝒚,𝒚)+i\sigma _{\gamma \beta }^3𝒰(𝒚,𝒚))`$
$`+(\delta _{\beta \gamma }𝒯(𝒚,𝒚)i\sigma _{\beta \gamma }^3𝒰(𝒚,𝒚))_{\gamma \alpha }(𝒚,𝒙)\},`$
where we have introduced quantities $`𝒯(𝒙,𝒚)𝒯_{\alpha \alpha }(𝒙,𝒚)`$ and $`𝒰(𝒙,𝒚)(\sigma ^3)_{\alpha \beta }𝒰_{\alpha \beta }(𝒙,𝒚)`$ so that $`\overline{\mathrm{\Psi }}\mathrm{\Psi }(x)=𝒯(𝒙,𝒙)`$ and $`\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }(x)=𝒰(𝒙,𝒙)`$. In actual calculation, it is more convenient to treat equations for the operators without spinor structure $`𝒯(𝒙,𝒚)`$ and $`𝒰(𝒙,𝒚)`$ because they form closed equations (see Appendix D). Once we solve them, we immediately obtain $`𝒯_{\alpha \beta }(𝒙,𝒚)`$ and $`𝒰_{\alpha \beta }(𝒙,𝒚)`$ from the above equations.
For systematic $`1/N`$ expansion of the bilocal fermionic constraints, one must know how to expand $`_{\alpha \beta }(𝒑,𝒒)`$. It is the boson expansion method, especially, the Holstein-Primakoff type expansion for large $`N`$ theories, that enables us to expand $`_{\alpha \beta }(𝒑,𝒒)`$ as operator quantities:
$$_{\alpha \beta }(𝒑,𝒒)=N\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{1}{\sqrt{N}}\right)^n\mu _{\alpha \beta }^{(n)}(𝒑,𝒒).$$
(13)
According to the Holstein-Primakoff expansion \[Eqs. (D12)$``$(D15)\], the first three terms are written in terms of bilocal bosonic variable $`B(𝒑,𝒒)`$ as
$`\mu _{\alpha \beta }^{(0)}(𝒑,𝒒)`$ $`=`$ $`\delta _{\alpha \beta }\delta ^{(3)}(𝒑+𝒒)\theta (p^+)\theta (q^+),`$ (14)
$`\mu _{\alpha \beta }^{(1)}(𝒑,𝒒)`$ $`=`$ $`B_{\beta \alpha }(𝒒,𝒑)\theta (p^+)\theta (q^+)+B_{\alpha \beta }^{}(𝒑,𝒒)\theta (p^+)\theta (q^+),`$ (15)
$`\mu _{\alpha \beta }^{(2)}(𝒑,𝒒)`$ $`=`$ $`{\displaystyle [d𝒌]\underset{\gamma }{}B_{\alpha \gamma }^{}(𝒑,𝒌)B_{\beta \gamma }(𝒒,𝒌)\theta (p^+)\theta (q^+)}`$ (16)
$``$ $`{\displaystyle [d𝒌]\underset{\gamma }{}B_{\gamma \beta }^{}(𝒌,𝒒)B_{\gamma \alpha }(𝒌,𝒑)\theta (p^+)\theta (q^+)}.`$ (17)
where
$$[d𝒒]=_0^{\mathrm{}}𝑑q^+_{\mathrm{}}^{\mathrm{}}d^2q_{}.$$
Any commutator between $`_{\alpha \beta }(𝒑,𝒒)`$’s \[such as Eq. (D10)\] is correctly reproduced if one uses the following bosonic commutators
$`[B_{\alpha \beta }(𝒑_1,𝒑_2),B_{\gamma \delta }^{}(𝒒_1,𝒒_2)]=\delta _{\alpha \gamma }\delta _{\beta \delta }\delta ^{(3)}(𝒑_1𝒒_1)\delta ^{(3)}(𝒑_2𝒒_2),`$ (18)
$`[B_{\alpha \beta }(𝒑_1,𝒑_2),B_{\gamma \delta }(𝒒_1,𝒒_2)]=0(p_i^+,q_i^+>0).`$ (19)
Note also that the state annihilated by $`B(𝒑,𝒒)`$ is identified with the original Fock vacuum:
$$B(𝒑,𝒒)|0=0.$$
(20)
More detailed discussions about the boson expansion method applied to LF field theories are found in Ref. and Appendix D of the present paper.
### C Solution to the bilocal fermionic constraint
We are ready to solve the bilocal fermionic constraint using the $`1/N`$ expansion. As we commented before, it is convenient to solve the equations for $`𝒯(𝒑,𝒒)`$ and $`𝒰(𝒑,𝒒)`$ (see Eqs. (D3) and (D6) in Appendix D for their explicit forms). Once we know $`𝒯(𝒑,𝒒)`$ and $`𝒰(𝒑,𝒒)`$, then it is straightforward to obtain $`𝒯_{\alpha \beta }(𝒑,𝒒)`$ and $`𝒰_{\alpha \beta }(𝒑,𝒒)`$.
Expanding $`𝒯(𝒑,𝒒)`$ and $`𝒰(𝒑,𝒒)`$ as
$`𝒯(𝒑,𝒒)=N{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{1}{\sqrt{N}}}\right)^nt^{(n)}(𝒑,𝒒),`$ (21)
$`𝒰(𝒑,𝒒)=N{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{1}{\sqrt{N}}}\right)^nu^{(n)}(𝒑,𝒒),`$ (22)
and inserting them into the fermionic constraints, we find for the lowest order $`𝒪(N)`$
$$\left(\begin{array}{c}t^{(0)}(𝒑,𝒒)\\ u^{(0)}(𝒑,𝒒)\end{array}\right)=\left(\begin{array}{c}m_0\frac{ϵ(p^+)}{q^+}\delta ^{(3)}(𝒑+𝒒)\\ 0\end{array}\right)g_0^2\frac{ϵ(p^+)}{q^+}_{\mathrm{}}^{\mathrm{}}\frac{d^3𝒌}{(2\pi )^3}\left(\begin{array}{c}t^{(0)}(𝒌,𝒑+𝒒𝒌)\\ u^{(0)}(𝒌,𝒑+𝒒𝒌)\end{array}\right),$$
(23)
where $`g_0^2=g^2N`$. Since there are no operators in these equations, $`t^{(0)}`$ and $`u^{(0)}`$ are $`c`$-numbers. Nonzero solutions give leading order contribution to $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ and $`\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }`$,
$`0\left|\overline{\mathrm{\Psi }}\mathrm{\Psi }\right|0=N{\displaystyle \frac{d^3𝒑}{(2\pi )^3}f_t(𝒑)}+\mathrm{},`$ (24)
$`0\left|\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }\right|0=N{\displaystyle \frac{d^3𝒑}{(2\pi )^3}f_u(𝒑)}+\mathrm{},`$ (25)
where $`t^{(0)}(𝒑,𝒒)=f_t(𝒑)\delta ^{(3)}(𝒑+𝒒)`$ and $`u^{(0)}(𝒑,𝒒)=f_u(𝒑)\delta ^{(3)}(𝒑+𝒒)`$. As the above equation (23) with $`m_0=0`$ is invariant under the chiral rotation, we can always take $`u^{(0)}(𝒑,𝒒)=0`$. For the massive case, we also take $`u^{(0)}(𝒑,𝒒)=0`$ and $`t^{(0)}(𝒑,𝒒)0`$. Now let us introduce a quantity $`M`$, which corresponds to the dynamical mass of fermion,
$$M=m_0g^2\overline{\mathrm{\Psi }}\mathrm{\Psi }.$$
(26)
Then, to obtain $`t^{(0)}(𝒑,𝒒)`$ is equivalent to determining $`M`$, viz.
$$t^{(0)}(𝒑,𝒒)=M\frac{ϵ(p^+)}{p^+}\delta ^{(3)}(𝒑+𝒒).$$
(27)
In terms of $`M`$, the leading order fermionic constraint (23) is rewritten as
$$\frac{Mm_0}{M}=g_0^2\frac{d^3𝒑}{(2\pi )^3}\frac{ϵ(p^+)}{p^+}.$$
(28)
Physically this equation should be interpreted as a gap equation. This is clarified in the next subsection.
Similarly higher order fermionic constraints are solved order by order. This is because the fermionic constraints for $`t^{(n)}(𝒑,𝒒)`$ and $`u^{(n)}(𝒑,𝒒)`$ are linear equations with respect to the highest order. The $`n=1,2`$ solutions are important for giving a nontrivial Hamiltonian and so on. More details are discussed in Appendix D.
### D Gap equation
Now let us discuss the physics meaning of Eq. (28). As we mentioned above, this equation should be regarded as a gap equation for chiral condensate. In several previous studies of ours, we have seen essentially the same kind of equations . Indeed, in Ref. it was pointed out that Eq. (28) itself is the gap equation. Also in the chiral Yukawa model , the zero-mode constraint for the scalar field reduced to the above equation and was interpreted as a gap equation. Since this identification is an indispensable step for our framework, let us again explain it within the NJL model.
First of all, consider a naive massless limit $`m_00`$ of Eq. (28):
$$M\left(1g_0^2\frac{d^3𝒑}{(2\pi )^3}\frac{ϵ(p^+)}{p^+}\right)=0.$$
Thus we find two possibilities: the first is $`M=0`$ and the second is
$$1g_0^2\frac{d^3𝒑}{(2\pi )^3}\frac{ϵ(p^+)}{p^+}=0$$
(29)
but $`M`$ is arbitrary. Of course the first case is not interesting because it corresponds to the symmetric phase. On the other hand, the second case with nonzero $`M`$ does not immediately mean the existence of broken phase. Since the equation (29) is independent of $`M`$ as it is, the dynamical mass $`M`$ is left undetermined, which is not a physically acceptable situation. However, this observation is not correct because the divergent integral in Eq. (29) is not regularized. Indeed, we can identify Eq. (29) with the gap equation only after we carefully treat the infrared (IR) divergence.
To see this, let us put an IR cutoff. First consider the same cutoff schemes as in the equal-time formulation, such as the covariant four-momentum cutoff. We can easily translate it into a cutoff on the light-cone momentum $`k^+`$ and $`k_{}^i`$ and obtain the correct gap equation. Indeed, in Ref. , a noncovariant (rotationally invariant) three-momentum cutoff was performed to obtain the known result. But such a cutoff is artificial as a light-front theory, and we here adopt another cutoff scheme, the parity invariant cutoff. Usually, it is natural and desirable to choose a cutoff so as to preserve symmetry of a system as much as possible. For the LF coordinates $`x^\pm `$ and $`x_{}^i`$, it would be natural to consider parity transformation ($`x^+x^{},x_{}^ix_{}^i`$) and two-dimensional rotation in the transverse plane. In the usual canonical formulation where $`x^+`$ is treated separately, the parity invariance is not manifest. However, we find it useful for obtaining the gap equation. In momentum space, the parity transformation is exchange of $`k^+`$ and $`k^{}`$ and replacement $`k_{}^ik_{}^i`$. Therefore the parity invariant cutoff is given by $`k^\pm <\mathrm{\Lambda }`$ and $`k_{}^2<\mathrm{\Lambda }^2`$. Using the dispersion relation $`2k^+k^{}k_{}^2=M^2`$, we find that the parity invariant regularization inevitably relates the ultraviolet (UV) and IR cutoffs:
$$\frac{k_{}^2+M^2}{2\mathrm{\Lambda }}<k^+<\mathrm{\Lambda }.$$
(30)
This also implies the planar rotational invariance $`k_{}^2<2\mathrm{\Lambda }^2M^2=\mathrm{\Lambda }^2`$. What is important here is the use of constituent mass $`M`$ in the dispersion relations. Physically it corresponds to imposing self-consistency conditions. Since the IR cutoff includes $`M`$, the r.h.s. of Eq. (28) has nontrivial dependence on $`M`$:
$$\frac{Mm_0}{M}=\frac{g_0^2\mathrm{\Lambda }^2}{4\pi ^2}\left\{2\frac{M^2}{\mathrm{\Lambda }^2}\left(1+\mathrm{ln}\frac{2\mathrm{\Lambda }^2}{M^2}\right)\right\}.$$
(31)
This is the gap equation and is equivalent to that of the previous result in the chiral Yukawa model . It has a nonzero solution $`M0`$ even in the $`m_00`$ limit. The somewhat unfamiliar equation (31) of the NJL model exhibits the same property as the standard gap equations of the equal-time quantization. For example, when $`m_0=0`$, there is a critical coupling $`g_{\mathrm{cr}}^2=2\pi ^2/\mathrm{\Lambda }^2`$, above which $`M0`$.
The essential and inevitable step to obtain the gap equation is the inclusion of mass information as the regularization rather than the fact that the UV and IR cutoffs are related to each other. If we regulate the divergent integral without mass information ( e.g., introducing the UV and IR cutoffs independently), we cannot reproduce the gap equation. The loss of mass information is closely related to the fundamental problem of the LF formalism , and the parity invariant regularization can be considered as one of the prescriptions for it. Reference discussed within scalar theory that the light-front quantization gives a mass-independent two-point function (at equal LF time), which contradicts the result from general arguments concerning the spectral representation. We have been encountered with the same problem in Eq. (28) because the integral is regarded as a naive estimation of $`\overline{\mathrm{\Psi }}\mathrm{\Psi }/M`$ by using fermion with mass $`M`$. And also the origin of mass-independent result can be traced back to the mode expansion (2). Even if we include the wave function for free fermion field, we do not have any mass-dependence on the mode expansion .
Let us give a brief comment on the chiral Gross-Neveu model. Of course the important difference of the 1+1 dimensional case is the renormalizability. Ignoring the transverse directions in the above calculation, we easily find the gap equation $`(Mm_0)/M=g_0^2/(2\pi )\mathrm{ln}(2\mathrm{\Lambda }^2/M^2)`$ where the parity-invariant cutoff $`M^2/2\mathrm{\Lambda }<k^+<\mathrm{\Lambda }`$ was used. Though it explicitly depends on the cutoff $`\mathrm{\Lambda }`$ and is divergent as $`\mathrm{\Lambda }\mathrm{}`$, we can remove the divergence by coupling constant renormalization .
### E Hamiltonian
Having the solution to the bilocal fermionic constraint, we can rewrite the fermion bilinear operators in terms of the bilocal bosons. Of special importance is the (Hermitian) Hamiltonian, which is easily expressed by $`𝒯_{\alpha \beta }(𝒑,𝒒)`$ and $`𝒰_{\alpha \beta }(𝒑,𝒒)`$ as follows
$`H=P^{}`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}{\displaystyle d^3x\left[\left(\psi ^{}\sigma ^i_i\chi +_i\chi ^{}\sigma ^i\psi \right)+m_0\left(\psi ^{}\chi +\chi ^{}\psi \right)\right]}`$ (32)
$`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle d^3𝒑d^3𝒒\delta ^{(3)}(𝒑+𝒒)iq_{}^i\sigma _{\alpha \beta }^i\left[\left(𝒯_{\alpha \beta }(𝒑,𝒒)+𝒯_{\beta \alpha }(𝒒,𝒑)\right)+i\left(𝒰_{\alpha \beta }(𝒑,𝒒)𝒰_{\beta \alpha }(𝒑,𝒒)\right)\right]}`$ (34)
$`+{\displaystyle \frac{m_0}{2}}{\displaystyle d^3𝒑d^3𝒒𝒯(𝒑,𝒒)\delta ^{(3)}(𝒑+𝒒)}.`$
Apparently this Hermitian version of the Hamiltonian seems equivalent to the free Hamiltonian, but the information of interaction enters through the bad component $`\chi `$. We find the $`1/N`$ expansion of the Hamiltonian
$$H=N\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{1}{\sqrt{N}}\right)^nh^{(n)},$$
(35)
by substituting the solutions of the fermionic constraints into the Hamiltonian. The zeroth order contribution is just a divergent constant and we discard it. The first order is strictly zero. Nontrivial contribution comes from the order $`𝒪(N^0)`$
$`h^{(2)}`$ $`=`$ $`{\displaystyle [d𝒑][d𝒒]\left(\frac{p_{}^2+M^2}{2p^+}+\frac{q_{}^2+M^2}{2q^+}\right)B_{\alpha \beta }^{}(𝒑,𝒒)B_{\alpha \beta }(𝒑,𝒒)}`$ (39)
$`+{\displaystyle \frac{g_0^2}{(2\pi )^3}}{\displaystyle [d𝒑][d𝒒][d𝒌][d𝒍]\delta ^{(3)}(𝒑+𝒒𝒌𝒍)\alpha (p^++q^+)}`$
$`\times [{\displaystyle }{}_{\gamma \delta }{}^{\alpha \beta }(𝒑,𝒒;𝒌,𝒍){\displaystyle }{}_{\gamma \delta }{}^{\alpha \beta }(𝒑,𝒒;𝒌,𝒍)]B_{\alpha \beta }^{}(𝒑,𝒒)B_{\gamma \delta }(𝒌,𝒍)`$
$`+c\mathrm{number},`$
where $`\alpha (p^++q^+)`$ is defined by Eq. (D20) and “kernels” of the interaction terms are
$`{\displaystyle }{}_{\gamma \delta }{}^{\alpha \beta }(𝒑,𝒒;𝒌,𝒍)[\mathrm{SS}(𝒑)\mathrm{SS}(𝒒)]_{\alpha \beta }[\mathrm{SS}(𝒌)\mathrm{SS}(𝒍)]_{\delta \gamma },`$ (40)
$`{\displaystyle }{}_{\gamma \delta }{}^{\alpha \beta }(𝒑,𝒒;𝒌,𝒍)[𝒫(𝒑)𝒫(𝒒)]_{\alpha \beta }[𝒫(𝒌)𝒫(𝒍)]_{\delta \gamma },`$ (41)
with
$$\mathrm{SS}_{\alpha \beta }(𝒑)=\left(\frac{ip^i\sigma ^iM}{2p^+}\right)_{\alpha \beta },𝒫_{\alpha \beta }(𝒑)=\left(\frac{ip^i\sigma ^iM}{2p^+}\sigma ^3\right)_{\alpha \beta }.$$
(42)
As is evident from the explicit forms of the kernels (40) and (41), they originate from the scalar interaction $`(\overline{\mathrm{\Psi }}\mathrm{\Psi })^2`$ and the pseudoscalar one $`(\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi })^2`$, respectively. If we substitute a nontrivial (trivial) solution of the gap equation (31) into the above Hamiltonian, then it governs dynamics of the broken (symmetric) phase. The first term of $`h^{(2)}`$ corresponds to a free part with mass $`M`$ and the second term to an interaction part. In the broken phase, $`M`$ is the dynamical mass and the Hamiltonian suggests a constituent picture.
As we mentioned before, the Hermite Hamiltonian of the chiral Gross-Neveu model has only an $`m_0`$-dependent term. Neglecting the transverse coordinates in Eq. (34), we have
$$P_{\mathrm{GN}}^{}=\frac{m_0}{2\sqrt{2}}𝑑x^{}\left(\psi ^{}\chi +\chi ^{}\psi \right).$$
Furthermore, the classical solution for the bad spinor component $`\chi `$ is proportional to $`m_0`$. Therefore naive $`m_00`$ limit gives just a zero Hamiltonian. However, if we solve the gap equation and substitute the nontrivial solution into the Hamiltonian, the resulting Hamiltonian turns out to be proportional to $`M^2`$ and survives even in the chiral limit. This is easily seen from the Hamiltonian of the NJL model (39). The (constituent) mass term in Eq. (39) comes from the bare mass term, whose factor $`m_0`$ cancels with a factor $`M^2/m_0`$ in the second order solution $`d^3𝒑t^{(2)}(𝒑,𝒑)`$. Of course this is not reached if we set $`m_0=0`$ from the beginning. Therefore inclusion of the bare mass term is necessary to obtain a correct (constituent) mass term of the Hamiltonian.
## IV Physics in the broken phase
By solving the fermionic constraint, we acquired necessary ingredients for discussing physics consequences of the chiral symmetry breaking. Basic quantities such as $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$, $`\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }`$, and the null-plane chiral charge (15) are expressed in terms of the bilocal bosons $`B_{\alpha \beta }(𝒑,𝒒)`$ and $`B_{\alpha \beta }^{}(𝒑,𝒒)`$ as
$`\overline{\mathrm{\Psi }}\mathrm{\Psi }(x)=𝒯(𝒙,𝒙)={\displaystyle \frac{N}{g_0^2}}(m_0M)+\sqrt{N}{\displaystyle \frac{d^3𝒑d^3𝒒}{(2\pi )^3}t^{(1)}(𝒑,𝒒)e^{i(𝒑+𝒒)𝒙}}+𝒪(N^0),`$ (1)
$`\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }(x)=𝒰(𝒙,𝒙)=\sqrt{N}{\displaystyle \frac{d^3𝒑d^3𝒒}{(2\pi )^3}u^{(1)}(𝒑,𝒒)e^{i(𝒑+𝒒)𝒙}}+𝒪(N^0),`$ (2)
$`Q_5^{\mathrm{LF}}={\displaystyle }d^3𝒑\sigma _{\alpha \beta }^3:_{\alpha \beta }(𝒑,𝒑):={\displaystyle }d^3𝒑\sigma _{\alpha \beta }^3\mu _{\alpha \beta }^{(2)}(𝒑,𝒑)+𝒪(N^{1/2}),`$ (3)
where $`t^{(1)}(𝒑,𝒒)`$ and $`u^{(1)}(𝒑,𝒒)`$ are given in Appendix D. Now that these are given as functions of the bilocal bosons at the operator level, all the calculation is done with the commutators (18) and (19).
### A Chiral transformation and nonconservation of chiral charge
Why could we obtain a nonzero fermion condensate? To understand this, let us rewrite the fermionic constraint (1) as
$$i_{}\chi _a=\frac{1}{\sqrt{2}}\left(\sigma ^i_i+m_0\right)\psi _a\frac{g^2}{\sqrt{2}}\left(\psi _a𝒯(𝒙,𝒙)+i\sigma ^3\psi _a𝒰(𝒙,𝒙)\right),$$
and substitute Eqs. (1) and (2) into $`𝒯(𝒙,𝒙)`$ and $`𝒰(𝒙,𝒙)`$, respectively<sup>§</sup><sup>§</sup>§In the leading order, this procedure corresponds to the mean-field approximation done by Heinzl et al. . They solved the fermionic constraint by simply linearizing the interaction parts as $`g^2/\sqrt{2}\psi _a𝒯(𝒙,𝒙)`$. By evaluating the vacuum expectation value $`𝒯(𝒙,𝒙)`$ self-consistently with the dynamical fermion mass $`M=m_0g^2𝒯(𝒙,𝒙)`$, they obtained the gap equation. If one uses the parity-invariant cutoff, the result coincides with ours.. Then the leading order equation turns out to be equivalent to the constraint equation for a free fermion with mass $`M`$,
$$i_{}\chi _a=\frac{1}{\sqrt{2}}\left(\sigma ^i_i+M\right)\psi _a.$$
Also at the same order, the equation of motion for the good component $`\psi `$ says that the fermion acquires a mass $`M`$. This means that the operator structure of the bad spinor $`\chi `$ changes depending on which solutions of the gap equation (28) is selected. For massive fermion, the fermion condensate $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ is no longer zero even if the vacuum is trivial. One can find an analogy between the chiral Yukawa model and the NJL model because in the chiral Yukawa model, operator structure of the longitudinal zero modes and subsequently of the bad spinor component changes depending on the phases.
One thing to be noted is the peculiarity of the mode expansion (2). It is evident that the mode expansion has no mass dependence in it. This caused the problem of identifying the lowest fermionic constraint with the gap equation. We had to supply mass information properly when we regularize the IR divergence. On the other hand, such independence of mass, in turn, implies that our mode expansion allows fermions with any value of mass. In other words, the LF vacuum does not distinguish mass of the fermion. Therefore we can regard the vacuum for massless fermion as that for massive one. Mass of the fields is determined by the Hamiltonian. This is the reason why we can live with the trivial vacuum while having a nonzero fermion condensate. This fact is not a limited phenomenon for our specific mode expansion but a common one for light-front field theories. Indeed, even if we expand a fermion field with free spinor wave functions, $`u(p)`$ and $`v(p)`$, we have no mass dependence .
The fact that the operator structure changes depending on the phases, also resolves a seeming contradiction between the triviality of the null-plane chiral charge and the nonzero chiral condensate $`0\left|\overline{\mathrm{\Psi }}\mathrm{\Psi }\right|00`$. In general, it is known that a null-plane charge always annihilates the vacuum irrespective of the presence of symmetry. This can be checked explicitly by the expression (3), viz.
$$Q_5^{\mathrm{LF}}|0=0.$$
(4)
However, the triviality of $`Q_5^{\mathrm{LF}}`$ in the presence of the chiral condensate immediately leads to a contradiction if an equation $`[Q_5^{\mathrm{LF}},\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }]=2i\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ could hold in the broken phase. In the previous analysis of the chiral Yukawa model , we were faced with exactly the same problem and resolved it by recognizing that in the broken phase the chiral transformation of dependent variables are different from the usual one simply because their operator structure changes. This is of course true of the NJL model. First of all, as we saw above, if we select the nontrivial solution of the gap equation, the fermion is no longer a massless fermion even in the chiral limit. Secondly, we can explicitly show that the usual transformation law $`[Q_5^{\mathrm{LF}},\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }]=2i\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ holds only in the symmetric phase ($`M=0`$). In the broken phase a simple calculation (up to $`𝒪(N^{1/2})`$) leads to
$`[Q_5^{\mathrm{LF}},\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }(x)]`$ (5)
$`=2i\overline{\mathrm{\Psi }}\mathrm{\Psi }(x)+2i{\displaystyle \frac{N}{g_0^2}}(m_0M)+2i\sqrt{N}M{\displaystyle \frac{d𝒑d𝒒}{(2\pi )^3}e^{i(𝒑+𝒒)𝒙}}`$ (6)
$`\times \left[{\displaystyle \frac{\mu _{\alpha \alpha }^{(1)}(𝒑,𝒒)}{q^+}}g_0^2{\displaystyle \frac{ϵ(p^+)}{q^+}}\alpha (p^++q^+){\displaystyle \frac{d^3𝒌}{(2\pi )^3}\frac{\mu _{\alpha \alpha }^{(1)}(𝒌,𝒑+𝒒𝒌)}{p^++q^+k^+}}\right]+𝒪(N^0).`$ (7)
Even if we take the chiral limit $`m_00`$, the extra term proportional to $`M`$ survives nonzero. This also implies that if we put $`M=0`$, the usual relation holds. The unusual chiral transformation, however, is consistent with the triviality of $`Q_5^{\mathrm{LF}}`$ because $`0\left|[Q_5^{\mathrm{LF}},\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }]\right|0=0`$.
A similar situation occurs for the Hamiltonian. Nonconservation of the null-plane chiral charge has been pointed out by several people as a characteristic feature of the chiral symmetry breaking on the LF . They discussed it under the assumption of the PCAC relation, but we can check it explicitly by using the broken Hamiltonian. After lengthy calculation, we find the commutator $`[Q_5^{\mathrm{LF}},H]`$ is really nonzero and again proportional to the dynamical mass $`M`$:
$`[Q_5^{\mathrm{LF}},H]`$ $`=`$ $`M{\displaystyle \frac{g_0^2}{16\pi ^3i}}{\displaystyle [d𝒑][d𝒒][d𝒌][d𝒍]\delta ^{(3)}(𝒌+𝒍𝒑𝒒)\alpha (p^++q^+)B_{\alpha \beta }^{}(𝒑,𝒒)B_{\gamma \delta }(𝒌,𝒍)}`$ (11)
$`\times [({\displaystyle \frac{p_{}^i}{p^+}}{\displaystyle \frac{q_{}^i}{q^+}})\sigma _{\alpha \beta }^i({\displaystyle \frac{1}{k^+}}+{\displaystyle \frac{1}{l^+}})\sigma _{\delta \gamma }^3({\displaystyle \frac{1}{p^+}}{\displaystyle \frac{1}{q^+}})\delta _{\alpha \beta }({\displaystyle \frac{k_{}^i}{k^+}}{\displaystyle \frac{l_{}^i}{l^+}})(\sigma ^i\sigma ^3)_{\delta \gamma }`$
$`({\displaystyle \frac{p_{}^i}{p^+}}{\displaystyle \frac{q_{}^i}{q^+}})(\sigma ^i\sigma ^3)_{\alpha \beta }({\displaystyle \frac{1}{k^+}}+{\displaystyle \frac{1}{l^+}})\delta _{\delta \gamma }+({\displaystyle \frac{1}{p^+}}+{\displaystyle \frac{1}{q^+}})\sigma _{\alpha \beta }^3({\displaystyle \frac{k_{}^i}{k^+}}{\displaystyle \frac{l_{}^i}{l^+}})\sigma _{\delta \gamma }^i]`$
$`+𝒪(N^{\frac{1}{2}}).`$
Therefore, the LF chiral charge is not conserved even in the chiral limit. In our framework it would be more understandable to mention that the Hamiltonian is not invariant under the LF chiral transformation in the broken phase. The broken phase Hamiltonian (39) has three terms: $`M`$-independent, linearly dependent and quadratically dependent terms. The quadratically dependent term, as well as the $`M`$-independent one, does not break the LF chiral symmetry. It is the term proportional to the dynamical fermion mass $`M`$ which breaks the LF chiral symmetry. And also, since Eq. (11) is proportional to $`g_0^2`$, the symmetry breaking term purely comes from the interactionFor a massive free fermion, we have $`[Q_5^{\mathrm{LF}},H]=0`$..
This result should be consistent with the current divergence relation Eq. (18). Integrating it over the space, we have
$$_+Q_5^{\mathrm{LF}}=\frac{1}{i}[Q_5^{\mathrm{LF}},H]=2m_0𝑑x^{}d^2x_{}\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }.$$
(12)
Therefore if the LF chiral charge is not conserved in the chiral limit, the r.h.s. must show a singular behavior
$$𝑑x^{}d^2x_{}\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }\frac{1}{m_0}.$$
(13)
This can be verified directly by using the solution of the fermionic constraint. Indeed we find that $`𝑑x^{}d^2x_{}\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }=𝑑𝒑u^{(2)}(𝒑,𝒑)`$ is proportional to $`M/m_0`$ and gives exactly the same result as Eq. (11). Importance of such singular behavior for making the Nambu-Goldstone boson meaningful was stressed by Tsujimaru et al. in scalar theories . Assuming the PCAC relation, they showed that the zero mode of the Nambu-Goldstone boson has a singularity of $`m_{\mathrm{NG}}^1`$ where $`m_{\mathrm{NG}}`$ is an explicit symmetry-breaking mass. Our result (13) is consistent with theirs because the operator $`\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }`$ is directly related to the Nambu-Goldstone boson. Later, we will prove that the PCAC relation is derived from the current divergence relation (18).
### B LF bound-state equations for mesons and their solutions
#### 1 Single bosonic state as a fermion-antifermion state
In our formulation with the boson expansion method, any bosonic excited state is described by the Fock states of the bilocal bosons constructed on the trivial vacuum:
$$\underset{i}{}B_{\alpha _i\beta _i}^{}(𝒑_i,𝒒_i)|0.$$
(14)
Since the Hamiltonian (39) is quadratic with respect to the bilocal bosons, the first excited state is given by a single bosonic state. In fermionic degrees of freedom, the one boson state corresponds to the leading contribution (of $`1/N`$ expansion) of a fermion-antifermion state. To see this, let us write a mesonic state only with a “color” singlet fermion-antifermion Fock component:
$$|\mathrm{meson};P^+,P_{}=\frac{1}{\sqrt{N}}_0^{P^+}𝑑k^+_{\mathrm{}}^{\mathrm{}}d^2k_{}\mathrm{\Phi }^{\alpha \beta }(𝒌)b_\alpha ^a(𝒌)d_\beta ^a(𝑷𝒌)|0,$$
(15)
where the LC wavefunction $`\mathrm{\Phi }^{\alpha \beta }(𝒌)`$ is normalized so as to satisfy the condition
$$\mathrm{meson};𝑷|\mathrm{meson};𝑸\text{ }=(2\pi )^32P^+\delta ^{(3)}(𝑷𝑸),$$
(16)
or equivalently,
$$_0^1𝑑x\frac{d^2k_{}}{16\pi ^3}\underset{\alpha \beta }{}|\mathrm{\Phi }^{\alpha \beta }(𝒌)|^2=1.$$
(17)
According to the Holstein-Primakoff type expansion (D15), the fermion-antifermion operator $`b_\alpha ^{}d_\beta ^{}`$ can be equivalently rewritten as
$`b_\alpha ^a(𝒌)d_\beta ^a(𝑷𝒌)`$ $`=`$ $`:_{\alpha \beta }^{}(𝒌,𝑷+𝒌):`$ (18)
$`=`$ $`\sqrt{N}B_{\alpha \beta }^{}(𝒌,𝑷𝒌)`$ (20)
$`{\displaystyle \frac{1}{2\sqrt{N}}}{\displaystyle [d𝒒][d𝒒^{}]B_{\gamma \beta }^{}(𝒒,𝑷𝒌)B_{\alpha \delta }^{}(𝒌,𝒒^{})B_{\gamma \delta }(𝒒,𝒒^{})}+\mathrm{}.`$
Therefore at the leading order of $`1/N`$ expansion, the mesonic state is described as a single (bilocal) boson state,
$$|\mathrm{meson};P^+,P_{}=_0^{P^+}𝑑k^+_{\mathrm{}}^{\mathrm{}}d^2k_{}\mathrm{\Phi }^{\alpha \beta }(𝒌)B_{\alpha \beta }^{}(𝒌,𝑷𝒌)|0+𝒪(N^{1/2}).$$
(21)
Besides this, it is evident from the normalization condition (16), a local operator $`a^{}(𝑷)=d^3k\mathrm{\Phi }^{\alpha \beta }(𝒌)B_{\alpha \beta }^{}(𝒌,𝑷𝒌)`$ satisfies the usual bosonic commutators.
The LC wavefunction $`\mathrm{\Phi }^{\alpha \beta }(𝒌)`$ and the mass of a meson $`M_{\mathrm{meson}}`$ is determined by solving the light-front eigen-value equation:
$$h^{(2)}|\mathrm{meson};P^+,P_{}=0=\frac{M_{\mathrm{meson}}^2}{2P^+}|\mathrm{meson};P^+,P_{}=0,$$
(22)
where we set $`P_{}^i=0`$, for simplicity.
#### 2 Scalar and Pseudoscalar mesons
In the leading order of $`1/N`$ expansion, the Hamiltonian has only quadratic terms of the bosonic operators. Therefore, diagonalization of the Hamiltonian, or equivalently, solving the light-cone bound state equation (22) is straightforward. First of all, if one notices the orthogonal property $`\left(𝒫(𝒌)𝒫(𝒍)\right)_{\alpha \beta }\left(\mathrm{SS}(𝒌)\mathrm{SS}(𝒍)\right)_{\beta \alpha }=0`$ where $`𝒌=(xP^+,k_{}^i)`$ and $`𝒍=𝑷𝒌=((1x)P^+,k_{}^i)`$, one can easily find the spinor structure for scalar ($`\sigma `$) and pseudoscalar ($`\pi `$) states should be
$`|\pi ;P^+,P_{}=0=P^+{\displaystyle _0^1}𝑑x{\displaystyle d^2k_{}\varphi _\pi (x,k_{}^i)\left\{\left(ik_{}^i\sigma ^i+M\right)\sigma ^3\right\}_{\alpha \beta }B_{\alpha \beta }^{}(𝒌,𝒍)|0},`$ (23)
$`|\sigma ;P^+,P_{}=0=P^+{\displaystyle _0^1}𝑑x{\displaystyle d^2k_{}\varphi _\sigma (x,k_{}^i)\left\{ik_{}^i\sigma ^i+(12x)M\right\}_{\alpha \beta }B_{\alpha \beta }^{}(𝒌,𝒍)|0}.`$ (24)
These two states are orthogonal to each other. Somewhat nonstandard spinor structure of the mesonic states is due to our specific choice of the mode expansion Eq. (2) and the representation for $`\gamma `$ matrices Eq. (4). For example, if one rewrites the pseudoscalar field $`\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }`$ in terms of the bilocal bosons, one finds the same spinor structure as that of Eq. (23). Note also that $`\left\{\gamma _5(\gamma _{}^ik_{}^i+M)\right\}_{\alpha \beta }=\left\{\left(ik_{}^i\sigma ^i+M\right)\sigma ^3\right\}_{\alpha \beta }`$ for $`\alpha ,\beta =1,2`$ in our two-component representation for the $`\gamma `$ matrices.
Spinor independent parts of the LC wavefunctions $`\varphi _{\pi ,\sigma }(x,k_{}^i)`$ are given as solutions of the following integral equations:
$`m_\pi ^2\varphi _\pi (x,k_{}^i)`$ $`=`$ $`{\displaystyle \frac{k_{}^2+M^2}{x(1x)}}\varphi _\pi (x,k_{}^i){\displaystyle \frac{g_0^2\alpha }{(2\pi )^3}}{\displaystyle \frac{1}{x(1x)}}{\displaystyle _0^1}𝑑y{\displaystyle d^2\mathrm{}_{}\frac{\mathrm{}_{}^2+M^2}{y(1y)}\varphi _\pi (y,\mathrm{}_{}^i)},`$ (25)
$`m_\sigma ^2\varphi _\sigma (x,k_{}^i)`$ $`=`$ $`{\displaystyle \frac{k_{}^2+M^2}{x(1x)}}\varphi _\sigma (x,k_{}^i){\displaystyle \frac{g_0^2\alpha }{(2\pi )^3}}{\displaystyle \frac{1}{x(1x)}}{\displaystyle _0^1}𝑑y{\displaystyle d^2\mathrm{}_{}\frac{\mathrm{}_{}^2+(12y)^2M^2}{y(1y)}\varphi _\sigma (y,\mathrm{}_{}^i)}.`$ (26)
Here the factor $`\alpha =\alpha (P^+)`$ defined by Eq. (D20) is given as a result of the gap equation,
$$\alpha =\left(\frac{m_0}{M}+\frac{2g_0^2}{(2\pi )^3}d^2q_{}_0^1\frac{dx}{x}\right)^1.$$
(27)
Since these integral equations are separable ones, solutions are easily found
$`\varphi _\pi (x,k_{}^i)=C_\pi {\displaystyle \frac{g_0^2}{(2\pi )^3}}{\displaystyle \frac{M}{m_0}}{\displaystyle \frac{1}{x(1x)(k_{}^2+M^2)/m_\pi ^2}},`$ (28)
$`\varphi _\sigma (x,k_{}^i)=C_\sigma {\displaystyle \frac{g_0^2}{(2\pi )^3}}{\displaystyle \frac{M}{m_0}}\left({\displaystyle \frac{m_\sigma ^24M^2}{m_\sigma ^2}}\right){\displaystyle \frac{1}{x(1x)(k_{}^2+M^2)/m_\sigma ^2}},`$ (29)
where $`C_\pi `$ and $`C_\sigma `$ are constants $`C_{\pi ,\sigma }=_0^1𝑑xd^2k_{}\varphi _{\pi ,\sigma }(x,k_{}^i)`$.
Equations for $`m_\pi `$ and $`m_\sigma `$ are derived from the normalization condition for the LC wavefunctions, viz.
$`1=g_0^2{\displaystyle \frac{M}{m_0}}{\displaystyle _0^1}𝑑x{\displaystyle \frac{d^2k_{}}{(2\pi )^3}\frac{m_\pi ^2}{k_{}^2+M^2m_\pi ^2x(1x)}},`$ (30)
$`1=g_0^2{\displaystyle \frac{M}{m_0}}{\displaystyle _0^1}𝑑x{\displaystyle \frac{d^2k_{}}{(2\pi )^3}\frac{m_\sigma ^24M^2}{k_{}^2+M^2m_\sigma ^2x(1x)}}.`$ (31)
These are exactly equivalent to the corresponding equations in the previous work on the chiral Yukawa model , where we obtained them by calculating pole masses of the scalar and pseudoscalar bosons. If one uses the same cutoff scheme (extended parity-invariant cutoff) as in Ref. :
$$\frac{k_{}^2+M^2}{x}+\frac{k_{}^2+M^2}{1x}<2\mathrm{\Lambda }^2,$$
(32)
the pseudoscalar mass for small bare mass $`m_0`$ is
$$m_\pi ^2=\frac{m_0N}{g_0^2M}Z_\pi +𝒪(m_0^2),$$
(33)
where a cutoff dependent factor
$$Z_\pi =\frac{1}{N}\left[\frac{1}{8\pi ^2}\mathrm{ln}\left(\frac{1+\sqrt{12M^2/\mathrm{\Lambda }^2}}{1\sqrt{12M^2/\mathrm{\Lambda }^2}}\right)\frac{\sqrt{12M^2/\mathrm{\Lambda }^2}}{4\pi ^2}\right]^1$$
(34)
is related to normalization of a pseudoscalar state \[see Eq. (44)\]. Clearly $`m_\pi `$ vanishes in the chiral limit $`m_00`$ and the pseudoscalar state is identified with the Nambu-Goldstone boson. In Eq. (25), the first term corresponds to a kinetic energy part of the fermion and antifermion with the constituent mass $`M`$ and the second term, a potential energy part. The masslessness of the pseudoscalar state in the chiral limit is realized by the exact cancellation between the kinetic energy and the potential energy. Indeed, if we integrate Eq. (25) over $`x`$ and $`k_{}^i`$, we find
$`m_\pi ^2{\displaystyle _0^1}𝑑x{\displaystyle d^2k_{}\varphi _\pi (x,k_{}^i)}`$ $`=`$ $`\left(1{\displaystyle \frac{2g_0^2_0^1\frac{dx}{x}\frac{d^2q_{}}{(2\pi )^3}}{\frac{m_0}{M}+2g_0^2_0^1\frac{dx}{x}\frac{d^2q_{}}{(2\pi )^3}}}\right){\displaystyle _0^1}𝑑y{\displaystyle d^2\mathrm{}_{}\frac{\mathrm{}_{}^2+M^2}{y(1y)}\varphi _\pi (y,\mathrm{}_{}^i)}`$ (35)
$``$ $`0(m_00).`$ (36)
Therefore, $`m_\pi =0`$ is fulfilled in the chiral limit even though the fermion has the constituent mass.
On the other hand, squared mass of the scalar state for small $`m_0`$ is
$$m_\sigma ^2=4M^2+𝒪(m_0).$$
(37)
At a first glance, the result $`m_\sigma =2M`$ in the chiral limit seems to suggest a static picture of a fermion and an antifermion, but actually the mass of the scalar meson comes from a part of the potential energy. The kinetic energy cancels with the rest of the potential energy.
Equations (28) and (29) have the same functional form with respect to the variables $`x`$ and $`k_{}^i`$. However, the difference between $`m_\pi `$ and $`m_\sigma `$ greatly affects the shape of the LC wavefunctions. This is most remarkable in the chiral limit: As $`m_00`$, Eq. (28) becomes independent of $`x`$:
$`\varphi _\pi (x,k_{}^i)i\sqrt{N}\sqrt{Z_\pi }{\displaystyle \frac{1}{k_{}^2+M^2}},`$ (38)
where the constant $`C_\pi `$ was evaluated as $`C_\pi i(2\pi )^3(NZ_\pi )^{1/2}.`$ On the other hand, Eq. (29) shows a narrow peak at $`x=1/2`$. Therefore, the pseudoscalar state is a highly collective state, while the scalar state shows an approximate constituent picture.
Now let us compare our result Eq. (28) with those of the literatures . First of all, equivalence with the result of Ref. is easily verified. As we commented before, the unfamiliar spinor structure in Eq. (23) is due to our specific choice of the mode expansion and the representation for the $`\gamma `$ matrices. If one uses the following mode expansion for the good component of the fermion:
$$\psi _+(x)=\underset{\lambda }{}_{\mathrm{}}^{\mathrm{}}\frac{d^2k_{}}{2\pi }_0^{\mathrm{}}\frac{dk^+}{\sqrt{2\pi k^+}}\left[\stackrel{~}{b}(𝒌,\lambda )u_+(𝒌,\lambda )e^{i𝒌𝒙}+\stackrel{~}{d}^{}(𝒌,\lambda )v_+(𝒌,\lambda )e^{i𝒌𝒙}\right],$$
one can obtain the same spinor structure as that of Ref. . Of course the two LC wavefunctions should coincide with each other for observable quantities. Indeed, both give the same (quark) distribution function $`q(x)=d^2k_{}/(2\pi )^3_{\alpha ,\beta }|\mathrm{\Phi }_{\alpha \beta }(𝒌)|^2`$ .
On the other hand, the result of Ref. seems different from ours Eq. (28). The possible origin of the discrepancy might be attributed to the following two points. First of all, the author of Ref. considered the Melosh transformation which relates the LF spinor and the usual spinor in the equal-time quantization. Such nonstatic spin effects might be important when we discuss phenomenological aspects of light mesons (for example, see Ref. ). However, even if we take it into account, it is hard to see the coincidence. Secondary but most importantly, he derived the pion LC wavefunction by projecting the Bethe-Salpeter amplitude on the equal LC time plane. Though this procedure should give the same result as that of the LF bound-state equation as far as we are considering only the ladder ($`1/N`$ leading) contribution, equivalence of the two is a highly nontrivial problem in our complicated analysis.
### C The Gell-Mann–Oakes–Renner and PCAC relations
Now that we have the LC wavefunction for the pseudoscalar meson, it is straightforward to obtain the decay constant $`f_\pi `$:
$$iP^\mu f_\pi =0\left|j_5^\mu (0)\right|\pi ;𝑷.$$
(39)
For actual calculation, it is safer and easier to treat the plus component. If we use the extended parity-invariant cutoff, the result is
$$f_\pi =2MZ_\pi ^{\frac{1}{2}}+𝒪(N^0).$$
(40)
Together with the pseudoscalar mass (33) in the chiral limit, we find the Gell-Mann, Oakes, and Renner relation,
$`m_\pi ^2f_\pi ^2`$ $`=`$ $`4m_0\left({\displaystyle \frac{NM}{g_0^2}}\right)`$ (41)
$`=`$ $`4m_00\left|\overline{\mathrm{\Psi }}\mathrm{\Psi }\right|0.`$ (42)
The PCAC relation is also checked by using the state $`|\pi ;𝑷`$. If we normalize the pseudoscalar field $`\pi _\mathrm{n}(x)\overline{\mathrm{\Psi }}(x)i\gamma _5\mathrm{\Psi }(x)`$ as
$$0\left|\pi _\mathrm{n}(0)\right|\pi ;𝑷=1,$$
(43)
we find that $`Z_\pi ^{1/2}`$ given in Eq. (34) is the normalization factor
$$\pi _\mathrm{n}(x)=Z_\pi ^{\frac{1}{2}}g^2\overline{\mathrm{\Psi }}(x)i\gamma _5\mathrm{\Psi }(x),$$
(44)
where we have used the gap equation. Therefore, we arrive at the PCAC relation
$`_\mu j_5^\mu `$ $`=`$ $`2m_0\overline{\mathrm{\Psi }}(x)i\gamma _5\mathrm{\Psi }(x)`$ (45)
$`=`$ $`m_\pi ^2f_\pi \pi _\mathrm{n}(x).`$ (46)
Note that the decay constant (40) and the normalization factor (44) are equivalent to the previous results (Eqs. (5.25) and (5.28) in Ref. ) in the infinitely heavy mass limit of bosons $`\mu \mathrm{}`$.
### D Symmetric phase
Here we consider the symmetric phase in the chiral limit $`m_0=0`$. When $`g_0^2<g_{\mathrm{cr}}^2=2\pi ^2/\mathrm{\Lambda }^2`$, the gap equation (31) has only a trivial solution $`M=0`$. A quantity which should be zero in the broken phase is now estimated as
$$1g_0^2\frac{d^3𝒌}{(2\pi )^3}\frac{ϵ(k^+)}{k^+}=1\frac{g_0^2}{g_{\mathrm{cr}}^2}0.$$
(47)
Subsequently the factor $`\alpha `$ defined by Eq. (D20) is different from that of the broken phase \[Eq. (27)\],
$$\alpha ^1\alpha _{\mathrm{sym}}^1=1\frac{g_0^2}{g_{\mathrm{cr}}^2}+\frac{2g_0^2}{(2\pi )^3}_0^1\frac{dx}{x}d^2q_{}.$$
(48)
Then, both of the LF bound-state equation for the scalar and pseudoscalar states are given by
$$m_{\mathrm{sym}}^2\varphi _{\mathrm{sym}}(x,k_{}^i)=\frac{k_{}^2}{x(1x)}\varphi _{\mathrm{sym}}(x,k_{}^i)\frac{g_0^2\alpha _{\mathrm{sym}}}{(2\pi )^3}\frac{1}{x(1x)}_0^1𝑑yd^2\mathrm{}_{}\frac{\mathrm{}_{}^2}{y(1y)}\varphi _{\mathrm{sym}}(y,\mathrm{}_{}^i).$$
(49)
The solution to the bound-state equation is
$$\varphi _{\mathrm{sym}}(x,k_{}^i)=C_{\mathrm{sym}}\frac{g_0^2}{(2\pi )^3}\left(1\frac{g_0^2}{g_{\mathrm{cr}}^2}\right)^1\frac{1}{x(1x)k_{}^2/m_{\mathrm{sym}}^2},$$
(50)
where $`C_{\mathrm{sym}}`$ is a normalization constant and $`m_{\mathrm{sym}}=m_\pi =m_\sigma `$ is given as a solution of the equation
$$\frac{1}{g_0^2}\frac{1}{g_{\mathrm{cr}}^2}=_0^1𝑑x\frac{d^2k_{}}{(2\pi )^3}\frac{m_{\mathrm{sym}}^2}{k_{}^2m_{\mathrm{sym}}^2x(1x)}.$$
(51)
Again this is equal to the previous result of the chiral Yukawa model with $`\mu ^2\mathrm{}`$ (Eq. (5.26) in Ref. ) and therefore if we use the same cutoff as before, we obtain the same result for $`m_{\mathrm{sym}}`$. Moreover, though the above calculation was intended only to $`g_0^2<g_{\mathrm{cr}}^2`$ case, if we increase the coupling constant over its critical value $`g_{\mathrm{cr}}^2`$, we find a negative solution $`m_{\mathrm{sym}}^2<0`$. This implies that the symmetric solution causes instability when $`g_0^2>g_{\mathrm{cr}}^2`$ and thus we must choose the broken solution.
## V Summary and Conclusion
We have investigated a description of D$`\chi `$SB on the LF in the NJL model. The importance of solving the fermionic constraint for the bad spinor component was greatly stressed in analogy with the zero-mode constraint of scalar models. The exact classical solution enabled us to check the properties of the LF chiral transformation. Though the chiral transformation is differently introduced on the LF, we finally found the equivalence to the usual chiral transformation.
For a description of D$`\chi `$SB of LF NJL model, it was very important to solve the fermionic constraint nonperturbatively in quantum treatment. To do so, we introduced a bilocal formulation and solved the bilocal fermionic constraint with a fixed operator ordering by the $`1/N`$ expansion. Systematic $`1/N`$ expansion of the fermion bilocal operator is realized by the boson expansion method as the Holstein-Primakoff expansion. The leading bilocal fermionic constraint was identified with the gap equation for the chiral condensate after we took care of the infrared divergence. If we choose a nontrivial solution of the gap equation, we have a Hamiltonian in the broken phase but with a trivial vacuum.
The physical role of the fermionic constraint in the LF NJL model is very similar to that of the zero-mode constraint for scalar models. We have seen a close parallel between these two constraints. Especially it should be noted that the gap equation came from the longitudinal zero mode of the bilocal fermionic constraint.
It is very natural that we can reach the broken phase by solving the quantum fermionic constraint by $`1/N`$ expansion because the fermionic constraint is originally a part of the Euler-Lagrange equation and thus must include relevant information of dynamics. What we did is very similar to the usual mean-field approximation for the Euler-Lagrange equations. Indeed the leading order in the $`1/N`$ expansion corresponds to the mean-field approximation. However, our way of solving the fermionic constraint with the boson expansion method can easily go beyond the mean-field level. Such higher order calculation enabled us to derive a correct broken Hamiltonian and to show the divergent behavior of the (spatial integration of) pseudoscalar field.
Independence of mass from the mode expansion has both desirable and undesirable aspects. The inclusion of correct mass dependence into the IR divergent integral was required when we identify the lowest fermionic constraint with the gap equation. This is the point we must always take care of. On the other hand, the Fock vacuum is defined independent of the value of mass. Due to this fact, it is enough to have only one vacuum, namely, the Fock vacuum even in the chirally broken phase. This is the favorable aspect. However, the cost of such a simple vacuum was payed by, for example, unusual chiral transformation of fields such as $`[Q_5^{\mathrm{LF}},\overline{\mathrm{\Psi }}i\gamma _5\mathrm{\Psi }]2i\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ and non-vanishing of the LF chiral charge $`[Q_5^{\mathrm{LF}},H]0`$ in the broken phase. We found that the both effects are proportional to the dynamical fermion mass $`M`$. We also insisted the necessity of a bare mass term which accurately produced the constituent mass term. Although the special role of the fermionic constraint might be restricted to the LF NJL model, the unusual chiral transformation and the nonconservation of the chiral charge are general features of the chiral symmetry breaking on the LF. This is because they are natural consequences of the coexistence of the chiral symmetry breaking and the Fock vacuum.
The leading order eigenvalue equation for a single bosonic state is equivalent to the leading order fermion-antifermion bound-state equation. The bound-state equations were solved analytically for scalar and pseudoscalar mesons and we obtained their light-cone wavefunctions and masses. The meson masses, the decay constant, and so on were fairly consistent with those of our previous analysis on the chiral Yukawa model. The leading order calculation was limited only to two body sector (fermion and antifermion). If we consider the higher order Hamiltonian such as $`h^{(3)}`$ or $`h^{(4)}`$, we will be able to discuss four or six body sectors. In other words, since we have bosonic meson states, we can expand the Fock space in terms of the mesonic degrees of freedom. Then, for example, we will be able to discuss the mixing of scalar state and two pseudoscalar fields ($`\pi `$-$`\pi `$ mixing with $`\sigma `$).
## Acknowledgments
The authors acknowledge W. Bentz for discussions on the cutoff schemes. One of them (K.I.) is thankful to K. Harada and T. Heinzl for helpful discussions at the very early stage of this work and to K. Yazaki and K. Yamawaki for their useful comments. He is also grateful to T. Hatsuda and H. Toki for their encouragements and interests in our work.
## A Conventions
We follow the Kogut-Soper convention . First of all, the light-front coordinates are defined as
$$x^\pm =\frac{1}{\sqrt{2}}(x^0\pm x^3),x_{}^i=x^i(i=1,2),$$
(A1)
where we treat $`x^+`$ as “time”. The spatial coordinates $`x^{}`$ and $`x_{}`$ are called the longitudinal and transverse directions respectively. Derivatives in terms of $`x^\pm `$ are defined by $`_\pm =/x^\pm `$. It is useful to introduce projection operators $`\mathrm{\Lambda }_\pm `$ defined by
$$\mathrm{\Lambda }_\pm =\frac{1}{2}\gamma ^{}\gamma ^\pm =\frac{1}{\sqrt{2}}\gamma ^0\gamma ^\pm .$$
(A2)
Indeed $`\mathrm{\Lambda }_\pm `$ satisfy the projection properties $`\mathrm{\Lambda }_\pm ^2=\mathrm{\Lambda }_\pm ,\mathrm{\Lambda }_++\mathrm{\Lambda }_{}=1`$, etc. Splitting the fermion field by the projectors as
$$\mathrm{\Psi }^a=\psi _+^a+\psi _{}^a,\psi _\pm ^a\mathrm{\Lambda }_\pm \mathrm{\Psi }^a,$$
(A3)
we find that for any fermion on the LF, $`\psi _{}`$ component is a dependent degree of freedom. $`\psi _+`$ and $`\psi _{}`$ are called the “good component” and the “bad component”, respectively.
As is commented in the text, for practical calculation, we use the two-component representation for the gamma matrices. The two-component representation is characterized by a specific form of the projectors $`(\text{2})`$. Then the projected fermions $`\psi _\pm `$ have only two components. There are many possibilities which realize Eq. $`(\text{2})`$. For example, a specific representation
$$\gamma ^0=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\gamma ^3=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\gamma ^i=\left(\begin{array}{cc}i\sigma ^i& 0\\ 0& i\sigma ^i\end{array}\right)$$
(A4)
is used in Ref. . In this paper, however, we choose a representation (4) from which it is easy to extract information of the 1+1 dimensional results. Two-component spinors $`\psi `$ and $`\chi `$ are defined by Eq. (3). Results of the chiral Gross-Neveu model can be easily obtained if we make a replacement for the Pauli matrices $`\sigma _31`$ and $`\sigma ^i0`$, and regard $`\psi `$ and $`\chi `$ as one component spinors.
Using this representation, the Lagrangian density of the NJL model is written as
$``$ $`=`$ $`i\psi ^{}_+\psi +i\chi ^{}_{}\chi {\displaystyle \frac{1}{\sqrt{2}}}\left(\psi ^{}\sigma ^i_i\chi \chi ^{}\sigma ^i_i\psi \right){\displaystyle \frac{m_0}{\sqrt{2}}}\left(\psi ^{}\chi +\chi ^{}\psi \right)`$ (A6)
$`+{\displaystyle \frac{g^2}{4}}\left\{(\psi ^{}\chi +\chi ^{}\psi )^2(\psi ^{}\sigma ^3\chi \chi ^{}\sigma ^3\psi )^2\right\}.`$
## B Classical Solutions to the fermionic constraints
To solve the fermionic constraints $`classically`$ means that we treat all the fermion fields (both good and bad components) as Grassmann numbers and neglect all the $`c`$-numbers which will emerge in quantum theory under the exchange of variables.
Before discussing a complicated equation of the NJL model, it would be better to go first with the chiral Gross-Neveu model. We solve the fermionic constraint of the chiral Gross-Neveu model with the antiperiodic boundary condition:
$`\left\{i_{}+g^2a(x^{})\right\}\chi ={\displaystyle \frac{m_0}{\sqrt{2}}}\psi ,`$ (B1)
$`\chi _a(x^{}=\mathrm{})=\chi _a(x^{}=\mathrm{}),`$ (B2)
where we used a matrix notation with a matrix $`a_{ab}(x^{})\psi _a(x^{})\psi _b^{}(x^{})`$. The solution to this equation is given by
$$\chi (x^{})=_{\mathrm{}}^{\mathrm{}}𝑑y^{}G_{\mathrm{GN}}(x^{},y^{})\frac{m_0}{\sqrt{2}}\psi (y^{}),$$
(B3)
where $`G_{\mathrm{GN}}(x^{},y^{})`$ is the Green function satisfying
$`\left\{i_{}^x+g^2a(x^{})\right\}G_{\mathrm{GN}}(x^{},y^{})=\delta (x^{}y^{}),`$ (B4)
$`G_{\mathrm{GN}}(x^{}=\mathrm{},y^{})=G_{\mathrm{GN}}(x^{}=\mathrm{},y^{}).`$ (B5)
Due to Eq. (B5), the solution of course satisfies the antiperiodic boundary condition. Equation (B4) is solved as
$`G_{\mathrm{GN}}(x^{},y^{})`$ $`=`$ $`G_{\mathrm{GN}}^{(0)}(x^{})\left[{\displaystyle \frac{1}{2i}}ϵ(x^{}y^{})+C\right]G_{\mathrm{GN}}^{(0)1}(y^{}),`$ (B8)
$`G_{\mathrm{GN}}^{(0)}(x^{})=\mathrm{P}e^{ig^2_{\mathrm{}}^x^{}a(y^{})𝑑y^{}},`$
$`C={\displaystyle \frac{1}{2i}}{\displaystyle \frac{G_{\mathrm{GN}}^{(0)}(\mathrm{})G_{\mathrm{GN}}^{(0)}(\mathrm{})}{G_{\mathrm{GN}}^{(0)}(\mathrm{})+G_{\mathrm{GN}}^{(0)}(\mathrm{})}},`$
where $`G_{\mathrm{GN}}^{(0)}(x^{})`$ is a solution of a homogeneous equation $`\left\{i_{}^x+g^2a(x^{})\right\}G_{\mathrm{GN}}^{(0)}(x^{})=0`$ and the integral constant $`C`$ has been determined so that $`G_{\mathrm{GN}}(x^{},y^{})`$ satisfies the antiperiodic boundary condition. When $`N=1`$, the solution (B3) is equivalent to Domokos’ solution to the Thirring model on the light front .
In two dimensions, the LF chiral transformation is not distinguishable with the $`U(1)`$ transformation. Indeed, the “LF chiral transformation” on the good component is $`\psi e^{i\theta }\psi `$ and equivalent to the $`U(1)`$ transformation. And also, the solution (B3) implies that the bad component rotates just the same way as the good component $`\chi e^{i\theta }\chi `$.
Next let us turn to the NJL model. If we explicitly write all the indices, the fermionic constraint (1) is
$$i_{}\left(\begin{array}{c}\chi _{1a}\\ \chi _{2a}\end{array}\right)=\frac{1}{\sqrt{2}}\left(\begin{array}{c}m_0\psi _{1a}_z\psi _{2a}\\ _{\overline{z}}\psi _{1a}+m_0\psi _{2a}\end{array}\right)g^2\left(\begin{array}{c}\psi _{1a}\psi _{1b}^{}\chi _{1b}\psi _{1a}\psi _{2b}\chi _{2b}^{}\\ \psi _{2a}\psi _{2b}^{}\chi _{2b}\psi _{2a}\psi _{1b}\chi _{1b}^{}\end{array}\right),$$
(B9)
where $`_z=_1i_2`$ and $`_{\overline{z}}=_1+i_2`$. Since the equation for $`\chi _1`$ (or $`\chi _2`$) includes $`\chi _1`$ and $`\chi _2^{}`$ (or $`\chi _2`$ and $`\chi _1^{}`$), it is useful to introduce a constraint equation for $`\chi _2^{}`$ instead of $`\chi _2`$. Then we have a more tractable equation
$$i_{}\left(\begin{array}{c}\chi _{1a}\\ \chi _{2a}^{}\end{array}\right)=\frac{1}{\sqrt{2}}\left(\begin{array}{c}m_0\psi _{1a}_z\psi _{2a}\\ _z\psi _{1a}^{}+m_0\psi _{2a}^{}\end{array}\right)g^2\left(\begin{array}{cc}\psi _{1a}\psi _{1b}^{}& \psi _{1a}\psi _{2b}\\ \psi _{2a}^{}\psi _{1b}^{}& \psi _{2a}^{}\psi _{2b}\end{array}\right)\left(\begin{array}{c}\chi _{1b}\\ \chi _{2b}^{}\end{array}\right).$$
(B10)
As in the 1+1 dimensional case, the solution is immediately given if we find the Green function $`G(x^{},y^{},x_{})`$ which satisfies
$$\left\{i_{}^x+g^2𝒜(x^{})\right\}G(x^{},y^{},x_{})=\delta (x^{}y^{}),$$
(B11)
with a matrix $`𝒜_{ijab}(x)`$ defined by Eq. (5). The result is very similar to the two-dimensional result and is given by Eqs. (2) and (3) in the text.
## C Problem of the operator ordering
Here we consider the problem of operator ordering within the chiral Gross-Neveu model with $`N=1`$. Following the standard procedure, the Dirac brackets are calculated as
$`\{\psi (x),\psi ^{}(y)\}_\mathrm{D}=i\delta (x^{}y^{}),`$ (C1)
$`\{\chi (x),\psi ^{}(y)\}_\mathrm{D}=iG_{\mathrm{GN}}(x,y)\left({\displaystyle \frac{m_0}{\sqrt{2}}}g^2\psi ^{}(y)\chi (y)\right),`$ (C2)
$`\{\chi (x),\psi (y)\}_\mathrm{D}=iG_{\mathrm{GN}}(x,y)g^2\psi \chi (y),`$ (C3)
where $`G_{\mathrm{GN}}(x,y)`$ is the Green function (B8) for $`N=1`$ case. To quantize the system we simply replace the Dirac bracket $`\{A,B\}_\mathrm{D}`$ by the anticommutation relation $`i\{A,B\}`$. This procedure has ambiguity of the operator ordering.
The operator ordering we took for the fermionic constraint (1) in the NJL model corresponds to the following one in the chiral Gross-Neveu model,
$$i_{}\chi +g^2\psi \psi ^{}\chi =\frac{m_0}{\sqrt{2}}\psi .$$
(C4)
We can easily find its quantum solution due to $`[\psi \psi ^{}(x),\psi \psi ^{}(y)]=0`$. The solution is
$$\chi _{\mathrm{sol}}(x^{})=_{\mathrm{}}^{\mathrm{}}𝑑y^{}G_{\mathrm{GN}}(x,y)\frac{m_0}{\sqrt{2}}\psi (y),$$
(C5)
where $`G`$ is again the Green function (B8) with $`N=1`$.
Now let us consider the consistency for the anticommutator $`\{\chi ,\psi ^{}\}`$. It can be calculated two different ways: (i) from the solution $`\chi _{\mathrm{sol}}`$ of the fermionic constraint, and (ii) from the Dirac bracket (C2). We fix the operator ordering of the fermionic constraint by Eq. (C4) and check whether the Dirac bracket can produce the same anticommutator or not.
Instead of the anticommutator itself, we present here the calculation of a quantity $`iD_{}^x\{\chi (x),\psi ^{}(y)\}`$ where $`iD_{}^x=i_{}+g^2\psi \psi ^{}`$. Using the solution (C5), we have
$$iD_{}^x\{\chi _{\mathrm{sol}}(x),\psi ^{}(y)\}=\delta (x^{}y^{})\left(\frac{m_0}{\sqrt{2}}g^2\psi ^{}\chi \right).$$
(C6)
On the other hand, if we take the simplest ordering in the r.h.s. of the Dirac bracket (C2), we obtain
$`iD_{}^x\{\chi (x),\psi ^{}(y)\}`$ $`=`$ $`iD_{}^xG_{\mathrm{GN}}(x,y)\left({\displaystyle \frac{m_0}{\sqrt{2}}}g^2\psi ^{}\chi (y)\right)`$ (C7)
$`=`$ $`\delta (x^{}y^{})\left({\displaystyle \frac{m_0}{\sqrt{2}}}g^2\psi ^{}\chi (y)\right).`$ (C8)
This is identical with the result (C6). Therefore we find our ordering Eq. (C4) is consistent with the anticommutation relation
$$\{\chi (x),\psi ^{}(y)\}=G_{\mathrm{GN}}(x,y)\left(\frac{m_0}{\sqrt{2}}g^2\psi ^{}\chi \right).$$
(C9)
Of course if we take other operator ordering, the two results do not coincide. We expect that even in the NJL model, we can select the r.h.s. of Dirac brackets so that they coincide with the direct result with the ordering defined by Eq. (1).
## D Bilocal Fermionic constraints and their solutions by the Boson Expansion Method
It is tractable to solve the equations for $`𝒯(𝒙,𝒚)`$ and $`𝒰(𝒙,𝒚)`$ rather than $`𝒯_{\alpha \beta }(𝒙,𝒚)`$ and $`𝒰_{\alpha \beta }(𝒙,𝒚)`$. In momentum representation, the fermionic constraints for $`𝒯`$ and $`𝒰`$ are
$`q^+𝒯(𝒑,𝒒)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(iq_{}^i\sigma ^i+m_0\right)_{\alpha \beta }\left(_{\alpha \beta }(𝒑,𝒒)_{\alpha \beta }(𝒒,𝒑)\right)`$ (D3)
$`{\displaystyle \frac{g^2}{2}}{\displaystyle }{\displaystyle \frac{d^3𝒑^{}d^3𝒒^{}}{(2\pi )^3}}\{_{\alpha \beta }(𝒑,𝒒𝒑^{}𝒒^{})(\delta _{\alpha \beta }𝒯(𝒑^{},𝒒^{})+\sigma _{\alpha \beta }^3i𝒰(𝒑^{},𝒒^{}))`$
$`(\delta _{\alpha \beta }𝒯(𝒑^{},𝒒^{})\sigma _{\alpha \beta }^3i𝒰(𝒑^{},𝒒^{}))_{\alpha \beta }(𝒒𝒑^{}𝒒^{},𝒑)\},`$
$`q^+i𝒰(𝒑,𝒒)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\left\{\sigma ^3(iq_{}^i\sigma ^i+m_0)\right\}_{\alpha \beta }_{\alpha \beta }(𝒑,𝒒)+\left\{(iq_{}^i\sigma ^i+m_0)\sigma ^3\right\}_{\alpha \beta }_{\alpha \beta }(𝒒,𝒑)\right]`$ (D6)
$`{\displaystyle \frac{g^2}{2}}{\displaystyle }{\displaystyle \frac{d^3𝒑^{}d^3𝒒^{}}{(2\pi )^3}}\{_{\alpha \beta }(𝒑,𝒒𝒑^{}𝒒^{})(\sigma _{\alpha \beta }^3𝒯(𝒑^{},𝒒^{})+\delta _{\alpha \beta }i𝒰(𝒑^{},𝒒^{}))`$
$`+(\sigma _{\alpha \beta }^3𝒯(𝒑^{},𝒒^{})\delta _{\alpha \beta }i𝒰(𝒑^{},𝒒^{}))_{\alpha \beta }(𝒒𝒑^{}𝒒^{},𝒑)\}.`$
In place of the quantization condition (1), the system with bilocal operators can be characterized by the following algebra
$`[:_{\alpha \beta }(𝒑_1,𝒑_2):,:_{\gamma \delta }(𝒒_1,𝒒_2):]`$ (D7)
$`=`$ $`:_{\alpha \delta }(𝒑_1,𝒒_2):\delta _{\beta \gamma }\delta ^{(3)}(𝒑_2+𝒒_1):_{\gamma \beta }(𝒒_1,𝒑_2):\delta _{\alpha \delta }\delta ^{(3)}(𝒑_1+𝒒_2)`$ (D10)
$`+N\delta _{\alpha \delta }\delta _{\beta \gamma }\delta ^{(3)}(𝒑_1+𝒒_2)\delta ^{(3)}(𝒑_2+𝒒_1)`$
$`\times \left(\theta (p_1^+)\theta (p_2^+)\theta (q_1^+)\theta (q_2^+)\theta (p_1^+)\theta (p_2^+)\theta (q_1^+)\theta (q_2^+)\right),`$
where the normal order of $``$ was defined with respect to the Fock vacuum (3)
$$:_{\alpha \beta }^+(𝒑,𝒒):|0=:_{\alpha \beta }^+(𝒑,𝒒):|0=:_{\alpha \beta }^{++}(𝒑,𝒒):|0=0.$$
(D11)
The upper indices stand for signs of the longitudinal momenta.
The complicated structure of the algebra for the bilocal operators which originates from the fermion statistics, is greatly reduced if one introduces the boson expansion method. We can represent the operators $`::`$ in terms of bilocal boson operators $`B(𝒑,𝒒)`$ of order $`𝒪(N^0)`$ so that they fulfill the original algebra (D10). Since the algebra has a bosonic feature in the large $`N`$ limit,
$$[:_{\alpha \beta }^{++}(𝒑_1,𝒑_2):,:_{\gamma \delta }^{}(𝒒_1,𝒒_2):]N\delta _{\alpha \delta }\delta _{\beta \gamma }\delta ^{(3)}(𝒑_1+𝒒_2)\delta ^{(3)}(𝒑_2+𝒒_1),$$
it would be better to choose a representation which satisfies thisActually there are many possibilities to express Eq. (D10) in terms of bosonic operators, corresponding to various “local expansions” of the Grassmannian manifold of the bilocal operators.. The Holstein-Primakoff type expansion satisfies the requirement.
Physically this procedure corresponds to extracting purely bosonic degrees of freedom from a fermion-antifermion system, i.e., a mesonic system. The power of the boson expansion method in the light-front field theories was first demonstrated by one of the authors . He applied the Holstein-Primakoff type expansion to 1+1 dimensional QCD and derived an effective Hamiltonian for mesons as local bosons whose masses are given by the ’t Hooft equation. Using the effective Hamiltonian, we can in principle study, say, scattering of mesons as $`q\overline{q}`$ bound states.
Since the essential structure of the algebra (D10) is determined only by the longitudinal momentum, it is straightforward to apply the Holstein-Primakoff expansion discussed in Ref. to four dimensional fermionic theories. Indeed the operators $`::`$ are represented as follows:
$`:_{\alpha \beta }^+(𝒑_1,𝒑_2):`$ $`=`$ $`{\displaystyle [d𝒒]\underset{\gamma }{}B_{\alpha \gamma }^{}(𝒑_1,𝒒)B_{\beta \gamma }(𝒑_2,𝒒)}A_{\beta \alpha }(𝒑_2,𝒑_1),`$ (D12)
$`:_{\alpha \beta }^+(𝒑_1,𝒑_2):`$ $`=`$ $`{\displaystyle [d𝒒]\underset{\gamma }{}B_{\gamma \beta }^{}(𝒒,𝒑_2)B_{\gamma \alpha }(𝒒,𝒑_1)},`$ (D13)
$`:_{\alpha \beta }^{++}(𝒑_1,𝒑_2):`$ $`=`$ $`{\displaystyle [d𝒒]\underset{\gamma }{}(\sqrt{NA})_{\beta \gamma }(𝒑_2,𝒒)B_{\gamma \alpha }(𝒒,𝒑_1)},`$ (D14)
$`:_{\alpha \beta }^{}(𝒑_1,𝒑_2):`$ $`=`$ $`{\displaystyle [d𝒒]\underset{\gamma }{}B_{\gamma \beta }^{}(𝒒,𝒑_2)(\sqrt{NA})_{\gamma \alpha }(𝒒,𝒑_1)}.`$ (D15)
These give the $`1/N`$ expansion of $`_{\alpha \beta }(𝒑,𝒒)`$. The first few terms are shown in the text \[Eqs. (14)-(17)\].
If we expand the bilocal operators $`𝒯(𝒑,𝒒),𝒰(𝒑,𝒒)`$, and $`_{\alpha \beta }(𝒑,𝒒)`$, the equation for order $`n`$ can be written in a compact form:
$$\left(\begin{array}{c}t^{(n)}(𝒑,𝒒)\\ u^{(n)}(𝒑,𝒒)\end{array}\right)=\left(\begin{array}{c}F^{(n)}(𝒑,𝒒)\\ G^{(n)}(𝒑,𝒒)\end{array}\right)g_0^2\frac{ϵ(p^+)}{q^+}_{\mathrm{}}^{\mathrm{}}\frac{d^3𝒌}{(2\pi )^3}\left(\begin{array}{c}t^{(n)}(𝒌,𝒑+𝒒𝒌)\\ u^{(n)}(𝒌,𝒑+𝒒𝒌)\end{array}\right),$$
(D16)
where quantities $`F^{(n)}(𝒑,𝒒)`$ and $`G^{(n)}(𝒑,𝒒)`$ are generally complicated functions of bilocal operators except for the lowest order \[see Eq. (23)\]. For example, $`F^{(1)}`$ and $`G^{(1)}`$ are
$`F^{(1)}(𝒑,𝒒)={\displaystyle \frac{1}{2q^+}}(iq_{}^i\sigma ^i+M)_{\alpha \beta }\left(\mu _{\alpha \beta }^{(1)}(𝒑,𝒒)\mu _{\alpha \beta }^{(1)}(𝒒,𝒑)\right),`$ (D17)
$`G^{(1)}(𝒑,𝒒)={\displaystyle \frac{i}{2q^+}}\left[\left\{\sigma _3(iq_{}^i\sigma ^i+M)\right\}_{\alpha \beta }\mu _{\alpha \beta }^{(1)}(𝒑,𝒒)+\left\{(iq_{}^i\sigma ^i+M)\sigma _3\right\}_{\alpha \beta }\mu _{\alpha \beta }^{(1)}(𝒒,𝒑)\right],`$ (D18)
where $`\mu _{\alpha \beta }^{(1)}(𝒑,𝒒)`$ is given by the boson expansion method Eq. (15). Since all of the orders of the operators are less than $`n`$, we can solve this equation order by order. The solution of this integral equation is
$$\left(\begin{array}{c}t^{(n)}(𝒑,𝒒)\\ u^{(n)}(𝒑,𝒒)\end{array}\right)=\left(\begin{array}{c}F^{(n)}(𝒑,𝒒)\\ G^{(n)}(𝒑,𝒒)\end{array}\right)g_0^2\frac{ϵ(p^+)}{q^+}\alpha (p^++q^+)_{\mathrm{}}^{\mathrm{}}\frac{d^3𝒌^{}}{(2\pi )^3}\left(\begin{array}{c}F^{(n)}(𝒌^{},𝒑+𝒒𝒌^{})\\ G^{(n)}(𝒌^{},𝒑+𝒒𝒌^{})\end{array}\right),$$
(D19)
where
$$\alpha (P^+)=\left(1+g_0^2_{\mathrm{}}^{\mathrm{}}\frac{d^3𝒌}{(2\pi )^3}\frac{ϵ(k^+)}{P^+k^+}\right)^1.$$
(D20)
The quantities $`t^{(2)}(𝒑,𝒒)`$ and $`u^{(2)}(𝒑,𝒒)`$ are necessary for obtaining a correct Hamiltonian of the system.
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# hep-ph/0004130April 2000 OUTP–00–15P A Brane–World Explanation of the KARMEN Anomaly
## 1 Introduction
Two recent developments, namely the discovery of brane–world models and the flexibility in choosing fundamental scales and sizes of additional dimensions , have led to considerable activity exploring new options for particle phenomenology. Both developments have a natural place in string– or M–theory which leads to interesting connections between the leading candidate for a fundamental theory and particle physics.
Much of the interesting phenomenology of brane–world models is associated with the Kaluza–Klein modes that originate from large, gravity–only additional dimensions. Particularly the graviton and its Kaluza–Klein excitations are of interest in this context . Neutrino physics is another branch that is likely to be effected in brane–world models. This is because higher–dimensional bulk fermions lead to Kaluza–Klein towers of standard model singlets that may be interpreted as sterile neutrinos . Most of the work in brane–world neutrino physics to date has been focusing on rather specific models. Concretely, most of the explicit studies have been carried out for simple five–dimensional models with a special form of the action that, for example, does not include bulk mass terms (see, however, ref. , where Majorana mass terms have been considered in the context of orbifold models).
Brane–world models with large additional dimensions have led to a number of exciting predictions, such as the modification of gravity at small distances , that may be testable in the future. In this paper, we would like to pursue the idea that a signal originating from brane–world physics may have already been observed. Specifically, we will show that the anomaly in the KARMEN neutrino experiment can be explained in the context of brane–world models. This will also lead us to explore new possibilities for brane–world models of neutrino masses. That is, we will examine a number of new model building aspects in brane–world theories of neutrino masses.
The KARMEN neutrino experiment studies neutral current interactions of neutrinos that originate from the decay chain $`\pi ^+\mu ^+`$ at rest. The KARMEN experiment has a unique time structure that allows one to study the number of neutrino events in the detector as a function of time after the production of the initial $`\pi ^+`$. It is this distribution that shows an anomalous peak at a specific time after $`\pi ^+`$ production. The KARMEN anomaly can be explained in terms of a hypothetical $`X`$ particle produced in a rare decay $`\pi ^+\mu ^+X`$. It has also been shown that the $`X`$ particle can be interpreted as a sterile neutrino. While postulating such a particle leads to a satisfactory fit of the experimental data it has a theoretically unappealing feature. To match the observed peak, the $`X`$ particle has to have a very specific mass $`m_X`$ that is fine–tuned to the mass difference $`m_\pi m_\mu `$ with an accuracy of approximately $`10^4`$. That is, $`m_X/(m_\pi m_\mu )1\mathrm{1.8\hspace{0.17em}10}^4`$. There is no apparent theoretical explanation for such a coincidence of masses.
Here is where brane–world physics might come into play. Suppose, that a brane–world model leads to a Kaluza–Klein tower of sterile neutrinos. Clearly, if the spacing is sufficiently small, it is relatively more probable for one of the particles in this tower to fall into the narrow mass range relevant for the KARMEN experiment than it is for a single particle. In the extreme case, where the spacing of the tower is of the order of the mass gap available for the $`X`$ particle, the fine–tuning problem is solved completely.
In this paper, we will show that a brane–world model explaining the KARMEN anomaly does indeed exist. An analysis of the basic requirements singles out two specific choices of dimensions and scales. The first one corresponds to a five–dimensional model with intermediate string scale and the scale of the fifth dimension in the $`\mathrm{keV}`$ range. The second one corresponds to a six–dimensional model with a $`\mathrm{TeV}`$ string scale and the scale of the two additional dimensions in the $`\mathrm{MeV}`$ range. We will focus on the second case and construct an explicit six–dimensional model that meets all the requirements. A major phenomenological constraint on this model is that the other Kaluza–Klein sterile neutrinos do not distort the KARMEN spectrum in a way inconsistent with the experiment. We will show that this can be avoided by making the spacing of the Kaluza–Klein tower sufficiently large. Unfortunately then, the fine–tuning problem cannot be solved completely in the way indicated above but is still improved by two orders of magnitude with respect to the single particle interpretation. At the same time, the effect of the Kaluza–Klein modes lighter than the $`X`$ particle, provide a way of how our proposal could be experimentally distinguished from the single particle proposal. These particles contribute to the KARMEN spectrum at short times and, depending on the parameters of the model, the corresponding signal may be detectable in future experiments.
Our six–dimensional model also shows a number of interesting properties that are of general importance for brane–world models of neutrino masses. In the simplest version of the model with conserved lepton number, for example, we find three massless eigenstates that are generally nontrivial linear combinations of the electroweak eigenstates and the bulk neutrinos. The admixture of bulk neutrinos is controlled by a six–dimensional Dirac mass and can be made small or large while keeping the states massless. In contrast, in the models considered so far, massless modes were either not present or decoupled from the ordinary neutrinos. Furthermore, in those models, light neutrino masses and their mixing with bulk neutrinos are usually controlled by the same parameters. Our model shows that it is actually possible to decouple these two quantities. More generally, it is illustrated that the inclusion of all Lorenz–invariant (or even Lorenz non–invariant) mass terms in the higher–dimensional theory opens up new possibilities for brane–world neutrino phenomenology.
## 2 The KARMEN anomaly and its brane–world interpretation
The KARMEN experiment at the Rutherford Appleton Laboratory studies charged and neutral current interactions of neutrinos from the $`\pi ^+\mu ^+`$ decay chain at rest. The unique feature of the neutrino flux is its time structure. The primary pions are produced in $`0.5\mu \mathrm{s}`$ long pulses<sup>1</sup><sup>1</sup>1More precisely, the $`0.5\mu \mathrm{s}`$ pulses are made of two shorter $`100\mathrm{ns}`$ long pulses separated in time by $`330\mathrm{ns}`$. with a frequency of $`50\text{Hz}`$ by the ISIS synchrotron proton beam. Due to the short pion lifetime, $`\tau _\pi 26\mathrm{ns}`$, the muon neutrinos from the $`\pi ^+\mu ^+\nu _\mu `$ decay represent the prompt component of the neutrino beam, which exhausts itself within a few $`\tau _\pi `$ after the pion pulse. Hence, in the time window $`(0.6`$$`10.6)\mu \mathrm{s}`$ after proton beam on target the neutrino beam is exclusively composed out of $`\nu _e`$ and $`\overline{\nu }_\mu `$ that originate from the slow decay ($`\tau _\mu 2.2\mu \mathrm{s}`$) of the muons produced in the first $`0.5\mu \mathrm{s}`$. Consequently, the time spectrum of the number $`n`$ of detector events in this time window is expected to be described by
$$\frac{dn}{dt}=Ae^{t/\tau _\mu }+B.$$
(1)
Here $`A`$ is a time–independent<sup>2</sup><sup>2</sup>2The convolution of the muon exponential decay rate with the time spectrum of muon production is again an exponential with the same time constant. rate constant depending, among other things, on the muon production rate, on the target cross section and on the detector efficiency. The constant $`B`$ represents the event rate associated with a time–independent background.
The KARMEN time–anomaly manifests itself in a discrepancy between the expected time spectrum (1) and the measured one. This discrepancy is mainly due to a peak at approximately $`3.6\mu \mathrm{s}`$ after beam on target. Such a peak is clearly visible in Fig. 1a. In this plot the time window $`(0.6\text{}\mathrm{\hspace{0.17em}10.6})\mu \mathrm{s}`$ has been divided into 20 bins and the number of events for each bin (falling into the energy range from $`10`$ to $`36\mathrm{MeV}`$) has been indicated by a dot. In order to estimate the significance of the discrepancy and for later comparison with results based on our model we have performed a simple analysis of the data shown in Fig. 1a. Specifically, we have carried out a $`\chi ^2`$ fit of the data set in Fig. 1 assuming an event rate as in eq. (1). Background measurements on a wider time interval fix the parameter $`B`$ as indicated by the dashed line in Fig. 1. The best fit results in $`\chi _{\text{min}}^2=29`$ for 19 degrees of freedom, corresponding approximately to a $`2\sigma `$ discrepancy. However, the significance of the discrepancy is higher if one considers only the two bins around $`t=3.6\mu \mathrm{s}`$. This is appropriate because an explanation of the anomaly around $`t=3.6\mu \mathrm{s}`$ by statistical fluctuation is strongly disfavored. The peak around $`t=3.6\mu \mathrm{s}`$, already clearly present in the 1995 data , has, in fact, been confirmed by the subsequent KARMEN2 data . Once the two bins around $`t=3.6\mu \mathrm{s}`$ have been excluded, the remaining data is very well described by the distribution in eq. (1). The best fit leads to $`\chi _{\text{min}}^2=14`$ with 17 degrees of freedom, corresponding to the solid line in Fig. 1a. According to such a fit, a total of 557 events are expected in the two unfitted bins. As shown in , the measured 659 events represent a deviation at the level of approximately $`4\sigma `$. Since the time–anomaly appears to be statistically significant and any attempt to explain its peak structure by systematic effects failed so far, it is worth to speculate about its physical origin.
A simple description of the time spectrum can be achieved by assuming that the peak around $`t=3.6\mu \mathrm{s}`$ is associated with the decay of a slowly moving massive neutral particle “$`X`$” in the detector. Such a particle could, for example, be produced in the rare decay $`\pi ^+\mu ^+X`$ . The position of the peak then measures the time of flight of the particle and hence its velocity. According to the recent KARMEN fits the velocity is given by $`\beta _X0.0162`$. This velocity, in turn, determines the time the particle spends in the detector and, together with the time structure of the proton beam, the width of the signal. This width is found to be in nice agreement with the observed width of the peak. More generally, the slow particle interpretation is supported by a detailed analysis of the whole data set. Due to the very good time and spatial resolution of the detector, $`\mathrm{\Delta }t<2\mathrm{ns}`$, $`\mathrm{\Delta }x`$, $`\mathrm{\Delta }y`$, $`\mathrm{\Delta }z<15\text{cm}`$, an analysis of the correlation between time and position of the events is possible. Such an analysis is in agreement with the slow particle hypothesis and provides the most precise measurement of the $`X`$ particle velocity leading to the value given above .
Once the signal is attributed to the decay of a slow particle, a rare $`\pi ^+\mu ^+X`$ decay represents the simplest possibility for the production of the particle. Since the time width of the anomaly is smaller than the muon lifetime, it is natural to assume that its production is not associated to muon decay (see however ). Moreover, the production in the $`\pi ^+`$ decay allows to account for the slow, fixed velocity in terms of a small available phase space (unlike, for example, the production by the proton interactions with the target). Finally, the energy deposited by the particle in the detector is found to be less than about $`35\mathrm{MeV}`$ . This is what is expected from the decay of a particle produced in the decay $`\pi ^+\mu ^+X`$, whose mass $`m_X`$ has to be less than $`m_\pi m_\mu 33.9\mathrm{MeV}`$. On the other hand, if the production process were $`\pi ^+e^+X`$, the slow $`X`$ particle would have to have a mass $`m_X137\mathrm{MeV}`$ and would release more energy than measured (unless having a muon or another massive particle in the final state is for some reason favored). As for the decay, the simplest possibilities for the visible decay of a particle produced in the $`\pi ^+\mu ^+X`$ channel are $`Xe^+e^{}\nu `$ and $`X\gamma \nu `$, where $`\nu `$ generically represents a neutral light particle. A visible decay mainly through $`X\gamma \nu `$ thereby producing mono–energetic photons is highly disfavored . Consequently, in this paper we will consider the possibility that the $`X`$ particle is produced in the decay $`\pi ^+\mu ^+X`$ decay and the anomaly is mainly due to its visible decay $`Xe^+e^{}\nu `$. If the $`X`$ particle lifetime is larger than $`3.6\mu \mathrm{s}`$, as it will be in our case, the KARMEN data requires that
$$\text{BR}(\pi ^+\mu ^+X)\mathrm{\Gamma }_{\text{vis}}=(1\text{}\mathrm{\hspace{0.17em}2})\mathrm{\hspace{0.17em}10}^{11}\text{s}^1,$$
(2)
where $`\text{BR}(\pi ^+\mu ^+X)`$ is the branching ratio of the decay $`\pi ^+\mu ^+X`$ and $`\mathrm{\Gamma }_{\text{vis}}`$ is the partial width of the visible $`X`$ decays.
The first proposal for the nature of the hypothetical $`X`$ particle was made in . There, it was identified with a sterile neutrino that mixes with the Standard Model (SM) neutrinos. The same possibility was considered in . Supersymmetric scenarios were studied in . From a theoretical point of view, light sterile neutrinos are welcome if their lightness can be accounted for. Whereas the mass of the SM neutrinos is protected by the electroweak symmetry (and the absence of fundamental isovector Higgses), the SM symmetries alone cannot explain why sterile neutrinos should be lighter than, say, the Planck mass. It is therefore useful to consider sterile neutrinos in the context of a specific framework. In this paper, we will be dealing with sterile neutrinos that arise in brane–world models with large additional dimensions. Concretely, we consider the possibility that the KARMEN anomaly is due to a sterile neutrino that is part of a tower of Kaluza–Klein excitations associated with a SM–singlet fermion propagating in $`4+d`$ dimensions. A possible explanation for small sterile neutrino masses is available if these models are viewed in the context of string theory. Then, bulk fermions arise naturally as superpartners of moduli fields. Perturbatively, moduli are flat directions of the theory, and, hence, their fermionic partners are massless before supersymmetry breaking. After including supersymmetry breaking effects those fermions will receive a mass related to the supersymmetry breaking scale. Depending on the details, this might well provide a reason for small bulk fermion masses. The next sections will be devoted to a detailed study of such models, particularly in six dimensions. In the remainder of this section, we will describe some general features of our six–dimensional model from a phenomenological viewpoint and summarize its implications for the KARMEN experiment.
Let us consider six–dimensional Dirac fermions as the origin of the sterile neutrinos. For simplicity we will consider only one such Dirac fermion here, although our later treatment will be more general. From a four–dimensional perspective, this field can be described by a Kaluza–Klein tower of sterile Dirac neutrinos $`\nu _n`$ labeled by two integer numbers $`n=(n_1,n_2)`$. Assuming, for simplicity, that the two extra dimensions have equal size $`2\pi R`$, the mass spectrum of the Kaluza–Klein tower is given by by
$$m_n=\sqrt{\mu ^2+\frac{n^2}{R^2}},$$
(3)
where $`n^2=n_1^2+n_2^2`$ and the mass parameter $`\mu `$ originates from a six–dimensional Dirac mass term. Of course, to have a candidate for the $`X`$ particle among these Kaluza–Klein fermions we have to require that the mass parameter $`\mu `$ is smaller than mass $`m_X`$ of the $`X`$ particle. This requirement precisely corresponds to the problem of keeping sterile neutrinos light. With the above string theory explanation for this in mind, we assume that $`\mu <m_X`$. Then, we may have a specific sterile neutrino $`X=\nu _{n_X}`$ in the Kaluza–Klein tower that we would like to associate with the KARMEN anomaly. From eq. (3) we can compute the density of states $`\rho (m_X)`$ at the scale $`m_X`$, defined as the average number of states $`dn`$ per mass interval $`dm`$. Neglecting the mass $`\mu `$ for an order of magnitude estimate, this density can be computed from eq. (3) to be <sup>3</sup><sup>3</sup>3In our simplified model the $`n`$–th state is eight times degenerate if $`n_10`$, $`n_20`$, $`n_1n_2`$. In a general situation where the size of the two extra dimensions are different and all possible six–dimensional mass terms are included, the degeneracy is lifted.
$$\rho (m_\mathrm{X})2\pi R^2m_\mathrm{X}.$$
(4)
In other words, the average separation between two subsequent states at $`m_\mathrm{X}`$ in a mass–ordered list is given by
$$\mathrm{\Delta }m_{\text{KK}}=\frac{1}{\rho (m_\mathrm{X})}\frac{1}{2\pi R^2m_\mathrm{X}}.$$
(5)
The last equation highlights an interesting feature of our framework. Besides being able to account for the existence of light sterile neutrinos, it allows us to alleviate the fine–tuning problem associated with the slow particle interpretation of the KARMEN anomaly. The fine–tuning problem arises because the slowness of the $`X`$ particle requires its mass to be only slightly lower than the difference $`m_\pi m_\mu 33.9\mathrm{MeV}`$. More precisely, in order to be produced with a velocity $`\beta _X0.0162`$, the $`X`$ particle mass $`m_X`$ must be
$$m_X=m_\pi m_\mu \delta m_X,$$
(7)
where
$$\delta m_X\beta _X^2(m_\pi m_\mu )\frac{m_\pi }{2m_\mu }6\mathrm{keV}.$$
(8)
This means that $`m_X`$ is fine tuned to the $`m_\pi m_\mu `$ mass difference with a precision of $`1.8\times 10^4`$! On the other hand the fine tuning problem is significantly alleviated in the approach we propose, since the $`X`$ particle is just one of the many excitation of the six–dimensional fermion. It is therefore relatively more probable that one out of the many states has a mass that falls into the critical range.
In principle, this framework allows us to completely eliminate the fine tuning problem. In fact, the problem disappears if the average separation between subsequent states is about twice the mass range available for the $`X`$ particle, that is, $`\delta m_X=m_\pi m_\mu m_X\mathrm{\Delta }m_{\text{KK}}/2`$. From eq. (5) this happens if
$$\frac{1}{R}=(2\pi m_X\mathrm{\Delta }m_{\text{KK}})^{1/2}(4\pi m_X\delta m_X)^{1/2}1.6\mathrm{MeV}.$$
(9)
In this case, the production of the particles heavier than $`X`$ in the $`\pi ^+\mu ^+\nu _n`$ decay would be kinematically forbidden<sup>4</sup><sup>4</sup>4This assumes that there is no local inhomogeneity in the density of states, so that all states heavier than $`X`$ are indeed heavier by a margin larger than the mass gap $`\delta m_X`$ of the X particle.. However, the states lighter than $`m_\mathrm{X}`$, being faster than the $`X`$ particle, modify the time spectrum in the window $`(0.6\text{}\mathrm{\hspace{0.17em}3.6})\mu \mathrm{s}`$. In order compute the effect of those lighter states on the time spectrum, one has to take into account that they are produced more copiously due to the larger phase space available. However, their signal is shorter because they spend less time inside the detector. A simulation taking into account the detailed time–shape of the signal associated with the lighter particles shows that the KARMEN data is in very good agreement with the mass spectrum associated with an inverse radius of about $`15\mathrm{MeV}`$, namely about 10 time larger than the natural value in (9)<sup>5</sup><sup>5</sup>5The simulation assumes that the degeneracy mentioned in footnote 3 is completely lifted.. Fig. 1b shows the corresponding best fit (solid line). Once the mass spectrum has been fixed by the choice of $`1/R`$, the overall normalization of the sterile neutrinos signal and the constant $`A`$ in eq. (1) are fitted to the data. The best fit shown in Figure has a $`\chi _{\text{min}}^2=15.8`$ for 18 degrees of freedom. The fine tuning problem is alleviated. The maximum density of states at the $`m_X`$ scale is about 40 times larger than $`(m_\pi m_\mu )^1`$, making this framework about 40 times more natural than the standard framework with a single particle. However, a significant fine tuning is still needed.
The discussion above points out a peculiar feature of our proposal. The $`X`$ particle comes with a number of particles lighter than $`m_X`$ whose number depends on the radius of the additional dimensions. For large values of the inverse radius, the lighter states are separated by a large mass gap. In this case our model is indistinguishable from models with a single sterile neutrino as far as the KARMEN experiment is concerned. However, from the theoretical point of view, it still provides a well motivated scenario for the existence of a sterile neutrino with a mass of $`33.9\mathrm{MeV}`$ which, at the same time, can account for the standard oscillation phenomenology (see Section 4). On the other hand, if the average mass gap $`\mathrm{\Delta }m_{\text{KK}}`$ between the states is smaller than about $`5\mathrm{MeV}`$, the particles lighter than $`X`$ could give a small contribution to the early KARMEN time spectrum. For example, with the parameters used for the fit shown in Fig. 1b, corresponding to an average separation of about $`1\mathrm{MeV}`$ between states, one gets a small contribution to the number of events in the first time bin. Therefore, a detailed study of the time distribution at early times after the end of the proton pulse could test the model if the mass separation $`\mathrm{\Delta }m_{\text{KK}}`$ happens not to be too large. An upgraded detector with tracking capability devoted to the study of the anomaly could certainly explore this possibility. Needless to say, the detection of a signal at early times would represent a strong hint for the existence of a brane–world.
## 3 Scales and couplings
In this section, we would like to analyze the scales and couplings of the brane–world theory that are required for a solution of the KARMEN anomaly along the lines explained above. This analysis will provide us with an overview over the various models that may be suitable for our purpose. In the following section, we will pick one of the models found in this way and develop it in detail.
We start with a ten–dimensional string theory with string scale $`M_\mathrm{s}`$. The six internal dimensions are schematically divided into two groups. The first group of $`d`$ dimensions are the ones that give rise to Kaluza–Klein neutrinos potentially relevant for the KARMEN experiment. The other $`6d`$ dimensions are those without any direct relevance for the experiment. For example, a particular dimension is not relevant in this sense if its inverse size is much larger than the mass of the $`X`$ particle. For simplicity we assume that the first $`d`$ dimensions have the same size characterized by the radius $`R`$. This is likewise assumed for the remaining $`6d`$ dimensions where the corresponding radius is called $`\rho `$. Given this setup, there are three basic requirements to be satisfied. First, the four–dimensional Planck constant $`M_{\text{Pl}}`$ specified by
$$\left(\frac{M_{\text{Pl}}}{M_\mathrm{s}}\right)^2=16\pi (2\pi RM_\mathrm{s})^d(2\pi \rho M_\mathrm{s})^{6d}$$
(10)
should have the correct value. Except for this general condition we have two more constraints that are specific to our problem. The first one is related to the probability of having a Kaluza–Klein particle with mass in the relevant range for the slow–particle interpretation of the KARMEN experiment. We would like to significantly reduce the amount of fine tuning required. This means that the parameter $`\mathrm{\Delta }`$ defined as
$$\frac{1}{\mathrm{\Delta }}=\rho (m_\mathrm{X})\mathrm{\Delta }m_{\text{KK}}$$
(11)
should be significantly smaller than $`𝒪\left(10^4\right)`$, corresponding to the fine–tuning of the single–particle solution. We remind the reader that $`\rho (m_\mathrm{X})`$ is the density of Kaluza–Klein states at $`m_X=33.9\mathrm{MeV}`$ and $`\mathrm{\Delta }m_{\text{KK}}6\mathrm{keV}`$ is the mass gap between the $`X`$ particle mass and the $`\pi `$-$`\mu `$ mass difference. Note that, from the distinction of dimensions made above, only the Kaluza–Klein modes associated to $`d`$ dimensions are relevant for this density. For $`N`$ Dirac spinors in $`4+d`$ dimensions with a Dirac mass $`\mu `$ it is given by
$$\rho (m_\mathrm{X})=\frac{d\mathrm{\hspace{0.17em}2}^{[d/2]}\pi ^{d/2}}{\mathrm{\Gamma }(d/2+1)}NR^dm_\mathrm{X}(m_\mathrm{X}^2\mu ^2)^{d/21}.$$
(12)
The final requirement originates from the constraint (2) on the branching ratio and the visible decay width of the $`X`$ particle. This constraint implies that there should be a non–negligible mixing between the $`X`$ particle and the left–handed neutrinos. In the models under consideration such a mixing is generated by a mass term between brane and bulk fields that resides on the brane. On dimensional grounds the associated mass parameter $`\widehat{m}`$ can be written as
$$\widehat{m}=\frac{hv}{(2\pi RM_\mathrm{s})^{d/2}}=\frac{hm_t}{h_t(2\pi RM_\mathrm{s})^{d/2}},$$
(13)
where $`h`$ is a brane–bulk Yukawa coupling, $`v`$ is the the standard model Higgs VEV and $`h_t`$ and $`m_t`$ are the top Yukawa–coupling and mass. Then, the mixing matrix element $`U_{\nu X}`$ between the left–handed neutrinos and the $`X`$ particle specified by
$$U_{\nu X}=\frac{\widehat{m}}{m_X}$$
(14)
should not be too small. Typically, from eq. (2) one has $`|U_{\nu X}|10^3`$.
Altogether, we now have three constraints for our three fundamental parameters $`M_\mathrm{s}`$, $`R`$ and $`\rho `$. All other quantities appearing in those constraint are either fixed or in a well–defined range. We can therefore simply solve for $`M_\mathrm{s}`$, $`R`$ and $`\rho `$ as a function of the number of dimensions $`d`$. The fact that we can basically determine all fundamental scales is quite remarkable and shows how constrained the problem is. A priori, it is not clear at all that a solution with sensible values of the scales exists. However, as Table 1 shows, this is indeed the case. The quantity $`y`$ appearing in the table is defined as
$$y=\frac{h}{10^3|U_{\nu X}|h_t}.$$
(15)
Let us discuss the various cases shown in table 1. For one additional dimensions $`d=1`$, we have a model with intermediate string scale, one large additional dimension in the $`\mathrm{keV}`$ range and the size of the other five dimensions close to the string scale. For two additional dimensions, $`d=2`$, the situation is quite different. We have a model with the string scale in the $`\mathrm{TeV}`$ range, two large additional dimensions in the $`\mathrm{MeV}`$ range and the scale of the other four internal dimensions being about two to three orders of magnitude larger. In the case of three additional dimensions, $`d=3`$, the typical string scale is already quite low and one needs a large factor $`N\mathrm{\Delta }y^2`$ to elevate it above the required $`\mathrm{TeV}`$ limit. Therefore, this case only represents a marginal possibility<sup>6</sup><sup>6</sup>6Moreover, unlike for the other two cases, the scale $`\rho ^1`$ is smaller than the mass $`m_\mathrm{X}`$ of the $`X`$ particle. Therefore, if the Kaluza–Klein neutrinos propagated in those dimensions they will be relevant for the experiment, contrary to our initial definition of those dimensions. Therefore, one has to assume that this is not the case. Such an assumption is probably not very appealing from the perspective of model–building.. Finally, one can show that models with more than three additional dimensions, $`d>3`$, lead to a string scale below $`\mathrm{TeV}`$ and are, therefore, not viable in our context.
In summary, we have isolated two cases which may lead to an explanation of the KARMEN anomaly in the context of a brane–world model. The first case represents an effective five–dimensional model with intermediate string scale and the remaining five dimensions being close to the string scale. The second one is a model with a $`\mathrm{TeV}`$ string scale, two large dimensions in the $`\mathrm{MeV}^1`$ range and the scale of the remaining four dimensions being two to three orders of magnitude larger. Therefore, this model is effectively six–dimensional in an intermediate energy range. We stress that the existence of these two possibilities, particularly the one with a $`\mathrm{TeV}`$ string scale, was by no means guaranteed and is due to a rather fortunate conspiracy of the various mass scales in the problem. While explicit examples can be constructed in both cases, in this paper we will focus on the six–dimensional case with a $`\mathrm{TeV}`$ string scale.
## 4 An explicit six–dimensional example
In this section, we would like to present an explicit brane–world model that realizes the ideas outlined above.
It is clear from the above discussion of scales that it is sufficient for our purpose to construct an effective six–dimensional brane–world model, valid for energies below $`\rho ^1`$. The Kaluza–Klein modes associated to the remaining four internal dimensions are too heavy to be relevant within our context. More specifically, we would like to couple a six–dimensional bulk with right–handed neutrinos and a three–brane that carries the standard–model fields. Such a model should then be analyzed in detail with respect to its consequences for the KARMEN anomaly. Furthermore, as we will see, the model illustrates some general features that arise in brane–world models of neutrino masses which have not been considered so far.
The six–dimensional bulk action of the model is specified by
$$S_{\mathrm{bulk}}=d^4xd^2y\left[\overline{\mathrm{\Psi }}_I\mathrm{\Gamma }^Ai_A\mathrm{\Psi }_I(\mu _{IJ}\overline{\mathrm{\Psi }}_{LI}\mathrm{\Psi }_{RJ}+\text{h.c.})\right].$$
(16)
Focusing on the relevant Yukawa couplings between bulk and brane fields, the brane action reads
$$S_{\mathrm{brane}}=_{\{y=0\}}d^4x\left[\frac{h_{aIi}}{M_\mathrm{s}}\overline{\mathrm{\Psi }}_{aI}L_iH+\text{h.c.}\right].$$
(17)
We use coordinates $`x^A`$ with indices $`A,B,\mathrm{}=0,\mathrm{},5`$ for the total six–dimensional space–time. Four–dimensional coordinates $`x^\mu `$ are indexed by $`\mu ,\nu ,\mathrm{}=0,1,2,3`$. The remaining two coordinates are denoted by $`y^\alpha `$, where $`\alpha ,\beta ,\mathrm{}=1,2`$. The signature of our metric is “mostly minus”. We have introduced $`N`$ six–dimensional Dirac fermions $`\mathrm{\Psi }_I`$, where $`I,J,\mathrm{}=1,\mathrm{},N`$, in the bulk. Their left– and right–handed components are defined in the usual way by $`\mathrm{\Psi }_{R/L,I}=(1\pm \mathrm{\Gamma }_7)\mathrm{\Psi }_I`$ where $`\mathrm{\Gamma }_A`$ are the six–dimensional gamma matrices and $`\mathrm{\Gamma }_7=\mathrm{\Gamma }_0\mathrm{}\mathrm{\Gamma }_5`$. The $`4+2`$ decomposition of these gamma matrices can be written in the form $`\mathrm{\Gamma }_I=\{\gamma _\mu \mathrm{𝟏}_2,i\gamma _5\sigma _\alpha \}`$. Here $`\gamma _\mu `$ are the four–dimensional Dirac matrices and $`\gamma _5=i\gamma _0\gamma _1\gamma _2\gamma _3`$. The two–dimensional Dirac–matrices $`\sigma _\alpha `$ can be identified with the first two Pauli matrices. Defining the two–dimensional chirality operator by $`\sigma _3i\sigma _1\sigma _2`$ we have the relation $`\mathrm{\Gamma }_7=\gamma _5\sigma _3`$ between six–, four– and two–dimensional chiralities. For later purposes we also define six–dimensional charge conjugation by $`\mathrm{\Psi }_I^c=C^1\overline{\mathrm{\Psi }}_I^T`$, where the charge conjugation matrix $`C`$ is specified by the relations $`(\mathrm{\Gamma }_A)^T=C\mathrm{\Gamma }_AC^1`$, as is usual in even dimensions. Furthermore, we have introduced Dirac mass terms with associated mass matrix $`\mu _{IJ}`$ in the bulk. By a suitable redefinition of the bulk fermions we can diagonalize this mass matrix. In the following, we will use this diagonalized form
$$\mu _{IJ}=\text{diag}(\mu _1,\mathrm{},\mu _N)$$
(18)
with real, positive eigenvalues $`\mu _I`$. We also choose the flavor basis for the lepton doublets $`L_i`$ such that the charged lepton Yukawa couplings are diagonal. In order to be able to couple the bulk spinors to brane field we decompose each six–dimensional field $`\mathrm{\Psi }_I`$ into two four–dimensional Dirac spinors $`\mathrm{\Psi }_{aI}`$ where $`a=+,`$ indicates the internal two–dimensional chirality of the component. The three–brane is taken to be located at $`y=0`$ in the internal space and carries, among the other standard model fields, the lepton doublets $`L_i`$, where $`i,j,\mathrm{}=e,\mu ,\tau `$, and the Higgs doublet $`H`$. These fields have a Yukawa coupling to the two components $`\mathrm{\Psi }_{aI}`$ of the bulk spinors introduced above where the dimensionless coupling constants are denoted by $`h_{aIi}`$. Note that, from the dimensionality of the bulk spinors, the corresponding operator are suppressed by one power of the string scale $`M_\mathrm{s}`$.
In summary, our model consists of $`N`$ “right–handed neutrino” bulk spinors with a Dirac mass in six dimensions and Yukawa couplings to the lepton doublets located on the three–brane. Our model, as stands, has a lepton number $`U(1)`$ symmetry with $`L_i`$ and $`\mathrm{\Psi }_I`$ each carrying one unit of charge. We would like to impose this $`U(1)`$ symmetry (or an appropriate discrete subgroup thereof) on our model. This forbids the other mass terms and couplings that we could have written in our action. In particular, it forbids the bulk Majorana mass terms $`\overline{\mathrm{\Psi }^c}_I\mathrm{\Psi }_J`$ and $`\overline{\mathrm{\Psi }^c}_I\mathrm{\Gamma }_7\mathrm{\Psi }_J`$ which would be allowed by six–dimensional Lorenz invariance. Furthermore, it forbids the brane–bulk coupling $`\overline{\mathrm{\Psi }^c}_{aI}L_iH`$ that would be allowed from four–dimensional Lorenz invariance. However, we stress that, at this stage, nothing forbids the bulk Dirac mass term. It has, therefore, been included in the above action.
Next, we would like to work out the four–dimensional effective action of our model. We compactify the additional dimensions on a two–dimensional torus that, for simplicity, is taken to be rectangular and with equal radii $`R`$ in both direction. With the Higgs vacuum expectation value $`v`$, we introduce the mass parameters
$$\widehat{m}_{aIi}=\frac{h_{aIi}v}{2\pi RM_\mathrm{s}}.$$
(19)
We expand the bulk spinors into Kaluza–Klein modes as
$$\mathrm{\Psi }(x,y)=\frac{1}{2\pi R}\underset{n^2}{}\mathrm{\Psi }_{In}(x)\mathrm{exp}\left(\frac{in_\alpha y^\alpha }{R}\right)$$
(20)
where $`n=(n_1,n_2)`$. As we did with with the spinor $`\mathrm{\Psi }_I`$ before, we decompose each of its Kaluza–Klein modes $`\mathrm{\Psi }_{In}`$ into two four–dimensional Dirac spinors $`\mathrm{\Psi }_{aIn}`$, where $`a=+,`$. Each of these four–dimensional Dirac spinors is decomposed further into Weyl spinors according to
$$\mathrm{\Psi }_{aIn}=\left(\begin{array}{c}\xi _{aIn}^c\\ \eta _{aIn}\end{array}\right).$$
(21)
The charge conjugated Weyl spinor $`\xi ^c`$ is defined as $`\xi ^c=ϵ\xi ^{}`$ where $`ϵ`$ is the two–dimensional epsilon–symbol. As a result, for each bulk spinor $`\mathrm{\Psi }_I`$ and each fixed mode number $`n`$ we obtain two four–dimensional Dirac spinors $`\mathrm{\Psi }_{aIn}`$ or, equivalently, four left–handed Weyl spinors $`\xi _{aIn}`$, $`\eta _{aIn}`$, where $`a=+,`$. With these definitions, the mass terms in the four–dimensional effective Lagrangian read
$`_{\mathrm{bulk}}`$ $`=`$ $`{\displaystyle \underset{n^2}{}}\left[m_n\xi _{+In}\eta _{In}m_n^{}\xi _{In}\eta _{+In}+\mu _I(\xi _{+In}\eta _{+In}+\xi _{In}\eta _{In})+\text{h.c.}\right]`$ (22)
$`_{\mathrm{brane}}`$ $`=`$ $`{\displaystyle \underset{n^2}{}}\left[m_{+Ii}\xi _{+In}+m_{Ii}\xi _{In}\right]\nu _i+\text{h.c.}.`$ (23)
Here, the Kaluza–Klein masses $`m_n`$ are defined by
$$m_n=\frac{in_1+n_2}{R}.$$
(24)
To interpret the structure of masses given above it is useful to diagonalize the bulk Lagrangian. Since for each mode number $`n`$ and bulk flavor $`I`$ we have two degenerate Dirac spinors, there is a degree of arbitrariness in the choice of the eigenstates. To keep the brane Lagrangian unchanged, we rotate the spinors $`\eta `$ only. We therefore introduce the linear combinations
$$\eta _{}^{}{}_{+In}{}^{}=\frac{\mu _I\eta _{+In}+m_n\eta _{In}}{\sqrt{\mu _I^2+|m_n|^2}},\eta _{}^{}{}_{In}{}^{}=\frac{\mu _I\eta _{In}m_n^{}\eta _{+In}}{\sqrt{\mu _I^2+|m_n|^2}}$$
(25)
which allows us to write the bulk Lagrangian in the form
$$_{\mathrm{bulk}}=\underset{n^2}{}\sqrt{\mu _I^2+|m_n|^2}\left[\xi _{+In}\eta _{}^{}{}_{+In}{}^{}+\xi _{In}\eta _{}^{}{}_{In}{}^{}+\text{h.c.}\right].$$
(26)
From this we read off the following mass matrix
$$=\begin{array}{c}\nu \\ \mathrm{}\\ \xi _{+n}\\ \eta _{}^{}{}_{+n}{}^{}\\ \xi _n\\ \eta _{}^{}{}_{n}{}^{}\\ \mathrm{}\end{array}\left(\begin{array}{ccccccc}0& \mathrm{}& \widehat{m}_+^T& 0& \widehat{m}_{}^T& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& & & & & \\ \widehat{m}_+& & 0& \sqrt{\mu ^2+|m_n|^2}& 0& 0& \\ 0& & \sqrt{\mu ^2+|m_n|^2}& 0& 0& 0& \\ \widehat{m}_{}& & 0& 0& 0& \sqrt{\mu ^2+|m_n|^2}& \\ 0& & 0& 0& \sqrt{\mu ^2+|m_n|^2}& 0& \\ \mathrm{}& & & & & & \mathrm{}\end{array}\right),$$
(27)
where for ease of notation we have used $`\nu (\nu _e,\nu _\mu ,\nu _\tau )^T`$, $`\xi _{an}(\xi _{a1n}\mathrm{}\xi _{aNn})^T`$, $`\eta _{an}^{}(\eta _{a1n}^{}\mathrm{}\eta _{aNn}^{})^T`$ and $`\mu =\text{diag}(\mu _1\mathrm{}\mu _N)`$. Given this matrix, it is easy to show that we have three exactly massless Weyl fermions $`\nu _{0i}`$ given by
$$\nu _{0i}=(N^1)_{ij}\left[\nu _j\underset{n,I,a}{}\frac{\widehat{m}_{aIj}^{}}{\sqrt{\mu _I^2+|m_n|^2}}\eta _{}^{}{}_{aIn}{}^{}\right].$$
(28)
Since in our context all the mass parameters are several orders of magnitude above the eV scale, it is natural to attribute the solar and atmospheric neutrino phenomenology to these massless eigenstates. The occurrence of massless eigenstates is due to the U(1) symmetry we have imposed. We will discuss below how a small breaking of this symmetry can generate the small masses necessary to account for the standard phenomenology of neutrino oscillations. As long as the states $`\nu _{0i}`$ are massless, the matrix $`N`$ in eq. (28) can be determined only up to a unitary transformation among the massless eigenstates. However, $`N`$ will be fixed (up to phase rotations) once small masses for the $`\nu _{0i}`$ have been generated. After integrating over the Kaluza–Klein modes up to the cut–off $`M_\mathrm{s}`$ one finds
$$(N^{}N)_{ij}\delta _{ij}+\pi \underset{a,I}{}\widehat{m}_{aIi}^{}\widehat{m}_{aIj}\left[\frac{1}{\mu ^2}+R^2\mathrm{ln}\frac{M_\mathrm{s}^2}{\mu ^2+1/R^2}\right].$$
(29)
We are interested in a situation where the mass matrix has a hierarchical structure satisfying $`|\widehat{m}_{aIi}|/\sqrt{\mu _I^2+|m_n|^2}1`$ for all states. In this case $`(N^{}N)_{ij}\delta _{ij}`$ and $`N`$ is approximately given by a unitary matrix $`U`$, that is
$$NU.$$
(30)
An exact diagonalization of the mass matrix $``$ shows that the massive eigenstates can be organized in two Dirac spinors for each bulk flavor index $`I`$ and each mode number $`n`$. Assuming a hierarchical structure of the mass matrix also simplifies the process of finding such eigenstates since it implies that the mass matrix (27) can be diagonalized perturbatively. While this reproduces eq. (28) for the massless states, it tells us that for each bulk flavor index $`I`$ and each mode number $`n`$ the two Dirac spinor eigenmodes $`(\stackrel{~}{\nu }_{aIn}^c,\nu _{aIn})^T`$ are given approximately by
$`\stackrel{~}{\nu }_{aIn}^c`$ $`\xi _{aIn}`$ (31a)
$`\nu _{aIn}`$ $`\eta _{aIn}^{}+{\displaystyle \underset{i}{}}{\displaystyle \frac{\widehat{m}_{aIi}}{\sqrt{\mu _I^2+|m_n|^2}}}\nu _i,`$ (31b)
where $`a=+,`$. This equation holds up to corrections of order $`|\widehat{m}_{aIi}|^2/(\mu _I^2+|m_n|^2)`$. Furthermore, in the same approximation, these two spinors have a degenerate Dirac mass
$$M_{In}=\sqrt{\mu _I^2+|m_n|^2}.$$
(32)
Given this information we can invert the formula (28) and express the weak eigenstates $`\nu _i`$ in terms of mass eigenstates $`\nu _{0i}`$, $`\nu _{aIn}`$ as
$$\nu _iU_{ij}\nu _{0j}+\underset{n,I,a}{}\frac{\widehat{m}_{aIi}^{}}{\sqrt{\mu _I^2+|m_n|^2}}\nu _{aIn}$$
(33)
up to small terms of order $`|\widehat{m}_{aIi}|^2/(\mu _I^2+|m_n|^2)`$. Eq. (33) is the starting point for studying the phenomenology of our model. Due to our convention for the flavor basis of the left–handed neutrinos (leading to diagonal charged lepton Yukawa couplings), the $`\nu _i`$ are the neutrino “flavor eigenstates”. They are expressed as a superposition of three light Majorana mass eigenstates ($`\nu _{0i}`$) and a series of left–handed components of heavy Dirac mass eigenstates ($`\nu _{aIn}`$). In the limit $`|\widehat{m}_{aIi}|^2/(\mu _I^2+|m_n|^2)1`$, the “heavy” contribution to the states $`\nu _i`$ is small, so that the flavor eigenstates are mainly light. Therefore, on one hand the model can potentially accommodate the standard oscillation phenomenology, which mainly involves the three light Majorana eigenstates $`\nu _{0i}`$ and the $`3\times 3`$ mixing matrix $`U`$. On the other hand, as we will see, the small admixture of the Kaluza–Klein states account for the KARMEN anomaly.
Let us discuss the results that we have obtained so far from the general perspective of brane–world models for neutrino masses. We have found three exactly massless Weyl fermions as given by eq. (28). In general they are superpositions of the electroweak eigenstates $`\nu _i`$ on the brane and the Kaluza–Klein states $`\eta _{aIn}^{}`$. The relative weight of these two contributions is controlled by the Dirac masses $`\mu _I`$. For vanishing Dirac masses $`\mu _I0`$ the right hand side of eq. (28) is dominated by the terms corresponding to the Kaluza–Klein zero mode $`\eta _{aI0}^{}`$. Hence, in this limit, the appearance of massless modes is not surprising. The massless states (28) are simply linear combinations of the bulk zero modes in this case. The situation becomes less trivial once we switch on $`\mu _I`$. Then, there are no massless bulk modes any more since, from eq. (32), the Kaluza–Klein spectrum is bounded from below by $`\mu _I`$. To summarize, despite the introduction of a bulk Dirac mass term, the symmetries of the higher–dimensional theory (Lorenz invariance and lepton number) lead to massless eigenstates which are non–trivial linear combinations of the brane and the bulk fields. We believe this can be of considerable relevance for brane–world models of neutrino masses, in general. A final point concerns the order of magnitude of the bulk Dirac mass $`\mu `$. As stands, the natural value of this mass is probably the string scale $`M_\mathrm{s}`$. Of course, in this case, the bulk neutrinos would be too heavy to be relevant for any neutrino physics. If the string scale is of order $`\mathrm{TeV}`$ some bulk states might receive smaller masses of order, say, $`\mathrm{MeV}`$ due to small couplings. In addition, if the bulk fermions are interpreted as superpartners of string moduli they receive masses only after supersymmetry breaking, which may also account for their smallness.
So far, we have simply assumed that the massless states acquire a small mass through a breaking of the U(1) symmetry. However, we have not explicitly shown how light neutrino masses and mixings needed for the standard neutrino phenomenology can be generated. While this is certainly appropriate for our main purpose it is, of course, important to show how small but non–vanishing neutrino masses can be incorporated into the model. This is what we will do in the remainder of the section.
In order to generate neutrino masses we need to break the lepton number symmetry that we have imposed on our model. Rather than presenting a complete model of how this can be realized spontaneously, we simply parameterize this breaking by adding the following Majorana mass terms to the bulk action (16)
$$S_{\mathrm{bulk},\mathrm{L}}=d^4xd^2y\left[\frac{1}{2}M_{IJ}^{(A)}\overline{\mathrm{\Psi }}_I^c\mathrm{\Psi }_J\frac{1}{2}M_{IJ}^{(S)}\overline{\mathrm{\Psi }}_I^c\mathrm{\Gamma }_7\mathrm{\Psi }_J+\text{h.c.}\right],$$
(34)
with $`M_{IJ}^{(A)}`$ antisymmetric and $`M_{IJ}^{(S)}`$ symmetric. They violate lepton number by two units. Furthermore, we have to amend the brane action (17) by
$$S_{\mathrm{brane},\mathrm{L}}=_{\{y=0\}}d^4x\left[\frac{l_{aIi}}{M_\mathrm{s}}\overline{\mathrm{\Psi }}_{aI}^cL_iH+\text{h.c.}\right]$$
(35)
which violates lepton number by two units as well. Recall that $`\mathrm{\Psi }_{aI}`$, where $`a=+,`$ are the two four–dimensional Dirac spinors contained in the six–dimensional bulk spinors $`\mathrm{\Psi }_I`$. As before, after electroweak symmetry breaking, we introduce the mass parameters
$$\stackrel{~}{m}_{aIi}=\frac{l_{aIi}v}{2\pi RM_\mathrm{s}}.$$
(36)
Using the expansion (20) of the bulk fermions we then find the additional four–dimensional terms
$`_{\mathrm{bulk},\mathrm{L}}`$ $`=`$ $`{\displaystyle \underset{n^2}{}}\left[M_{IJ}\eta _{+In}\eta _{Jn}M_{IJ}^{}\xi _{+In}\xi _{Jn}+\text{h.c.}\right]`$ (37)
$`_{\mathrm{brane},\mathrm{L}}`$ $`=`$ $`{\displaystyle \underset{n^2}{}}\left[\stackrel{~}{m}_{+Ii}\xi _{+In}+\stackrel{~}{m}_{Ii}\xi _{In}\right]\nu _i+\text{h.c.}`$ (38)
where the Majorana mass matrix $`M`$ is defined by $`M=M^{(S)}+M^{(A)}`$. Now our total four–dimensional effective action consists of the original lepton number preserving parts (22), (23) and the above lepton number violating parts (37), (38). It would clearly be interesting to thoroughly investigate the neutrino phenomenology resulting from this action that includes all obvious bulk mass terms allowed by higher–dimensional Lorenz invariance. We hope to return to this problem in a future publication. For the purpose of the present paper, we would merely like to check the order of magnitude of masses that are induced. To do this, we focus on the case of one flavor in the bulk as well as on the brane. Assuming that $`\widehat{m}_a/\mu 1`$ and $`\stackrel{~}{m}_a/\mu 1`$ we can apply the see–saw formula to find the light neutrino masses. After summing over all Kaluza–Klein states with cut–off $`M_\mathrm{s}`$ we find for the mass $`m_0`$ of the previously massless states $`\nu _{0i}`$
$`m_0`$ $`=`$ $`(\stackrel{~}{m}_{}\widehat{m}_{}\mu \stackrel{~}{m}_{}\stackrel{~}{m}_+M+\widehat{m}_{}\widehat{m}_+M^{}\stackrel{~}{m}_+\widehat{m}_+\mu )`$ (39)
$`\times \left({\displaystyle \frac{2}{\mu ^2+|M|^2}}+2\pi R^2\mathrm{ln}{\displaystyle \frac{M_\mathrm{s}^2}{\mu ^2+|M|^2+1/R^2}}\right).`$
As expected, this vanishes if we restore lepton number by setting $`\stackrel{~}{m}_a=M=0`$. The term in the first parenthesis together with the first term in the second parenthesis correspond to what one normally expects for a see–saw suppressed neutrino mass. In fact, these terms describe the contribution from the lightest Kaluza–Klein state with mode number $`n=0`$. However, in addition we have a large number of higher Kaluza–Klein states each of which contributes to the light neutrino mass. The net effect of all these states is encoded in the last term in eq. (39). This term is proportional to $`R^2`$ rather than the explicit mass scales $`\mu `$ or $`M`$. Hence it corresponds to a suppression by the mass scale $`1/R`$ associated to the size of the additional dimensions. As we will see, in our context, this suppression is sufficient to obtain reasonable small neutrino masses under plausible assumptions. Eq. (39) also illustrates the above mentioned decoupling of neutrino masses and neutrino–bulk mixing. The former depend in particular on the lepton number violating quantities $`\stackrel{~}{m}_\pm `$ and $`M`$. As long as these quantities are sufficiently small the light neutrinos are still approximately given by eq. (28) and, hence, they depend on lepton number conserving parameters only.
## 5 Quantitative analysis of the model
We would now like to demonstrate more quantitatively that the six–dimensional model presented in the previous section can explain the KARMEN anomaly and is, at the same time, compatible with other phenomenological constraints.
Let us start by explaining the requirements on our model that follow from the KARMEN anomaly. This amounts to specifying the details of the fit that we have presented in Section 2 and work out some of its consequences. Let us start with the constraint (2) on the branching ratio and on the decay width of the $`X`$ particle that our model has to satisfy. We should express this constraint in terms of the parameters of our model. The $`X`$ particle corresponds to a mode $`n_X`$ such that
$$m_X=\sqrt{\mu ^2+|m_{n_\mathrm{X}}|^2}.$$
(40)
As we saw in the previous Section, two almost degenerate Dirac mass eigenstates with opposite internal chirality are associated to this mode $`n_X`$, namely $`(\stackrel{~}{\nu }_{+n_X}^c,\nu _{+n_X})^T`$ and $`(\stackrel{~}{\nu }_{n_X}^c,\nu _{n_X})^T`$. The mixing $`U_{iX}^a`$, where $`a=+,`$ between the SM flavor eigenstates $`\nu _i`$, $`i=e,\mu ,\tau `$, and these two states can be read off from eq. (33) as
$$|U_{iX}^a|=\frac{|\widehat{m}_{ai}|}{\sqrt{\mu ^2+|m_{n_X}|^2}}=\frac{|\widehat{m}_{ai}|}{m_X}.$$
(41)
However, the two states are not completely degenerate but split at higher order in the see–saw approximation. For some of the parameter space this splitting will be smaller than the $`X`$ particle mass gap $`\delta m_X6\mathrm{keV}`$. In this case, both states contribute to the KARMEN anomaly and can effectively be accounted for by a single state $`X`$ with mass parameter
$$\widehat{m}_i^2=|\widehat{m}_{+i}|^2+|\widehat{m}_i|^2.$$
(42)
If, on the other hand, the splitting is larger than $`\delta m_X6\mathrm{keV}`$ one mass eigenstate will be above the threshold for production in the KARMEN experiment and the anomaly is due to the other state. We will call the corresponding mass parameter $`\widehat{m}`$ as well, but it is now a function of $`\widehat{m}_{\pm i}`$ generally different from the one given above. For both cases, we define the mixing matrix elements by
$$|U_{iX}|=\frac{\widehat{m}_i}{m_X}.$$
(43)
Through those mixings, the mass parameters $`\widehat{m}_i`$ control the $`X`$ particle production and decay and therefore the number of anomalous events in the KARMEN experiment. The branching ratio for the decay $`\pi ^+\mu ^+X`$ is given by
$$\text{BR}(\pi ^+\mu ^+X)=|U_{\mu X}|^2\frac{\lambda (m_\pi ^2,m_\mu ^2,m_X^2)}{\lambda (m_\pi ^2,m_\mu ^2,0)}0.028|U_{\mu X}|^2,$$
(44)
where
$$\lambda (a,b,c)\left[a(b+c)(bc)^2\right]\sqrt{a^2+b^2+c^22ab2ac2bc}.$$
(45)
Possible contributions from bulk dynamics to the visible decay width are highly suppressed compared to those from standard $`W`$ and $`Z`$ exchange. The impact of bulk dynamics on the total width can in principle be non–negligible but it would not change drastically the lifetime, which means that it would not affect our conclusions at all. Neglecting the sub–dominant $`X\gamma \nu `$ decay, the visible decay width is given by
$$\begin{array}{c}\mathrm{\Gamma }_{\text{vis}}\mathrm{\Gamma }(Xe^+e^{}\nu )\hfill \\ \hfill =\left(\frac{m_X}{m_\mu }\right)^5\left[|U_{eX}|^2\left(\frac{1}{4}+s_W^2+2s_W^4\right)+(|U_{\mu X}|^2+|U_{\tau X}|^2)\left(\frac{1}{4}s_W^2+2s_W^4\right)\right]\tau _\mu ^1\\ \hfill 890\text{sec}^1|U_{eX}|^2+195\text{sec}^1(|U_{\mu X}|^2+|U_{\tau X}|^2).\end{array}$$
(46)
Putting these results together, the constraint (2) can be expressed in terms of the mixings $`|U_{iX}|`$ as follows
$$\begin{array}{c}(1\text{}\mathrm{\hspace{0.17em}2})\mathrm{\hspace{0.17em}10}^{11}=\text{BR}(\pi ^+\mu ^+X)\mathrm{\Gamma }_{\text{vis}}\text{sec}\hfill \\ \hfill 25|U_{\mu X}|^2|U_{eX}|^2+5.5|U_{\mu X}|^2(|U_{\mu X}|^2+|U_{\tau X}|^2).\end{array}$$
(47)
If we make the plausible assumption $`\widehat{m}_e^2\widehat{m}_\mu ^2,\widehat{m}_\tau ^2`$, it follows that the first term on the right hand side of the previous equation is sub–dominant. Eq. (47) then represents a constraint on $`|U_{\mu X}|`$ and $`|U_{\tau X}|`$, or, equivalently on $`\widehat{m}_\mu `$ and $`\widehat{m}_\tau `$ which reads
$$\widehat{m}_\mu ^2(\widehat{m}_\mu ^2+\widehat{m}_\tau ^2)=\left[(40\text{}\mathrm{\hspace{0.17em}47})\mathrm{keV}\right]^4.$$
(48)
This constraint relates $`\widehat{m}_\mu `$ and $`\widehat{m}_\tau `$ and allows us to compute $`\text{BR}(\pi ^+\mu ^+X)`$ and $`\mathrm{\Gamma }_{\text{vis}}`$ as a function of the ratio
$$r\widehat{m}_\mu /\widehat{m}_\tau .$$
(49)
The result is shown in Fig. 2 for the range $`10^4<\widehat{m}_\mu /\widehat{m}_\tau <1`$ which corresponds to the plausible situation $`\widehat{m}_\mu <\widehat{m}_\tau `$. The width of the band in these figures corresponds to the uncertainty on the right hand side of eq. (48). If the “atmospheric” mixing angle measured by SuperKamiokande originated in the diagonalization of the charged lepton mass matrix, we would expect $`\widehat{m}_\mu /\widehat{m}_\tau =𝒪\left(1\right)`$. This is because the brane–bulk Yukawa coupling was written in a basis where the charged lepton mass matrix is diagonal. However, $`\widehat{m}_\mu /\widehat{m}_\tau =𝒪\left(1\right)`$ gives a too high branching ratio $`\text{BR}(\pi ^+\mu ^+X)`$ according to the present PSI upper limit (solid horizontal line in Fig. 2a). The new expected limit is also shown (dashed horizontal line in Fig. 2a). The latter would imply $`\widehat{m}_\mu /\widehat{m}_\tau 10^2`$. Therefore, we focus on the case where $`\widehat{m}_\mu /\widehat{m}_\tau 1`$.
Furthermore, a lower constraint on the $`\widehat{m}_\mu /\widehat{m}_\tau `$ range is given by the requirement that $`\widehat{m}_\tau `$ is not too large. If $`\widehat{m}_\tau `$ were of the order of the Dirac mass $`\mu `$ or more, the $`\tau `$ neutrino would predominantly consist of a heavy mass eigenstate, which would be phenomenologically problematic. Consequently, $`r`$ should be preferably in the range
$$r=10^4\text{to}\mathrm{\hspace{0.33em}10}^2.$$
(50)
Then, the mass parameters $`\widehat{m}_i`$ of our model as well as the branching ratio and the decay width of the $`X`$ particle can be expressed in terms of $`r`$ as
$$\widehat{m}_\mu =(4.0\text{}\mathrm{\hspace{0.17em}4.7})\mathrm{keV}\left(\frac{r}{10^2}\right)^{1/2}$$
(51a)
$$\widehat{m}_\tau =(0.4\text{}\mathrm{\hspace{0.17em}0.47})\mathrm{MeV}\left(\frac{10^2}{r}\right)^{1/2}$$
(51b)
$$\text{BR}(\pi ^+\mu ^+X)=(3.8\text{}\mathrm{\hspace{0.17em}5.3})\mathrm{\hspace{0.17em}10}^{10}\left(\frac{r}{10^2}\right)$$
(51c)
$$\tau _X<\mathrm{\Gamma }_{\text{vis}}^1=(27\text{}\mathrm{\hspace{0.17em}38})\text{sec}\left(\frac{r}{10^2}\right).$$
(51d)
It remains to fix the radius $`R`$ of the two additional dimensions. Of course, $`1/R`$ should be smaller than the mass $`m_\mathrm{X}33.9\mathrm{MeV}`$ of the $`X`$ particle. Furthermore, as we have seen in Section 2, a value of $`1/R1.6\mathrm{MeV}`$ leads to a complete solution of the fine–tuning problem. However, as explained there, we need to choose a somewhat larger value $`1/R15\mathrm{MeV}`$ to avoid significant contributions from modes lighter than the $`X`$ particle to the KARMEN spectrum. In summary, we therefore require
$$15\mathrm{MeV}1/R<33.9\mathrm{MeV}.$$
(52)
The fine–tuning is smallest at the lower end of this range which is why we have chosen $`1/R15\mathrm{MeV}`$ in the fit performed in Section 2. Finally, we have to fix the order of magnitude of our Dirac mass $`\mu `$. To avoid too large admixture of the lightest Kaluza–Klein mode we have to require that $`\mu \widehat{m}_\tau `$. In addition, of course, $`\mu `$ has to be smaller than the $`X`$ particle mass. Together, this leads to
$$\widehat{m}_\tau \mu <33.9\mathrm{MeV}.$$
(53)
We remark that this allows for the plausible value $`\mu =O(1/R)`$ for the Dirac mass. Finally, we should check the assumption made in the previous section about the unitarity of the matrix $`N`$ and the smallness of the heavy states contribution to the SM neutrinos. Both issues are related to the square of the normalization matrix $`N^{}N`$ defined in eq. (29) which should be close to unity. For our favorite values $`1/R\mu 15\mathrm{MeV}`$ and $`M_\mathrm{s}=1\mathrm{TeV}`$ this square is given by
$$N^{}N1+0.05\left(\frac{10^2}{r}\right)$$
(54)
It constraints $`r`$ to be in a subset of the range (50) which, however, is still comfortably large.
Now that we have basically fixed all the lepton number conserving quantities in our model we should analyze what we need to do to generate small masses for the neutrinos. In the previous section we have shown that we can generate neutrino masses by introducing a lepton number breaking bulk Majorana mass $`M`$ and lepton number violating brane–bulk masses $`\stackrel{~}{m}_\pm `$. Both mass terms violate lepton number by two units. For the simple case of one neutrino flavor the resulting neutrino mass has been given in eq. (39). Using our favorite values $`1/R\mu 15\mathrm{MeV}`$ and $`\widehat{m}_\pm \widehat{m}_\tau `$ where $`\widehat{m}_\tau `$ has been specified in eq. (50) we find
$$m_0\left(\frac{10^2}{r}\right)\left(2\frac{\stackrel{~}{m}_\tau }{\widehat{m}_\tau }+\frac{M}{\mu }\right)(1.5\text{}\mathrm{\hspace{0.17em}2.0})\mathrm{MeV}.$$
(55)
Note that the two ratios $`\stackrel{~}{m}_\tau /\widehat{m}_\tau `$ and $`M/\mu `$ of lepton number violating and conserving quantities on the brane and in the bulk, respectively, enter with the same strength. This fits nicely to the fact that they both break lepton number by two units. More quantitatively, to have a neutrino mass $`m_01\mathrm{eV}`$ we need both ratios to be less than roughly $`10^4r`$. Such a number may, for example, arise from a suppression of the form $`<\varphi >/M_\mathrm{s}`$ where $`\varphi `$ is a boson that carries lepton number and takes a VEV of the order, say $`1/R`$.
Although the main focus of the paper is on the particle physics aspects of our model, we would like to briefly address constraints from big–bang nucleosynthesis and astrophysics. In the early universe, the Kaluza–Klein mode with mode number $`n`$ and mixing $`U_{in}`$ decouples at a temperature of
$$T_{\mathrm{dec},\mathrm{n}}1\mathrm{MeV}|U_{in}|^{2/3}.$$
(56)
From eq. (50) the largest mixing angles are those for $`i=\tau `$. Using these angles and $`\mu =1/R15\mathrm{MeV}`$ as above we find
$$T_{\mathrm{dec},\mathrm{n}}(10\text{}11)\mathrm{MeV}\left(\frac{r}{10^2}\right)^{1/3}(1+n^2)^{1/3}.$$
(57)
It follows than the decoupling temperature is lower than the mass for all Kaluza–Klein modes. In particular, the decoupling temperature increases with a smaller power of the mode number $`n`$ than the masses $`M_n15\mathrm{MeV}\sqrt{1+n^2}`$ of the Kaluza–Klein modes. As a consequence, modes with large $`n`$ decouple when they are highly non–relativistic and are strongly diluted. Unfortunately, this suppression is not strong enough for low $`n`$ modes. Therefore, we have to assume that the bulk is empty at some temperature $`T<T_{\mathrm{dec},0}10\mathrm{MeV}`$ as it is customary for models with large additional dimensions . Taking this temperature sufficiently low (but above $`1\mathrm{MeV}`$, of course) recreation of Kaluza–Klein particles is Boltzmann–suppressed. As a consequence, their thermal production rate can always be kept below the Hubble rate.
Supernova cooling provides another constraint on our model. For the single–particle explanation of the KARMEN anomaly the supernova energy loss induced by the $`X`$ particle has been estimated in ref. . From this estimate, it is concluded there, that the single particle explanation of the KARMEN anomaly is ruled out. Our model is even more problematic since we have a tower of particles rather than just a single one. However, a closer analysis of the supernova energy loss is likely to weaken the bound given in ref. . This together with the assumption of a moderately lower supernova temperature might render our model consistent with the supernova constraint. In any case, we believe that our model is of interest from the viewpoint of particle physics and should be compared with the direct experimental information on neutrino properties.
## 6 Conclusions
In this paper, we have shown that the KARMEN anomaly can be understood in the context of a six–dimensional brane–world model. The slow particle $`X`$, responsible for the anomaly, is identified with a specific Kaluza–Klein excitation of a bulk fermion that, from a four–dimensional point of view, appears as a sterile neutrino. We have pointed out that in any interpretation of the KARMEN anomaly based on a single slow particle produced in the $`\pi ^+\mu ^+X`$ decay, the $`X`$ mass is fine–tuned to the mass difference $`m_\pi m_\mu `$ with an accuracy of order $`10^4`$. This means that $`\delta m_X/(m_\pi m_\mu )\mathrm{1.8\hspace{0.17em}10}^4`$, where $`\delta m_X=(m_\pi m_\mu )m_X`$. Such a problem can be significantly alleviated, although only partially solved, in the approach we propose, since the $`X`$ particle is just one of the many excitation of the bulk fermion. It is therefore relatively more probable that one out of the many states has a mass that falls into the critical range.
The phenomenology of the model with respect to the KARMEN experiment depends on the average separation $`\mathrm{\Delta }m_{\text{KK}}`$ between two Kaluza–Klein states at the scale $`m_X`$. In the limit where $`\mathrm{\Delta }m_{\text{KK}}`$ is large, that is $`\mathrm{\Delta }m_{\text{KK}}=𝒪\left(m_\pi m_\mu \right)`$, the model is indistinguishable from models with a single sterile neutrino. In this limit, although the fine–tuning is not improved, we still have a model that provides a theoretically well–motivated origin for the $`X`$ particle. Moreover, as a quite non–trivial feature, the model allows to incorporate the large mass scale of the $`X`$ particle as well as the small neutrino masses needed for the standard oscillation phenomenology. The limit of small $`\mathrm{\Delta }m_{\text{KK}}`$, that is $`\mathrm{\Delta }m_{\text{KK}}=𝒪\left(\delta m_X\right)`$, is ruled out because it would lead to a modification of the KARMEN time spectrum in the region $`(0.6\text{}\mathrm{\hspace{0.17em}3.1})\mu \mathrm{s}`$. For intermediate values of $`\mathrm{\Delta }m_{\text{KK}}`$, $`\mathrm{\Delta }m_{\text{KK}}=𝒪\left(1\right)\mathrm{MeV}`$ or more, the fine–tuning problem is significantly alleviated and the particles lighter than $`X`$ could give a small contribution to the early KARMEN time spectrum. Therefore, a detailed study of the time distribution of events at early times after the end of the proton pulse could test the model if the mass separation $`\mathrm{\Delta }m_{\text{KK}}`$ happens not to be too large. An upgraded detector with tracking capability devoted to the study of the anomaly could certainly explore this possibility. Needless to say, the detection of a signal at early times would represent a strong hint for the existence of a brane–world.
On the model building side, we have identified two possible structures for the sizes of the additional dimensions. The first case represents an effectively five–dimensional model with intermediate string scale and the remaining five dimensions being close to the string scale. The second one has a $`\mathrm{TeV}`$ string scale, two large dimensions in the $`\mathrm{MeV}^1`$ range and the scale of the remaining four dimensions being two to three orders of magnitude larger. Therefore, in this case the space–time is effectively six–dimensional in an intermediate energy range. We have described in detail a six–dimensional example. As in all brane–world neutrino models, we have considered bulk neutrinos whose Lorenz invariant Dirac and Majorana masses are suppressed relative to the string scale. In particular, we have taken the magnitude of the Dirac mass term to be of the order of the inverse size of the additional dimensions and we have further suppressed the Majorana mass term by means of a U(1) symmetry (lepton number). Despite the introduction of the bulk Dirac mass term, we have found that the model accommodates three light Majorana neutrinos which would be massless in the limit of unbroken U(1) symmetry. The light states are predominantly given by the electroweak eigenstates with small but sizeable admixture of Kaluza–Klein modes. We have therefore attributed the solar and atmospheric neutrino phenomenology to those light states. On the other hand the small heavy component of the flavor eigenstates accounts for the KARMEN anomaly. Therefore, the introduction of the Dirac mass term makes our model rather different from the models considered in the literature so far, where the massless modes either do not exist or decouple from the left–handed neutrinos. Moreover, the neutrino masses and the mixing of neutrinos and bulk states are controlled by different parameters. We believe that such a scheme can be of general interest for brane–world models of neutrino masses.
## 7 Acknowledgments
We thank Pierre Ramond, Graham Ross and Subir Sarkar for useful discussions and suggestions. This work is supported by the TMR Network under the EEC Contract No. ERBFMRX–CT960090.
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# O VI and Multicomponent H I Absorption Associated with a Galaxy Group in the Direction of PG0953+415: Physical Conditions and Baryonic Content1footnote 11footnote 1Based on observations obtained with the WIYN Observatory, which is a joint facility of the University of Wisconsin, Indiana University, Yale University, and the National Optical Astronomy Observatories.,2footnote 22footnote 2Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-2555.
## 1 Introduction
Hydrodynamic simulations of cosmological structure growth predict that when the initial density perturbations collapse, gas should be shock-heated to temperatures of $`10^5`$$`10^7^{}`$K (Ostriker & Cen 1996; Davé et al. 1999; Cen & Ostriker 1999a). The fraction of the gas which has been heated to these temperatures increases with decreasing redshift, and at the present epoch, the model of Cen & Ostriker (1999a) predicts that 47 % of the baryons (by mass) are in this shock-heated phase, hereafter referred to as warm/hot gas (to distinguish it from the hotter gas in rich galaxy clusters which are readily detected X-ray sources). This warm/hot gas prediction has not been adequately tested by observations because the soft X-rays emitted by gas at these temperatures are difficult to detect with current instrumentation, especially at lower temperatures where corrections for foreground absorption and emission are complicated. However, it may be possible to detect gas in the lower half of this temperature range via absorption lines of species such as O VI, Ne VIII, or Mg X in the spectrum of a background QSO (Verner, Tytler, & Barthel 1994). It is important to search for this warm/hot gas as part of the census of matter in the universe and because it could affect the formation and evolution of galaxies and galaxy groups and clusters (e.g., Blanton et al. 2000).
There are some indications that warm/hot gas is present in some galaxy groups. For example, Mulchaey et al. (1996) have noted that ROSAT observations show that poor galaxy groups which are rich in elliptical galaxies tend to exhibit X-ray emission (E $`>`$ 0.5 keV) while spiral-rich groups do not. They suggest that spiral-rich groups contain cooler ($`4\times 10^6`$ K) intragroup gas which is not easily detected in X-rays, and they predict that the intragroup medium of spiral-rich groups will produce absorption lines of O VI, but not C IV or N V because their column densities are too low at $`T>5\times 10^5^{}`$K. Interestingly, Savage, Tripp, & Lu (1998) have recently identified a QSO absorber in the spectrum of H1821+643 ($`z_{\mathrm{QSO}}`$ = 0.297) which fits this description: O VI absorption lines with two nearby spiral galaxies but no accompanying C IV or N V lines. However, the absorption could be due to the halo of the closer spiral galaxy (which is at a projected distance of 105 $`h_{75}^1`$ kpc) rather than the intragroup medium, and Savage et al. show that the O VI absorption could plausibly arise in very low density photoionized gas. The O VI doublet has been identified in several other intervening absorption systems<sup>3</sup><sup>3</sup>3Strong O VI absorption lines have also been detected in “associated” absorption line systems with $`z_{\mathrm{abs}}`$ $`z_{\mathrm{QSO}}`$ (e.g., Papovich et al. 2000 and references therein). These are a rather different class of absorber which are often known to be very close to the QSO (Hamann & Ferland 1999). In this paper, we have focused our analysis and discussion on the intervening systems. both at moderate redshifts (e.g., Bergeron et al. 1994; Burles & Tytler 1996; Jannuzi et al. 1998; Lopez et al. 1999; Churchill & Charlton 1999) and at high redshifts (Kirkman & Tytler 1997,1999). Composite spectra and statistical techniques have also been used to show that O VI absorption is present at high $`z`$ (Lu & Savage 1993; Davé et al. 1998). In most cases it has been difficult to pin down the ionization mechanism definitively, partly due to the low resolution of the observations made with first-generation Hubble Space Telescope (HST) spectrographs, but in several cases there is evidence that the O VI systems occur in multiphase absorbing media. An important step in this approach to the search for warm/hot gas is to determine whether the O VI absorbers trace collisionally ionized gas or photoionized gas.
As part of the program described by Tripp, Lu, & Savage (1998) to study low $`z`$ Ly$`\alpha `$ absorption line systems, we have recently observed the low redshift QSO PG0953+415 ($`z_{\mathrm{QSO}}`$ = 0.239) with the E140M echelle mode of the Space Telescope Imaging Spectrograph (STIS). This high resolution FUV spectrum has revealed another highly ionized O VI absorber associated with a group of spiral galaxies, and in this paper we present our analysis of this particular absorbing system. In §2 we review the observations and data reductions including measurements of galaxy redshifts with the WIYN telescope. We present in §3 the absorption line measurements. We constrain the temperature of the gas and examine its ionization in §4, and we discuss the implications of the observations in §5. Throughout this paper we assume $`H_0=75h_{75}`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_0`$ = 0.0. Also, all wavelengths and redshifts reported here are heliocentric, but in this direction heliocentric and LSR velocities are nearly identical.<sup>4</sup><sup>4</sup>4Assuming the Sun is moving in the direction $`l=53^{},b=25^{}`$ at 16.5 km s<sup>-1</sup> (Mihalas & Binney 1981), $`v_{\mathrm{LSR}}=v_{\mathrm{helio}}0.1`$ km s<sup>-1</sup>.
## 2 Observations
### 2.1 STIS Spectroscopy of PG 0953+415
PG 0953+415 was observed with STIS on 1998 December 4 and 1998 December 11 resulting in a total integration time of 24478 seconds.<sup>5</sup><sup>5</sup>5The QSO was observed for 10743 seconds on December 4 (HST archive ID no. O4X002010) and for 13735 seconds on December 11 (HST archive ID nos. O4X001010 and O4X001020). The observations were obtained with the medium resolution FUV echelle mode (E140M) with the 0.2$`\times `$0.2” aperture. Kimble et al. (1998) report that this STIS mode provides a resolution of $`R=\lambda /\mathrm{\Delta }\lambda `$ = 46000 (FWHM $``$ 7 km s<sup>-1</sup>) with the 0.2$`\times `$0.06” slit. With the 0.2$`\times `$0.2” slit, the FWHM is nearly identical but the broad wings of the line spread function are stronger (see Figure 13.87 in the Cycle 9 STIS Instrument Handbook). Photoevents were recorded by the FUV-MAMA detector in accumulation mode with on-board orbital Doppler compensation, and individual exposures ranged from 2215 to 2880 seconds in duration. Short exposures of a wavelength calibration lamp were obtained with each individual observation of the QSO.
The data were reduced with the software developed by the STIS Instrument Definition Team (IDT). Pixel-to-pixel sensitivity variations in the individual exposures were first corrected with a postlaunch flatfield, and then the individual images were aligned by cross-correlation<sup>6</sup><sup>6</sup>6The images can be aligned by cross-correlating the target images or the corresponding wavelength calibration lamp images. In principle direct cross-correlation of the target images is preferred because drift of the target in the aperture could cause shifts derived from the comparison lamp exposures to differ from the shifts appropriate for the target exposures. However, in this case the target images must contain enough sharp features to produce a sharp peak in the cross-correlation. We tried both methods, and we found that for this QSO observation, slightly better resolution was always obtained by cross-correlating the wavelength lamp exposures to determine the image shifts. and then coadded. To prolong the lifetime of the MAMA detector, the position of the dispersed spectrum on the detector is changed every few weeks. This was done between the first observation of PG 0953+415 and the second, which is beneficial because when the data are aligned and coadded, any residual fixed-pattern noise not adequately removed by the flatfield is reduced. This also enables the identification of spurious features by comparison of the wavelength-calibrated spectra obtained on the two occasions. After coaddition of the individual images, the scattered light correction described by Bowers and Lindler (1999) was applied. Then the spectra from each order were extracted with a simple unweighted slit with the standard height of 11 “lores” pixels,<sup>7</sup><sup>7</sup>7The MAMAs are designed to support half-pixel centroiding (see Timothy 1994), and these half-pixels are usually referred to as “hires” pixels. This capability was not employed here since it does not enhance the data for this particular observation. and smoothing of the background region also followed standard procedures for the STIS echelle modes. During the extraction, obviously hot pixels were fixed by interpolation between the adjacent pixels in the dispersion direction. The wavelength scale was computed with a postlaunch dispersion relation (which is updated when the position of the spectrum on the detector is changed) and put on a heliocentric basis. Finally, the data were flux calibrated and overlapping regions of adjacent orders were coadded with weighting based on signal-to-noise (S/N). The wavelength coverage is $``$1150-1730 Å with 5 small gaps between orders at $`\lambda >`$ 1630 Å.
Samples of the final spectrum are shown in Figure 1. We identify the lines at 1178.8 and 1185.3 Å as the O VI $`\lambda \lambda `$1032,1038 doublet at $`z_{\mathrm{abs}}`$ = 0.14232, and the corresponding H I Ly$`\alpha `$ line is well-detected at 1388.7 Å. A marginal absorption line is apparent at the expected wavelength of Ly$`\beta `$ as well, but it is recorded in a region where the spectrum is quite noisy. Note that the Ly$`\alpha `$ profile shows several components, the strongest of which is well-aligned with the O VI lines. There are no convincing alternative identifications of these lines. They are unlikely to be due to the ISM of the Milky Way since there are no resonance absorption lines at these wavelengths (see Morton 1991). The lines at 1178.8 and 1185.3 Å cannot be Ly$`\beta ,\gamma ,`$ or $`\delta `$ lines due to extragalactic absorption systems at different redshifts because the corresponding Ly$`\alpha `$ lines are not detected at the expected wavelengths. Similarly, we have searched for alternative metal line identifications at other redshifts, and we find no convincing candidates.
Since the O VI doublet was recorded in a spectral region which is relatively noisy, and since there are slightly hot spurious pixels apparent in this portion of the spectrum (see Figure 1), we inspected the spectra obtained on December 4 and December 11 independently to confirm that the O VI lines are real. Both O VI lines are detected in the individual December 4 and December 11 spectra, and the relative line depths and wavelength difference of the lines are consistent with the O VI identification. However, we noticed that the O VI $`\lambda 1037.6`$ line is broader in the December 4 spectrum than the December 11 spectrum. For the reader’s inspection, Figure 2 compares the O VI $`\lambda `$1037.6 profiles derived from the 1998 December 4 and 1998 December 11 data and from all of the data combined. We suspect that the 1037.6 Å line is broader in the December 4 spectrum due to noise, perhaps due to fixed pattern noise which was not adequately removed by the flatfield available when the data were reduced. Of course, this impacts the final spectrum as well, and the 1037.6 Å line is broader than the 1031.9 Å line in the final spectrum shown in Figure 1, though the magnitude of the discrepancy is reduced. It is alternatively possible that the O VI $`\lambda `$1031.9 line is artificially narrow due to noise or an unrecognized warm pixel which fills in the profile somewhat – we cannot rule out or favor any of these possibilities with the current STIS data. Consequently, we provide below results from measurement of the 1031.9 Å line only as well as results from measurement of both O VI lines.
### 2.2 WIYN Galaxy Redshift Measurements
As part of our continuing study of the relationship between low $`z`$ Ly$`\alpha `$ absorbers and galaxies (see Tripp, Lu, & Savage 1998), we have also measured the redshifts of galaxies in the $`1^{}`$ field centered on PG 0953+415 with the fiber-fed multiobject spectrograph (Hydra, see Barden & Armandroff 1995) on the WIYN telescope. The observations and measurement techniques are described in Tripp et al. (1998), and the new redshifts will be provided in a subsequent paper. Here we note that this survey has revealed four galaxies within $``$130 km s<sup>-1</sup> of the O VI absorber at $`z_{\mathrm{abs}}`$ = 0.14232. The measured redshifts of these galaxies are summarized in Table 1 along with some of their properties including projected distance from the sight line ($`\rho `$), velocity displacement from the absorber, and $`O`$ and $`E`$ magnitudes from the revised APS digitization of the POSS I plates (Pennington et al. 1993). We also list $`B_J`$ magnitudes calculated using the relation between $`O_{\mathrm{APS}}`$ and $`B_J`$ derived by Odewahn & Aldering (1995) and the corresponding absolute magnitudes in Table 1. These are all luminous galaxies ranging from 0.6$`L`$ to 4.0$`L`$ (assuming $`M_B`$ = –19.5 from Loveday et al. 1992). The velocity dispersion of this group is very uncertain since we do not know the redshift of the group center of mass and because of the small sample. However, if we assume that the O VI absorption arises at the center of mass of the group, then the radial velocity dispersion is roughly 100 km s<sup>-1</sup>, consistent with typical values observed in poor groups (e.g., Zabludoff & Mulchaey 1998).
The spectrum of the galaxy in Table 1 with the smallest projected separation from the sight line ($`\rho `$ = 395 kpc) shows strong Balmer lines in absorption as well as the other usual absorption lines, e.g., Ca II H & K, and well detected \[O II\] and H$`\beta `$ emission lines with observed equivalent widths of 10.5$`\pm `$2 Å and 4.5$`\pm `$0.5 Å, respectively. However, the \[O III\] emission lines are not apparent (H$`\alpha `$ is redshifted beyond the red end of the spectrum). Based on its luminosity and the Kennicutt (1992a,b) atlas of integrated galaxy spectra, this implies that the galaxy is a normal Sb-Sc spiral or a peculiar S0-Sa galaxy, without an active nucleus. The galaxy at $`\rho `$ = 3.0 Mpc is also an emission line galaxy with \[O II\], \[O III\], H$`\beta `$, and even H$`\gamma `$ seen in emission. This object has strong emission lines; the observed \[O II\] equivalent width is 37$`\pm `$10 Å. There are no indications of an active nucleus (e.g., broad H$`\beta `$ or large \[O III\]/H$`\beta `$ ratios), so given its high luminosity we classify this galaxy as a starbursting spiral or a normal Sc-Sd spiral. The other galaxies in Table 1 do not have emission lines and cannot be unambiguously classified.
It is important to note that the galaxy redshift survey we have been able to carry out to date is not very complete except at relatively bright limiting magnitudes, and it is quite possible that there are additional galaxies at $`z`$ 0.142 which are closer to the sight line than those in Table 1. Consequently, we cannot comment on whether or not a dwarf galaxy could give rise to the O VI system, for example. As we shall discuss in §5.1, it is difficult to discriminate between the hypothesis that the O VI absorption originates in an intragroup medium and the hypothesis that it occurs in the ISM of an individual galaxy given the limited information currently available. More galaxy redshift measurements and deep imaging would be valuable.
## 3 Absorption Line Measurements
The absorption line profiles of the O VI and H I Ly$`\alpha `$ lines at $`z_{\mathrm{abs}}`$ = 0.14232 are plotted versus restframe velocity in Figure 3. This figure also shows the spectral regions of lines of interest which are not detected. Restframe equivalent widths ($`W_\mathrm{r}`$) of the O VI and H I absorption lines were measured using the software of Sembach & Savage (1992), but adapted to use the statistical uncertainties directly calculated from the signal and background counts recorded by STIS. This software accounts for errors due to uncertainty in the height and curvature of the continuum as well as a 2% uncertainty in the flux zero point in the overall uncertainty in $`W_\mathrm{r}`$ \[see the appendix in Sembach & Savage (1992) for details\]. These equivalent widths are listed in Table 2. This software was also used to set upper limits on the N V, Si III, Si II, and C II equivalent widths, which are also listed in Table 2.
To measure column densities, we have used two methods: the apparent column density technique (e.g., Savage & Sembach 1991) and Voigt profile fitting. For Voigt profile fitting, we have used the software of Fitzpatrick & Spitzer (1997) with the line spread functions for the E140M mode with the 0.2$`\times `$0.2” aperture shown in Figure 13.87 of the Cycle 9 STIS Handbook. This software provides Doppler parameters ($`b`$), velocities ($`v`$), and column densities ($`N`$) for a number of components which is subjectively specified by the user. The profile fitting results are summarized in Table 3. Since the O VI $`\lambda `$1037.6 Å line may be artificially broadened by a noise feature (see §2.1), we present two alternative fits in Table 3: a fit to the O VI $`\lambda `$1031.9 line only and a fit to both of the O VI lines.
In the apparent column density approach, the apparent column density per unit velocity, $`N_\mathrm{a}(v)`$, is calculated directly from the apparent optical depth,
$$N_\mathrm{a}(v)=(m_\mathrm{e}c/\pi e^2)(f\lambda )^1\tau _\mathrm{a}(v)=3.768\times 10^{14}(f\lambda )^1\mathrm{ln}[I_\mathrm{c}(v)/I(v)]$$
(1)
in atoms cm<sup>-2</sup> (km s<sup>-1</sup>)<sup>-1</sup>, where $`f`$ is the oscillator strength, $`\lambda `$ is the wavelength in Å, $`I(v)`$ is the observed line intensity and $`I_\mathrm{c}(v)`$ is the estimated continuum intensity at velocity $`v`$, and the other symbols have their usual meanings. If the lines are fully resolved or are not affected by saturation, then the total column density can be obtained by simple integration, $`N=N_\mathrm{a}(v)𝑑v`$. We used the software of Sembach & Savage (1992) to measure the apparent column densities, so again, uncertainties due to the continuum placement and a 2% uncertainty in the flux zero point are included in the final error bars.
The H I Ly$`\alpha `$ and O VI $`\lambda `$1031.9 $`N_\mathrm{a}(v)`$ profiles are plotted in Figure 4, and their integrated apparent column densities are provided in Table 2. This table also lists $`4\sigma `$ upper limits on column densities of several species of interest obtained by integrating $`N_\mathrm{a}(v)`$ over the velocity range of the undetected lines. We also include a $`4\sigma `$ upper limit on C IV $`\lambda `$1548.2 from the Faint Object Spectrograph (FOS) observations of Jannuzi et al. (1998). Figure 4 shows that there appears to be a small velocity difference between the centroid of the strongest component of the H I profile and that of the O VI line. However, the H I and O VI velocities from profile fitting agree within their 1$`\sigma `$ uncertainties. More importantly, Figure 4 shows that the main component of the H I profile is broader than the O VI 1031.9 Å line, which suggests that the line widths are (at least partially) controlled by thermal motions (see §4.1). However, it is alternatively possible that the two strongest components of the H I profile include contributions from narrow features which are not detected in the O VI line. Such additional narrow lines, if present, are not well-constrained by the current data.
## 4 Physical Conditions
We now turn to the physical conditions of the absorbing gas. As discussed in §1, we are especially interested in testing the warm/hot gas prediction of cosmological simulations of structure growth. For this purpose we seek to constrain the gas temperature in §4.1, and we examine the ionization mechanism in §4.2.
### 4.1 Gas Temperature
If the line broadening is entirely due to thermal motions, then the gas temperature can be determined from the Doppler parameter of a given line,
$$T=mb^2/2k=A(b/0.129)^2,$$
(2)
where $`m`$ is the mass and $`A`$ is the atomic mass number of the element, the Doppler parameter $`b`$ is in km s<sup>-1</sup>, and $`T`$ is in K. However, it is usually true that additional factors such as gas turbulence or multiple unresolved components contribute to the width of a line, and consequently the gas temperature from equation (2) must be treated as an upper limit. If lines of two or more elements with adequately different masses are available, then $`b`$ can be expressed as
$$b^2=b_{nt}^2+(0.129)^2T/A,$$
(3)
which can be solved for $`T`$ and $`b_{nt}`$, the component of the broadening due to non-thermal motions, assuming the different absorption lines arise in the same gas and that the non-thermal motions (turbulence) can be assumed to have a Gaussian profile.
The temperature upper limits implied by equation (2) and the Doppler parameters in Table 3 are $`T`$(O VI) $`3.5\times 10^5^{}`$K and $`T`$(H I) $`5.8\times 10^4^{}`$K for the strongest component of the H I profile, which is also the closest component to the O VI line in velocity space (see Figure 4). Since the velocity centroids of the O VI and the H I are quite close and, in fact, in agreement within the 1$`\sigma `$ error bars, it is plausible that these O VI and H I lines originate in the same gas. Adopting this assumption, we derive from equation (3) $`T`$ = $`3.9\times 10^4^{}`$K and $`b_{nt}`$ = 18 km s<sup>-1</sup>.
While such low temperatures clearly favor photoionization (see below), we must bear in mind that there are considerable uncertainties in the Doppler parameters. For example, for the strongest component of the H I profile, increasing the $`b`$-value in Table 3 by just 1$`\sigma `$ increases the upper limit on $`T`$(H I) to $`8.7\times 10^4^{}`$K, much closer to the temperature range where O VI can be produced by collisional ionization. This issue is exacerbated by the non-uniqueness of the fitted profile when dealing with multiple blended components, i.e., assuming a different number of components or a different mix of broad and narrow components can lead to substantially different $`b`$-values. We provide an example of this in §4.2.2. It is also possible that the absorption arises in non-equilibrium collisionally ionized gas in which cooling is faster than recombination (see §4.2.2).
### 4.2 Ionization
Given the measurements in the previous sections, can we favor photoionization or collisional ionization in this O VI absorber? In the Milky Way ISM, O VI likely traces collisionally ionized hot gas (Jenkins 1978a,b), at least in the disk and lower halo. However, in extragalactic O VI absorbers, the EUV ionizing radiation field may be substantially harder, and the absorption might arise in very low density gas with a long path length. Both of these factors make photoionization more viable in the extragalactic case. We first consider the plausibility of photoionized models (§4.2.1), and then we test the collisional scenario (§4.2.2).
#### 4.2.1 Photoionization
We have explored the photoionized scenario using CLOUDY (version 90.04; Ferland et al. 1998) with the standard assumptions, in particular that there has been time for thermal and ionization equilibrium to prevail (see below). The absorber is treated as a constant density plane-parallel slab photoionized by the extragalactic radiation from QSOs and AGNs, as calculated by Haardt & Madau (1996) for $`z`$ 0.12. We set the mean intensity at the H I Lyman limit to $`J_\nu `$(LL) = 1 $`\times 10^{23}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>, a value in agreement with observational constraints (e.g., Kulkarni & Fall 1993; Maloney 1993; Tumlinson et al. 1999) and theoretical expectations (Shull et al. 1999). With these assumptions, we varied the metallicity<sup>8</sup><sup>8</sup>8We express the linear abundance of element X relative to element Y as (X/Y) and the logarithmic abundance in the usual fashion, \[X/Y\] = log (X/Y) – log (X/Y), and we indicate the overall metallicity with the variable $`Z`$. and ionization parameter ($`U=n_\gamma /n_\mathrm{H}`$ = H ionizing photon density/total hydrogen number density, neutral + ionized) to search for models which are consistent with the constraints set above: the H I and O VI column densities, limits on various column density ratios which can be derived using the upper limits in Table 2, and the gas temperature. The relative heavy element abundances (e.g., \[N/O\]) were initially set to the solar values from Grevesse & Anders (1989) and Grevesse & Noels (1993), then we considered an alternative N/O abundance (see below).
The basic results of the photoionization modeling are encapsulated in Figure 5. This figure shows various column density ratios involving O VI and other high ions of interest, plotted as function of the ionization parameter (bottom axis) and $`n_\mathrm{H}`$ (top axis). A metallicity of $`Z`$ = 1/10 $`Z_{}`$ was used to construct this model, but at low metallicities the column densities scale nearly directly with $`Z`$.<sup>9</sup><sup>9</sup>9As $`Z`$ approaches the solar value, the increased cooling provided by the metals affects $`T`$ and, in turn, the ionization balance of the gas. The gas temperature of the model as a function of $`U`$ is also indicated by a dotted line in Figure 5 with the linear temperature scale on the right axis. Of the various species that we are able to place limits on using the STIS E140M spectrum (see Table 2), only N V usefully constrains the photoionization model assuming \[N/O\] = 0; when the lower limit on $`N`$(O VI)/$`N`$(N V) is satisfied, the other species covered in the STIS E140M spectrum are predicted to have undetectably small column densities. Note that the STIS spectrum does not cover the C IV doublet at this redshift, and we discuss the constraint set by the FOS upper limit on C IV below. The O VI/N V column density ratio requires log $`U>`$ –1.26. Given this limit on $`U`$, the photoionized model requires $`n_\mathrm{H}10^{5.2}`$ cm<sup>-3</sup> and $`T`$ 31,200 K. At this ionization parameter, the O VI ionization fraction $`f`$(O VI) = O VI/O<sub>total</sub> = 0.164, and therefore the path length $`L`$ through the constant density absorbing region must be greater than $``$460 kpc to produce the observed O VI column density, log $`N`$(O VI) = 14.04, in 1/10 solar metallicity gas. With such a long path length, expansion of the universe should lead to an O VI line which is broader than the 1031.9 Å line shown in Figure 4, and the photoionized model predicts a substantially larger H I column density than observed. While the narrowness of the O VI line may be due to noise, the incorrect $`N`$(H I) rules out the 1/10 solar metallicity photoionized model. To match the observed $`N`$(O VI) and $`N`$(H I) with log $`U`$ = –1.26 and $`J_\nu `$(LL) = $`10^{23}`$ ergs<sup>-1</sup> s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>, the photoionization model requires $`Z1.6Z_{}`$ and $`L`$ 25 kpc, which also alleviates the problem with Hubble broadening and the width of the line. While this is an uncomfortably high metallicity to require, recent ASCA observations of X-ray bright galaxy groups indicate that in some groups, the intragroup gas metallicity is $``$1/3–1/2 solar (Davis, Mulchaey, & Mushotzky 1999; Hwang et al. 1999), so this may be marginally plausible. However, in such intragroup gas collisional ionization should be important.
Are these reasonable physical properties? Observations have established that some QSO absorbers can be quite extended. For example, the characteristic Ly$`\alpha `$ absorber “size” inferred from close QSO pairs is several hundred kpc at moderate redshifts (e.g., Dinshaw et al. 1995; Fang et al. 1996), and this is consistent with cosmological simulations of structure formation (e.g., Rauch, Haehnelt, & Steinmetz 1997; Davé et al. 1999). However, as various authors have argued on different grounds, these large absorbers are likely filamentary structures composed of smaller clouds rather than a single monolithic gas cloud (e.g., Rauch, Weymann, & Morris 1996; Cen & Simcoe 1997). Cen & Simcoe report that individual clouds have sizes of $``$50 kpc. Therefore the path length required by the photoionized model, $`L`$ 25 kpc, is reasonable. The low densities inferred for the O VI absorber at $`z_{\mathrm{abs}}`$ = 0.14232 are also consistent with cosmological simulations: according to equation (2) in Davé et al. (1998), the strongest component of the $`z_{\mathrm{abs}}`$ = 0.14232 absorber is expected to have $`n_\mathrm{H}10^{5.2}`$ cm<sup>-3</sup> based on its measured H I column density. Such low densities are also possible in the outer halo of a galaxy like the Milky Way. Murali (2000) has recently argued that the existence of the Magellanic Stream requires $`n_\mathrm{H}10^5`$ cm<sup>-3</sup> in the Milky Way halo at a Galactocentric distance of $``$50 kpc. Evidently, the size and density of this O VI absorber inferred from the photoionized model are plausible.
The most problematic property of the photoionized model is the required metallicity, $`Z1.6Z_{}`$. Most environments in which high metallicities are expected should involve collisional ionization or else should have substantially higher H I column densities. If we use expression (7) from Davé et al. (1999) to estimate the approximate overdensity associated with the $`z_{\mathrm{abs}}`$ = 0.14232 absorber based on its observed $`N`$(H I), we find that $`\rho _\mathrm{H}/\overline{\rho }_\mathrm{H}`$ 20. Based on this overdensity, the cosmological simulations of Cen & Ostriker (1999b) predict that the metallicity of this absorber should be $`Z0.1Z_{}`$. Similarly, to produce the absorption in a high metallicity region of the ISM of a galaxy like the Milky Way, the H I column density would have to be substantially higher than the observed $`N`$(H I).
We note that the H<sup>+</sup> recombination timescale, $`t_{\mathrm{rec}}(`$H$`{}_{}{}^{+})=1/\alpha (T)n_\mathrm{e}`$ where $`\alpha (T)`$ is the recombination coefficient, exceeds the age of the universe by a factor of several at such low densities, and one may wonder if modeling the gas with a photoionization equilibrium code is valid. However, the timescale for the gas to approach equilibrium, $`t_{\mathrm{eq}}=[t_{\mathrm{ion}}^1+2t_{\mathrm{rec}}^1]^1`$ where $`t_{\mathrm{ion}}`$ is the photoionization timescale (see Appendix A in Dove & Shull 1994), is vastly shorter than the recombination timescale in the conditions considered here, and photoionization equilibrium is a good approximation for the low density, highly ionized gas (see, e.g., Vedel, Hellsten, & Sommer-Larsen 1994).
A caveat in the photoionization models above is that we have assumed the solar \[N/O\] abundance. If instead we assume \[N/O\] $``$ –1.5 based on the trend of \[N/O\] vs. metallicity observed in giant extragalactic H II regions (Vila-Costas & Edmunds 1993), then the lower limit on $`N`$(O VI)/$`N`$(N V) is satisfied at a lower ionization parameter: the model requires log $`U>`$ –1.90. However, at lower values of $`U`$ the model predicts that Si III and Si IV are detectable, and in fact in this case the lower limit on $`N`$(O VI)/$`N`$(Si III) from Table 2 provides a more stringent constraint requiring log $`U>`$ –1.69 and $`n_\mathrm{H}10^{4.77}`$ cm<sup>-3</sup>. Similarly, the FOS upper limit on $`N`$(C IV) can be used to set tighter constraints on the photoionized model in this case: the $`N`$(O VI)/$`N`$(C IV) limit requires log $`U>`$ –1.33 and $`n_\mathrm{H}10^{5.13}`$ cm<sup>-3</sup>.
#### 4.2.2 Collisional Ionization
We next consider whether the gas could be collisionally ionized. In §4.1 we derived $`T5.8\times 10^4^{}`$K for the gas giving rise to the main component of the H I profile (the component at $`v`$ 0 km s<sup>-1</sup>). According to the collisional ionization equilibrium calculations of Sutherland & Dopita (1993), no O VI is produced by collisional ionization at this temperature, and in fact O II is the dominant ionization stage of oxygen. If we assume that the O VI and H I absorption lines arise in the same gas with the same turbulent broadening, we obtain $`T`$ = 3.9 $`\times 10^4^{}`$K, which makes collisional ionization seem rather unlikely. However, the absorber could be out of ionization and thermal equilibrium or it could be a multiphase medium, possibilities which are discussed below.
Furthermore, as noted in §4.1, there are substantial uncertainties in the $`b`$-values, and the H I Doppler parameter is consistent with gas having $`T1.6\times 10^5^{}`$K at the 3$`\sigma `$ level. From the O VI Doppler parameter we obtain $`T3.5\times 10^5^{}`$K. These temperatures are more in line with collisionally ionized gas. Furthermore, it is entirely possible that the H I profile contains a broad component at $`v`$ 0 km s<sup>-1</sup> superposed on a more narrow component. To demonstrate this, we show in Figure 6 two independent fits to the H I Ly$`\alpha `$ line. The model profile indicated with a solid line is the initial best fit which has the profile parameters summarized in Table 3. The model profile plotted with a dotted line used a very similar initial guess at the component parameters but was forced to include an additional component with the same velocity as the O VI line and $`b`$ = 76 km s<sup>-1</sup>, the Doppler parameter of H I in gas with $`T=3.5\times 10^5^{}`$K and no turbulent broadening. The two profile fits are nearly indistinguishable — by adjusting the parameters of the other components, the profile fitting code is able to compensate for the presence of the broad component to produce a very similar final result. The final column density required by the profile fitting software for the broad hot component is $`N`$(H I) = 1.4 $`\times 10^{13}`$ cm<sup>-2</sup>. Unless the spectrum has very good S/N, it would be easy to miss such a broad hot H I component. Therefore we cannot rule out the possibility that the gas is collisionally ionized on the basis of the H I Ly$`\alpha `$ profile properties.
Figure 6 is not shown in order to claim that such a broad component is present, but rather to show that it is allowed by the data. However, this exercise clearly shows that if the O VI absorption occurs in hot gas with $`T3\times 10^5^{}`$K, then the absorber must be a multiphase medium because cooler gas is required to account for the rest of the H I Ly$`\alpha `$ absorption including much of the absorption centered at $`v`$ 0 km s<sup>-1</sup>. Given sufficient sensitivity, metal lines such as C III $`\lambda `$977.02 and Si III $`\lambda `$1206.50 should be detectable in these cooler phases, but this may require substantial improvement over the current signal-to-noise due to the low H I column densities.
We next examine the high ion column density ratios predicted for collisionally ionized hot gas. Several useful high ion column density ratios from the collisional ionization equilibrium calculations of Sutherland & Dopita (1993) are plotted as a function of gas temperature in Figure 7. From this figure we see that if the gas is in collisional ionization equilibrium, then $`T`$ must be greater than $`2.3\times 10^5^{}`$K to satisfy the lower limit on $`N`$(O VI)/$`N`$(N V), assuming \[N/O\] = 0. If \[N/O\] = –1.5, then $`T1.8\times 10^5^{}`$K. These are plausible temperatures as discussed above.
However, the temperature range shown in Figure 7 is near the peak of the cooling function, and gas at these temperatures is likely to cool faster than it can recombine, even if it has low metallicity (Sutherland & Dopita 1993), leading to an overionized condition. Consequently, it may be more appropriate to compare the observed ratios to non-equilibrium calculations such as cooling fountain and cluster flow models (e.g., Edgar & Chevalier 1986; Benjamin & Shapiro 2000). Table 4 compares the observed lower limits on the O VI/Si IV, O VI/C IV, and O VI/N V column density ratios in the $`z_{\mathrm{abs}}`$ = 0.14232 absorber to the ratios predicted by non-equilibrium collisionally ionized models including the cooling fountain calculation of Benjamin & Shapiro (2000), the cluster cooling flow model of Edgar & Chevalier (1986), the turbulent mixing layer model of Slavin, Shull, & Begelman (1993), and the magnetized thermal conduction interface modeled by Borkowski, Balbus, & Fristrom (1990). From this table we see that the cooling fountain and conductive interface models are fully consistent with the current observational constraints on this absorber given the right choice of model parameters. The turbulent mixing layer is marginally inconsistent with the O VI/C IV limit, but this should be confirmed with higher sensitivity observations of $`N`$(C IV). However, we note that if the turbulent mixing layer post-mixing gas temperature is increased, then the model ironically becomes inconsistent with the O VI/Si IV limit as well, presumably due to increased production of Si IV by self-photoionization.
Given the available observational data, it is difficult to definitively constrain the ionization mechanism. It would be interesting to search for the Ne VIII $`\lambda \lambda `$ 770.41, 780.32 doublet in O VI absorbers. From Figures 5 and 7, one can see that in the equilibrium ionization models, O VI/Ne VIII $``$ 1 unless the gas has rather low density (in the photoionized case) or is relatively hot (if collisionally ionized). On the other hand, $`N`$(O VI) $`N`$(Ne VIII) in the non-equilibrium cooling fountain model (Table 4). Therefore detection of Ne VIII would provide evidence in favor of non-equilibrium collionsional ionization. Given sufficient $`z`$, the Ne VIII doublet will be redshifted into the bandpass of the Far Ultraviolet Spectroscopic Explorer or even the STIS $`\lambda `$ range. Unfortunately, $`z_{\mathrm{abs}}`$ = 0.14232 is inadequate, and the Ne VIII doublet is unobservable in this system. It is possible to look for the S VI $`\lambda \lambda `$ 933.38, 944.52 doublet at $`z_{\mathrm{abs}}`$ = 0.14232 with FUSE. However, in most models this doublet is predicted to be substantially weaker than the O VI lines (see Figures 5 and 7 and Table 4).
## 5 Discussion
### 5.1 Association with a Galaxy Group
The fact that the $`z_{\mathrm{abs}}`$ = 0.14232 O VI absorption line system is associated with a group of galaxies strongly indicates that the absorber is an intervening system rather than an intrinsic absorber which was ejected or somehow accelerated to high displacement velocity by the QSO. While this may seem a trivial conclusion since the O VI is displaced from the QSO redshift by $``$23,000 km s<sup>-1</sup>, during the last few years observations of absorption variability have established that some highly ionized intrinsic QSO absorbers are separated from $`z_{\mathrm{QSO}}`$ by such large velocities (e.g. Hamann, Barlow, & Junkkarinen 1997) and yet are relatively narrow (i.e., not traditional broad absorption line outflows). Therefore, it is important to find evidence that a given O VI absorber is indeed intervening even if $`\mathrm{\Delta }v`$ is large. Of course, this association with a galaxy group does not necessarily indicate that the absorption arises in the intragroup medium; it could be due to gas within one of the galaxies in the group. We briefly discuss some possibilities below.
The close proximity of the O VI system to a galaxy group provides a theoretical prejudice in favor of collisional ionization. The presence of galaxies requires that intergalactic gas has collapsed in this region of space, and simple arguments (e.g., §3 in Cen & Ostriker 1999a) suggest that substantial shock-heating probably occurred as a result. Therefore collisionally ionized gas is expected in the vicinity (if not along the pencil-beam probed by the QSO). It is interesting that this absorption system and associated galaxies fit the prediction of Mulchaey et al. (1996) discussed in §1: a group which is possibly spiral-rich and has associated O VI absorption, as expected based on their postulated collisionally ionized intragroup medium. However, we really are not sure that this is a spiral-rich group (or even a bound group) with only four known galaxies, and it is also possible that the O VI absorption arises in the gaseous halo of a single galaxy. However, models of galaxy gaseous halos (e.g., Mo & Miralda-Escudé 1996) also usually produce O VI absorption in a collisionally ionized hot phase. Similarly, models which produce QSO absorption line gas in supernova-driven winds from dwarf galaxies (e.g. Wang 1995) also involve substantial shock-heating and collisionally ionized gas.
The kinematics of the absorption also provide useful information about the possible origins of the system at $`z_{\mathrm{abs}}`$ = 0.14232. A fundamental question is whether the absorption is due to the intragroup medium or to the ISM of a single galaxy that happens to intercept the line-of-sight. The H I Ly$`\alpha `$ profile has the “leading edge asymmetry” that various authors (e.g., Lanzetta & Bowen 1992) have discussed as the signature of a moderately edge-on rotating disk. Of course, this is not a unique interpretation of such profile asymmetry, but it provides some evidence in favor of the single-galaxy interpretation. The simplicity of the O VI profile also favors this interpretation. The velocity dispersion of poor galaxy groups is typically one to a few hundred km s<sup>-1</sup> (Zabludoff & Mulchaey 1998), so one might expect the O VI to be spread over a larger velocity range if the absorption is due to the intragroup medium. Here, though, we must recognize that the observed O VI profiles are noisy and we only detect O VI at the velocity of the strongest component of the H I profile. If the hot gas is concentrated at the center of the group, then this may be the only location along the line-of-sight where we have sufficient sensitivity to detect it, and there may be O VI absorption at other velocities which has fallen below our detection threshold. To sort out the various possibilities, it would be very helpful to obtain additional STIS observations to improve the signal-to-noise and search for O VI at other velocities.
### 5.2 Number Density and Cosmological Mass Density
One means to test the warm/hot gas prediction of cosmological simulations (§1) is to compare the number of O VI absorbers observed per unit redshift, $`dN/dz`$, to the number statistically predicted from many random pencil-beams through the cosmological simulations. Is the observed $`dN/dz`$ consistent with the number predicted by the cosmological models? We can also estimate the mean cosmological mass density traced by the O VI systems at low $`z`$.
For these purposes, we combine the STIS observations of PG0953+415 with the GHRS observations of H1821+643 from Tripp et al. (1998).<sup>10</sup><sup>10</sup>10After this paper was completed, Tripp, Savage, & Jenkins (2000) carried out an analysis of new STIS echelle observations of H1821+643 with the E140M mode, and we refer the reader to that paper for the O VI number density and cosmological mass density derived from an independent data set. We emphasize that the manner in which these QSOs were selected should not bias the sample to enhance the number of O VI systems detected compared to sampling many random directions. Both QSOs were originally observed to study the relationship between weak Ly$`\alpha `$ clouds and galaxies<sup>11</sup><sup>11</sup>11$`HST`$ program IDs 6155 and 7747, see Tripp et al. (1998) for details. at $`z<`$ 0.3. These two objects were selected for this project simply because they are among the brightest known QSOs with $`z>`$ 0.2, criteria which were required to substantially increase the sample of weak Ly$`\alpha `$ absorbers with a minimal amount of HST time. No consideration was given to factors which might indicate an enhanced likelihood of detecting the O VI doublet when the targets were selected.
To estimate $`dN/dz`$ of the O VI systems, we define a sample of O VI lines with $`W_\mathrm{r}>`$ 60 mÅ for both lines of the doublet,<sup>12</sup><sup>12</sup>12The 3$`\sigma `$ detection limit throughout the region where O VI absorbers can be detected is 60 mÅ or better in the spectrum of PG0953+415 and 50 mÅ or better in the spectrum of H1821+643. and set the maximum absorber redshift, $`z_{\mathrm{max}}`$, for each sight line to exclude any absorbers within $`\mathrm{\Delta }v`$ 5000 km s<sup>-1</sup> of the QSO redshift and thereby avoid contaminating the sample with associated/intrinsic absorbers which are close to the QSO. One might argue that 5000 km s<sup>-1</sup> is insufficient (see above). However, we find that the two O VI systems in the final sample are associated with galaxy groups, which suggests that these are indeed intervening absorbers. The lower redshift cutoff for each sight line, $`z_{\mathrm{min}}`$, was determined by the lowest wavelength in the observed spectrum (with a small buffer to ensure that a line would be recognized if at that $`\lambda `$) in the case of H1821+643 and by the wavelength at which the S/N is unacceptably low in the case of PG0953+415. These criteria resulted in a sample of two O VI absorbers within a total redshift path of $`\mathrm{\Delta }z`$ = 0.100 (after a significant correction of 0.067 for regions of the spectra in which we cannot detect either of the O VI lines<sup>13</sup><sup>13</sup>13Both of the O VI lines were required to fall in unblocked regions of the spectra so that the doublet could be securely identified. because they would be blocked by strong ISM or extragalactic lines from other redshift systems). Therefore, $`dN/dz`$ 20 for O VI systems with $`W_\mathrm{r}`$ 60 mÅ at $`z<`$ 0.3. Using the confidence limits from Gehrels (1986) for a sample of two absorbers, we derive 4 $`<dN/dz<`$ 63 for these O VI absorbers at the 90% confidence level.
For comparison, Tripp et al. (1998) find $`dN/dz=102\pm 16`$ for H I Ly$`\alpha `$ lines with $`W_\mathrm{r}`$ 50 mÅ at $`z<`$ 0.3, and Weymann et al. (1998) report $`dN/dz=33\pm 4`$ for Ly$`\alpha `$ absorbers at $`z`$ = 0 with $`W_\mathrm{r}`$ 240 mÅ based on a large sample at $`z<`$ 1.5. In the case of low to moderate redshift Mg II absorbers, $`dN/dz=0.97\pm 0.10`$ for $`W_\mathrm{r}`$ 300 mÅ (Steidel & Sargent 1992) and $`dN/dz=2.65\pm 0.15`$ for $`W_\mathrm{r}`$ 20 mÅ (Churchill et al. 1999; see also §6 in Tripp, Lu, & Savage 1997). Evidently, these weak low $`z`$ O VI systems have a substantially larger cross section and/or covering factor than the Mg II absorbers. Similarly, the stronger O VI absorbers at higher redshifts are less common: Burles & Tytler (1996) report $`dN/dz=1.0\pm 0.6`$ for O VI systems with $`W_\mathrm{r}`$ 210 mÅ at $`<z_{\mathrm{abs}}>`$ = 0.9, and similar results (with smaller uncertainties) are derived from the larger sample provided by the FOS Quasar Absorption Line Key Project (B. Jannuzi & R. Weymann 2000, private communication). This comparison with the O VI $`dN/dz`$ derived from FOS data must be interpreted carefully, however, because the FOS samples are dominated by lines which are substantially stronger and at substantially higher redshifts than the O VI absorbers discussed in this paper. With a larger sample of low $`z`$ weak O VI absorbers and smaller uncertainties in their number density, comparison of $`dN/dz`$ to the space density of objects such as dwarf galaxies may provide insight on the nature of the O VI systems.
Next we estimate the baryonic content of the O VI absorbers, expressed as the cosmological mass density $`\mathrm{\Omega }_b`$(O VI), following previous analogous calculations for damped Ly$`\alpha `$ systems (e.g., Lanzetta et al. 1991) as well as O VI absorbers (Burles & Tytler 1996). To estimate the density of baryons in the gaseous component of the universe traced by O VI, we require information about the metallicity of the gas and the O VI ionization fraction, $`f`$(O VI) = O VI/O<sub>total</sub>. In collisional ionization equilibrium, $`f`$(O VI) peaks at $``$0.2 (Sutherland & Dopita 1993), and similarly low peak fractions are predicted by non-equilibrium collisional models (e.g., Shapiro & Moore 1976; Benjamin & Shapiro 2000). The O VI ion fraction is not much larger at peak value in photoionized gas (see appendix). Therefore we will adopt $`f`$(O VI) $``$ 0.2 to set a lower limit on the baryonic content of the O VI absorbers, and the following calculation is relatively independent of how the gas is ionized or whether or not it is due to the intragroup medium. For the mean cosmic metallicity of the O VI absorbers, it is less clear what value to adopt, but to set a lower limit on $`\mathrm{\Omega }_b`$(O VI) we should set (O/H) to a high but plausible value. As noted in §4.2.1, Davis et al. (1999) and Hwang et al. (1999) have derived metallicities of 1/3–1/2 solar for intragroup gas in several X-ray bright groups. Therefore we will initially use 1/2 solar metallicity for the calculation of $`\mathrm{\Omega }_b`$(O VI) and then discuss how it scales with (O/H). The mean cosmological mass density in the O VI absorbers, in units of the current critical density $`\rho _c`$, can be estimated as
$$\mathrm{\Omega }_b(\mathrm{O}\mathrm{VI})=\frac{\mu m_\mathrm{H}H_0}{\rho _ccf(\mathrm{O}\mathrm{VI})}\left(\frac{\mathrm{O}}{\mathrm{H}}\right)_{\mathrm{O}\mathrm{VI}}^1\frac{_iN_i(\mathrm{O}\mathrm{VI})}{_i\mathrm{\Delta }X_i}$$
(4)
where $`\mu `$ is the mean atomic weight (taken to be 1.3), (O/H)$`_{\mathrm{O}\mathrm{VI}}`$ is the assumed mean oxygen abundance by number in the O VI absorption systems, $`m_\mathrm{H}`$ is the mass of hydrogen, $`N_i`$(O VI) is the total O VI column density and $`\mathrm{\Delta }X_i`$ is the absorption distance interval (Bahcall & Peebles 1969) probed to the $`i`$th QSO,
$$\mathrm{\Delta }X_i=\frac{1}{2}\{[(1+z_{\mathrm{max}})^21][(1+z_{\mathrm{min}})^21]\}$$
(5)
assuming $`q_0`$ = 0.<sup>14</sup><sup>14</sup>14Over the redshift range probed by the sight lines to PG0953+415 and H1821+643, results are insensitive to the value assumed for $`q_0`$. As in the calculation of $`dN/dz`$, we correct $`\mathrm{\Delta }X_i`$ for spectral regions blocked by strong lines. Combining the PG0953+415 STIS data and the GHRS observations of H1821+643 from Tripp et al. (1998) with $`z_{\mathrm{min}}`$ and $`z_{\mathrm{max}}`$ set to the same values used for the derivation of $`dN/dz`$, we obtain $`\mathrm{\Omega }_b(\mathrm{O}\mathrm{VI})0.0006h_{75}^1`$ assuming the mean O abundance is 1/2 solar. This is a lower limit not only because $`f`$(O VI) and (O/H) were set to their approximate upper limits, but also because we have applied an equivalent width cutoff to define the sample; if O VI absorbers with $`W_\mathrm{r}`$ 60 mÅ significantly increase $`_iN_i`$(O VI), then the true $`\mathrm{\Omega }_b`$(O VI) will be higher. Note that $`\mathrm{\Omega }_b(\mathrm{O}\mathrm{VI})`$ is inversely proportional to (O/H). Decreasing the metallicity to 1/10 solar, for example, increases the lower limit on the O VI absorber baryon content to $`\mathrm{\Omega }_b(\mathrm{O}\mathrm{VI})0.003h_{75}^1`$.
To demonstrate the level of uncertainty in the cosmological mass density estimate due to small number statistics, we can recalculate $`\mathrm{\Omega }_b`$(O VI) using an alternative expression analogous to equation (9) from Rao & Turnshek (2000),
$$\mathrm{\Omega }_b(\mathrm{O}\mathrm{VI})=\frac{\mu m_\mathrm{H}H_0}{\rho _ccf(\mathrm{O}\mathrm{VI})}\left(\frac{\mathrm{O}}{\mathrm{H}}\right)_{\mathrm{O}\mathrm{VI}}^1\left(\frac{dN}{dz}\right)\frac{<N_{\mathrm{O}\mathrm{VI}}>}{(1+z)}$$
(6)
where $`<N_{\mathrm{O}\mathrm{VI}}>`$ is the mean O VI column density of the absorption systems in the sample and again we have assumed $`q_0`$ = 0. The advantage of this alternative expression for $`\mathrm{\Omega }_b`$(O VI) is that we can employ the Gehrels (1986) small sample statistics to estimate the uncertainty in $`dN/dz`$ and then propagate the uncertainty into the estimate of $`\mathrm{\Omega }_b`$(O VI). With $`dN/dz`$ = $`20_{13}^{+26}`$ and the other parameters set to the values assumed above, we find for the 1/10 solar metallicity case $`\mathrm{\Omega }_b(\mathrm{O}\mathrm{VI})0.003_{0.002}^{+0.004}h_{75}^1`$ (error bars are $`1\sigma `$ uncertainties).
Bearing in mind that there is still considerable uncertainty in the lower limit on $`\mathrm{\Omega }_b`$(O VI) due to the small sample, small redshift path, and uncertain mean metallicity, this preliminary estimate suggests that the O VI absorbers may indeed harbor a significant fraction of the baryons in the universe at low $`z`$. The lower limit assuming (O/H) = 1/10 solar is comparable to the cosmological mass density of stars, H I, and X-ray emitting galaxy group and cluster gas at low redshift (Fukugita, Hogan, & Peebles 1998), for example.
## 6 Summary
The paper is summarized as follows.
(1) We have observed the low redshift QSO PG0953+415 with the E140M mode of STIS, and this has revealed an O VI absorption line system at $`z_{\mathrm{abs}}`$ = 0.14232. This O VI absorber is highly ionized: multicomponent H I absorption is also detected at this redshift, but no other species are detected in the STIS spectrum including the N V doublet, Si III, and C II.
(2) We have measured galaxy redshifts in the field of the QSO using the WIYN telescope, and there are at least four galaxies within $``$130 km s<sup>-1</sup> of the O VI absorber with projected distances ranging from 395 kpc to 3.0 Mpc. Two of these galaxies appear to be spiral galaxies.
(3) A review of the observational constraints shows that we cannot definitively assert that the gas is collisionally ionized or photoionized, but photoionization requires an uncomfortably high metallicity which is inconsistent with theoretical expectations. If the gas is collisionally ionized, it is likely that it is not in equilibrium or that it is a multiphase absorber. Non-equilibrium collisional ionization models are consistent with the observations.
(4) Combining the STIS data on PG0953+415 with the high S/N low resolution GHRS observations of H1821+643 from Tripp et al. (1998), we identify two intervening O VI systems over a redshift path $`\mathrm{\Delta }z`$ of only 0.10. This implies that $`dN/dz`$ for O VI systems with $`W_\mathrm{r}>`$ 60 mÅ and $`z<`$ 0.3 is $``$20 with a large uncertainty due to the small number of systems so far detected. This represents a large density for metal line systems since $`dN/dz`$ for Ly$`\alpha `$ absorbers with $`W_\mathrm{r}>`$ 50 mÅ and $`z<`$ 0.3 is $`102\pm 16`$ (Tripp et al. 1998). We stress that the sample should not be biased in favor of O VI detection. If further observations confirm that $`dN/dz`$ is as large as 20 for O VI systems, then these absorbers may be an important baryon reservoir at low redshift, although this depends on the metallicity of the gas. If the mean metallicity is 1/2 solar, then $`\mathrm{\Omega }_b(\mathrm{O}\mathrm{VI})0.0006h_{75}^1`$. However, if the mean metallicity is 0.1 solar, then $`\mathrm{\Omega }_b(\mathrm{O}\mathrm{VI})0.003h_{75}^1`$, which is comparable to the baryonic content of other known constituents of the low $`z`$ universe such as galaxies and gas in galaxy clusters.
This research has made use of software developed by the STIS Instrument Definition Team for the reduction of STIS data, and we thank the STIS team for allowing us to use their software. We also benefitted from the use of CLOUDY, and we are indebted to Gary Ferland for sharing this program in which he has invested years of development effort. Similarly, we extend our thanks to Ed Fitzpatrick for the use of his profile fitting code and to Ken Sembach for his apparent column density software. Valuable comments were provided by Dave Bowen, Romeel Davé, Bruce Draine, Ed Jenkins, John Mathis, and Linda Sparke. We especially appreciate the careful review of the manuscript provided by Jane Charlton which clarified and improved the paper. We also thank Francisco Haardt for private communication of the Haardt & Madau background radiation fields in a convenient digital format, and Buell Jannuzi and Ray Weymann for private communication regarding the number density of O VI systems in the FOS Quasar Absorption Line Key Project. Finally, this research has made use of the APS Catalog of POSS I, which is supported by the NSF, NASA, and the University of Minnesota. The APS databases can be accessed at http://aps.umn.edu/. T. M. T. acknowledges support from NASA through grant NAG5–30110. B. D. S. acknowledges NASA support through grant GO–06499.02–95A from the Space Telescope Science Institute.
## Appendix A Maximum Ionization Fraction of O VI
In §5.2 we adopted $`f`$(O VI) = 0.2 to set a lower limit on the baryonic content of the O VI absorbers. Intuitively, $`f`$(O VI) = 1.0 would seem to set a more conservative lower limit on $`\mathrm{\Omega }_b`$(O VI). However, as noted in §5.2, O VI is not a preferred ionization stage of oxygen, and in collisionally ionized gas $`f`$(O VI) peaks at $``$0.2 in equilibrium and non-equilibrium calculations.
To evaluate the maximum O VI ionization fraction which can be attained if the O VI is created by photoionization, we have calculated with CLOUDY an O VI ionization fraction grid as a function of the temperature and ionization parameter of a parcel of gas assuming the Haardt & Madau (1996) UV background radiation at $`z`$ = 0.12. To show how $`f`$(O VI) depends on $`T`$ and $`U`$, we operated CLOUDY in a different mode than that employed in §4.2.1: instead of allowing CLOUDY to determine the gas temperature given the various heating and cooling processes which are important (as was done in §4.2.1), we fixed the temperature and ionization parameter at particular values in each grid cell and calculated the ensuing $`f`$(O VI) for a grid with 3.0 $``$ log $`T`$ 7.0 and –3.5 $``$ log $`U`$ 0.5. The resulting grid is shown in Figure 8.
As expected based on previous calculations, when the ionization parameter is low enough so that photoionization is unimportant and O VI only has a significant ionization fraction in gas which is hot enough to be collisionally ionized, Figure 8 shows that $`f`$(O VI) has a maximum value of $``$0.2 at log $`T`$ = 5.5. On the other hand, when the ionization parameter is high and photoionization dominates, Figure 8 indicates that $`f`$(O VI) could be somewhat larger if the gas is very cool. However, this is rather unlikely to occur because the photoionization by the UV background which creates the O VI will also heat the gas to log $`T`$ 4. For example, the CLOUDY model in Figure 5 which satisfies the O VI and N V constraints requires log $`T`$ 4.5 (see §4.2.1), and with this temperature constraint the maximum $`f`$(O VI) is 0.28. Therefore $`f`$(O VI) = 0.2 is a reasonable value for placing a lower limit on the baryonic content of O VI absorbers regardless of whether the gas is photoionized or collisionally ionized.
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# Weak localization correction to the FS interface resistance
## I Introduction
Novel phenomena in mesoscopic systems have been the subject of intense theoretical and experimental interest for many years. Observed effects such as weak localization and universal conductance fluctuations are due to the quantum interference of electrons at low temperatures. An important branch of mesoscopic physics is that of hybrid nanostructures where the influence of a superconductor (S) on a phase coherent normal (N) region is studied. At low temperatures, the role of Andreev reflection is important, whereby particles in the N region with excitation energies $`\epsilon `$ lower than the superconducting gap energy $`\mathrm{\Delta }`$ are reflected from the N/S interface as holes. Studies of subgap transport have shown that superconducting condensate penetration into the N region, known as the proximity effect, may substantially change the resistance of an N/S circuit.
Lately, improvements in the microprocessing of metals has led to the fabrication of nanostructures combining ferromagnetic (F) and S materials. The presence of a large exchange field, $`\epsilon _{\mathrm{ex}}\mathrm{\Delta }`$, suppresses electron-hole correlations in ferromagnets so that the role of the proximity effect is reduced. However a separate mechanism has been predicted to produce a resistance increase in a circuit consisting of a mono-domain F wire connected to an S electrode instead of an N electrode. This is a robust, classical effect which originates from the need to match a spin-polarized current in the ferromagnet to a spinless current in the superconductor.
The present work analyses the weak localization correction to the contact F/S resistance. In a semiclassical language, weak localization arises from an enhanced backscattering caused by the quantum interference of pairs of coherent quasi-particles. Constructive interference of a quasi-particle paired with its time reversed partner is destroyed by a magnetic field so we assume that the width of the F wire, $`L_{}`$, is small, $`L_{}100`$nm, in order to limit the influence of the intrinsic magnetic field. Weak localization is affected by the boundary conditions of the system. Particles that escape into an N electrode suffer dephasing which reduces the return probability. However, multiple Andreev reflections from an S electrode may enable a particle to return coherently. For a polydomain F wire, with no net average polarization, the return probability of an F/S system is greater than that of an F/N system.
In general, the spin channels have different conductivities, $`\sigma _\alpha =e^2\nu _\alpha D_\alpha `$, where $`\nu _\alpha `$ and $`D_\alpha `$ are the density of states and the classical diffusion coefficient for electrons in the spin state $`\alpha `$, $`\alpha =(,)`$. The degree of spin polarization, $`\zeta `$, is defined as $`\zeta =\left(\sigma _{}\sigma _{}\right)/\left(\sigma _{}+\sigma _{}\right)`$ and we calculate the return probability in an F/S system with an arbitrary degree of polarization, $`0\zeta 1`$. In addition we introduce a finite spin relaxation length, $`L_s`$, shorter than the length, $`L`$, of the F wire, $`L_sL`$. As the spin polarization increases, the majority spin channel experiences a further increase in return probability (as compared to the F/N case), whereas the minority spin channel suffers a reduction. Majority carriers at the superconducting interface cannot find minority carrier states to Andreev reflect into and are more likely to be reflected normally whereas minority carriers at the interface have an enhanced probability to escape. For $`\zeta 1`$ the reservoir appears to be totally insulating to majority spins but normal as far as minority spins are concerned.
After calculating the correction to the classical diffusion coefficient that is due to the enhanced return probability, we find the corresponding change in the spin polarized particle distribution and we determine the weak localization correction to the classical resistance. Generally this is a sum of a bulk term that is independent of the state of the reservoir and a contact term, $`\delta r_c^{FS(FN)}`$, that depends on whether the reservoir is in the superconducting or the normal state. We find that $`\delta r_c^{FS}>\delta r_c^{FN}`$, so that a change in the state of the reservoir from N to S results in an increase of the weak localization contribution to the resistance,
$$\delta r_c^{FS}\delta r_c^{FN}\gamma \left(\frac{8e^2}{\pi }\right)\frac{\left(\sigma _{}+\sigma _{}\right)^2}{\sigma _{}\sigma _{}}\left(R_{\mathrm{}}\frac{L_s}{L_{}}\right)^2,$$
(1)
where $`R_{\mathrm{}}=1/(\sigma _{}+\sigma _{})L_{}^{d2}`$ is the resistance of a cube of length equal to the width of the wire $`L_{}`$. The combination $`R_{\mathrm{}}L_s/L_{}`$ corresponds to the resistance of a piece of ferromagnet of length $`L_s`$. We evaluate the prefactor, $`\gamma `$, numerically and find that it is almost constant, $`\gamma 1/2`$, for all experimentally relevant values of spin polarization.
## II Classical Resistance
We begin by briefly describing the approach of Refs. which enables one to calculate the classical contact resistance of an F/S interface and we refer the reader to those papers for further details. We consider a single domain, disordered ferromagnetic wire of length $`L`$ as depicted in Fig. 1. The coordinate $`x`$ is used to describe the position along the wire and throughout the paper we consider there to be a normal reservoir at the left hand side, $`x=L`$, whereas the reservoir at the right hand side, $`x=0`$, is either N or S. The resistance of the disordered F wire is determined by solving diffusion equations for the isotropic part of the electron distribution function $`n_\alpha (\epsilon ,x)=𝑑\mathrm{\Omega }_𝐩n_\alpha (𝐩,x).`$ Using electron-hole symmetry, it is possible to consider only the symmetrised function $`N_\alpha (\epsilon ,x)=[n_\alpha (\epsilon ,x)+n_\alpha (\epsilon ,x)]/2,`$ where $`\epsilon `$ is the energy determined with respect to the chemical potential in the S (N) reservoir. The current density in any given spin channel is related to the spatial derivative of the distribution function as
$$j_\alpha =\sigma _\alpha _0^{\mathrm{}}𝑑\epsilon _xN_\alpha (\epsilon ,x).$$
(2)
The distribution function obeys a differential equation describing diffusion in the ferromagnetic wire,
$$D_\alpha _x^2N_\alpha (\epsilon ,x)=w_{}\nu _{\overline{\alpha }}\left[N_\alpha (\epsilon ,x)N_{\overline{\alpha }}(\epsilon ,x)\right],$$
(3)
where $`\overline{\alpha }=(,)`$ for $`\alpha =(,)`$. This equation may be expressed conveniently as
$`_x^2\left[\sigma _{}N_{}+\sigma _{}N_{}\right]`$ $`=`$ $`0,`$ (4)
$`\left(_x^2L_s^2\right)\left[N_{}N_{}\right]`$ $`=`$ $`0.`$ (5)
Spin relaxation, which may arise from spin-orbit scattering or spin-flip scattering at defects, is described by the right-hand side of Eq. (3). One may define an effective spin relaxation rate, $`L_s`$, as $`L_s^2=w_{}[\nu _{}/D_{}+\nu _{}/D_{}]`$ which accounts for the relaxation of the difference between the spin channels. In a similar way, the spin relaxation length in an individual channel, $`L_\alpha `$, is defined as $`L_\alpha ^2=w_{}\nu _{\overline{\alpha }}/D_\alpha `$ such that $`L_s^2=L_{}^2+L_{}^2`$. An equivalent expression is $`L_\alpha ^2=L_s^2(\sigma _\alpha +\sigma _{\overline{\alpha }})/\sigma _{\overline{\alpha }}`$
The above equations should be complemented by four boundary conditions, two at each end of the ferromagnetic wire. The boundary conditions at the left hand side of the wire are given by the equilibrium distribution of electrons in the normal reservoir,
$$N_\alpha (\epsilon ,L)=\frac{1}{2}\left[n_T(\epsilon eV)+n_T(\epsilon eV)\right],$$
(6)
such that the distribution of up and down spins is equal. The boundary conditions at the right hand side of the ferromagnetic wire at the superconducting reservoir have been discussed in detail in Refs. . For quasiparticles with energies below the superconducting gap they may be written as
$`\sigma _{}_xN_{}|_0`$ $`=`$ $`\sigma _{}_xN_{}|_0,`$ (7)
$`N_{}(\epsilon ,0)+N_{}(\epsilon ,0)`$ $`=`$ $`\left[n_T(\epsilon )+n_T(\epsilon )\right].`$ (8)
The first of these, Eq. (7), describes Andreev reflection such that the spin-up current must equal the spin-down current and the second, Eq. (8), states that the sum of the distributions is equal to the equilibrium value. On solving the above differential equations with the appropriate boundary conditions it is possible to calculate the current density in the wire. In particular, the spatial derivative of the distribution is
$`_xN_\alpha (\epsilon ,x)=`$ $`{\displaystyle \frac{(12N_L)}{2\mathrm{\Gamma }L}}`$ (10)
$`\times \left[1+{\displaystyle \frac{(\sigma _{\overline{\alpha }}\sigma _\alpha )}{2\sigma _\alpha }}{\displaystyle \frac{\mathrm{cosh}[(x+L)/L_s]}{\mathrm{cosh}(L/L_s)}}\right],`$
where
$$\mathrm{\Gamma }=1+\frac{(\sigma _{}\sigma _{})^2}{4\sigma _{}\sigma _{}}\frac{L_s}{L}\mathrm{tanh}(L/L_s),$$
(11)
and $`N_LN_\alpha (\epsilon ,L)`$.
The effect of the superconducting reservoir, for strong spin relaxation $`LL_s`$, is expressed in terms of a contact resistance $`r_c^{FS}`$ so that the total classical resistance, $`R_{cl}`$, may be written as a sum of resistances in series
$$R_{cl}=\frac{L}{L_{}}R_{\mathrm{}}+r_c^{FS},$$
(12)
$$r_c^{FS}R_{\mathrm{}}\frac{L_s}{L_{}}\frac{\left(\sigma _{}\sigma _{}\right)^2}{4\sigma _{}\sigma _{}}.$$
(13)
For a normal reservoir on the right hand side, the boundary conditions are similar to those in Eq. (6), namely $`N_\alpha (\epsilon ,0)=\left[n_T(\epsilon )+n_T(\epsilon )\right]/2`$. In this case, there is no contact resistance. A fall in temperature, precipitating a change in state of the reservoir from N to S, would therefore result in a resistance increase of the circuit. This is a robust classical effect which originates from the need to match a spin-polarized current in the ferromagnet to a spinless current in the superconductor. It is specific to mono-domain wires since in poly-domain wires the tranport properties of spin-up and spin-down particles do not differ and, on the average, the current is not spin-polarized.
## III Weak localization correction to the contact resistance
The weak localization correction to the classical diffusion coefficient, $`\delta D_\alpha `$, is related to a two particle Green’s function known as the Cooperon propagator. The Cooperon consists of a quasi-particle following a diffusive path that interferes with a quasi-particle traversing the same path in the opposite direction. When the particles follow a closed path, the interference results in an enhanced return probability. In a ferromagnet, singlet and $`S_z=0`$ triplet Cooperons are suppressed since the Fermi momentum for spin-up and spin-down particles is different, producing a difference in the phase accumulated by the particles along the path. However, there is no such suppression of the phase correlation in the triplet channel where same spin particles are paired. Thus we assume that the return probability is given by the triplet Cooperon propagator, $`C_{\alpha \alpha }(x,x^{})`$, for equal spatial coordinates,
$$\frac{\delta D_\alpha (x)}{D_\alpha ^0}\frac{8}{\pi L_{}^2}C_{\alpha \alpha }(x,x),$$
(14)
which obeys a diffusion equation in the disordered ferromagnet,
$$\left(\nu _\alpha D_\alpha _x^2+\nu _\alpha \tau _\alpha ^1\right)C_{\alpha \alpha }(x,x^{})=\delta \left(xx^{}\right).$$
(15)
Here the spin relaxation time, $`\tau _\alpha `$, is related to the spin relaxation length in the spin channel $`\alpha `$, $`\tau _\alpha =L_\alpha ^2/D_\alpha `$.
We stress that the existence of a Cooperon is not a result of the proximity effect as we do not consider triplet pairing induced in the ferromagnet by the presence of the superconductor. Nevertheless, weak localization is affected by the boundary conditions of the system. A schematic of the change in the boundary conditions of the Cooperon at the F/S interface as compared to the F/N interface is shown in Fig. 2. A process which pairs two spin-up quasi-particles is depicted in Fig. 2a. When the reservoir is in the N state, particles escaping into it suffer dephasing, thus destroying the Cooperon. At the left hand reservoir ($`x=L`$), which is in the N state, we therefore have
$$C_{}(L,x^{})=C_{}(L,x^{})=0.$$
(16)
When the right hand reservoir ($`x=0`$) is in the S state we apply
$`C_{}(0,x^{})`$ $`=`$ $`\mathrm{e}^{i\chi }C_{}(0,x^{}),`$ (17)
$`\sigma _{}_xC_{}(x,x^{})|_{x=0}`$ $`=`$ $`\sigma _{}\mathrm{e}^{i\chi }_xC_{}(x,x^{})|_{x=0},`$ (18)
where $`\chi `$ is a phase accumulated on reflection from the superconductor. These boundary conditions, similar to those for the density Eqs.(7,8), account for Andreev reflection. Multiple Andreev reflections at an S electrode may enable a particle to return coherently to the original place in the F wire with the same spin polarization. Fig. 2b shows a process whereby a pair of spin-up quasi-particles are Andreev reflected as spin-down quasi-holes which are subsequently Andreev reflected as spin-up quasi-particles. The calculation of $`\delta D_\alpha `$ by solving the Cooperon diffusion equation with the above boundary conditions is described in Appendix A.
The weak localization correction to the diffusion coefficient, $`\delta D_\alpha `$, leads to a correction to the current density, $`\delta j_\alpha =e^2\nu _\alpha \delta D_\alpha _0^{\mathrm{}}𝑑\epsilon _xN_\alpha (\epsilon ,x)`$. However, $`\delta D_\alpha `$ also produces a correction to the particle distribution function. We take this into account perturbatively by assuming that the correction, $`\delta D_\alpha `$, to the classical diffusion coefficient, $`D_\alpha ^0`$, is small and that there is a corresponding small correction, $`\delta N_\alpha `$, to the classical particle distribution, $`N_\alpha ^0`$. Expressions for the total diffusion coefficient, $`D_\alpha =D_\alpha ^0+\delta D_\alpha `$, and the total particle distribution, $`N_\alpha =N_\alpha ^0+\delta N_\alpha `$, are substituted into the previous differential equations describing diffusion in the ferromagnet, Eqs. (4,5), and the boundary conditions, Eqs. (6,7,8), to provide new equations for $`\delta N_\alpha `$. Hence the differential equations for the correction to the distribution function, $`\delta N_\alpha `$, in terms of the correction to the classical diffusion coefficient, $`\delta D_\alpha `$, are
$`_x^2\left[\sigma _{}\delta N_{}+\sigma _{}\delta N_{}\right]`$ $`=`$ $`_x\left[\delta 𝒥_{}+\delta 𝒥_{}\right],`$ (19)
$`\left(_x^2L_s^2\right)\left[\delta N_{}\delta N_{}\right]`$ $`=`$ $`_x\left[{\displaystyle \frac{\delta 𝒥_{}}{\sigma _{}}}{\displaystyle \frac{\delta 𝒥_{}}{\sigma _{}}}\right],`$ (20)
where $`\delta 𝒥_\alpha (x)=e^2\nu _\alpha \delta D_\alpha (x)_xN_\alpha ^0(x)`$. The boundary conditions at the ferromagnetic reservoir on the left hand side are
$$\delta N_{}(\epsilon ,L)=\delta N_{}(\epsilon ,L)=0.$$
(21)
whereas the boundary conditions at the superconducting reservoir on the right hand side are
$`\sigma _{}_x\delta N_{}|_0\sigma _{}_x\delta N_{}|_0`$ $`=`$ $`\delta 𝒥_{}\delta 𝒥_{},`$ (22)
$`\delta N_{}(\epsilon ,0)+\delta N_{}(\epsilon ,0)`$ $`=`$ $`0.`$ (23)
General solutions to the differential equations, Eqs. (19,20), are
$$\sigma _{}\delta N_{}+\sigma _{}\delta N_{}=Ux+V^{y=x}\left[\delta 𝒥_{}(y)+\delta 𝒥_{}(y)\right]𝑑y,$$
(24)
and
$`\delta N_{}\delta N_{}`$ $`=`$ $`We^{x/L_s}+Ye^{+x/L_s}`$ (26)
$`{\displaystyle ^{y=x}}\left[{\displaystyle \frac{\delta 𝒥_{}(y)}{\sigma _{}}}{\displaystyle \frac{\delta 𝒥_{}(y)}{\sigma _{}}}\right]\mathrm{cosh}\left({\displaystyle \frac{xy}{L_s}}\right)𝑑y.`$
where $`U`$,$`V`$,$`W`$, and $`Y`$ are unknown factors to be determined by the four boundary conditions. It is thus possible to evaluate $`\delta N_\alpha `$ in terms of the correction to the diffusion coefficient, $`\delta D_\alpha `$, and the classical part of the distribution, $`N_\alpha ^0`$.
Taking into account both $`\delta D_\alpha `$ and $`\delta N_\alpha `$, the weak localization correction to the current density summed over both spin channels is
$$\delta j=_0^{\mathrm{}}𝑑\epsilon \left[\delta 𝒥_{}+\delta 𝒥_{}+\sigma _{}_x\delta N_{}+\sigma _{}_x\delta N_{}\right].$$
(28)
This is determined by the factor $`U`$ in the general solution Eq. (24). Using the boundary conditions, Eqs. (21,22,23), to evaluate $`U`$, the weak localization correction to the current density may be expressed as
$`\delta j`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\epsilon }{\mathrm{\Gamma }}}{\displaystyle _L^0}{\displaystyle \frac{dx}{L}}\{[\delta 𝒥_{}+\delta 𝒥_{}]+`$ (30)
$`{\displaystyle \frac{(\sigma _{}\sigma _{})}{2}}[{\displaystyle \frac{\delta 𝒥_{}}{\sigma _{}}}{\displaystyle \frac{\delta 𝒥_{}}{\sigma _{}}}]{\displaystyle \frac{\mathrm{cosh}\left[(x+L)/L_s\right]}{\mathrm{cosh}\left[L/L_s\right]}}\}.`$
The term in brackets contains spatial dependences arising from the $`\mathrm{cosh}`$ term and from $`\delta 𝒥_\alpha (x)=e^2\nu _\alpha \delta D_\alpha (x)_xN_\alpha ^0(x)`$, where $`_xN_\alpha ^0(x)`$ depends on $`\mathrm{cosh}[(x+L)/L_s]`$, Eq.(10), and $`\delta D_\alpha (x)`$ is given in Appendix A. Taking these terms into account, it is straightforward to perform the integration with respect to the coordinate $`x`$. The final step is to evaluate $`\delta D_\alpha (x)`$ numerically as detailed in Appendix A. For a long wire, we express the weak localization correction as a sum of resistances,
$$\delta R=\delta R^\mathrm{B}+\delta r_c^{FS}$$
(32)
where $`\delta R^\mathrm{B}`$ is the bulk contribution, $`\delta R^\mathrm{B}=(8e^2/\pi )R_{\mathrm{}}^2L(L_{}+L_{})/L_{}^2`$ and $`\delta r_c^{FS}`$ is a contact term,
$$\delta r_c^{FS}(1\gamma )\left(\frac{8e^2}{\pi }\right)\frac{\left(\sigma _{}+\sigma _{}\right)^2}{\sigma _{}\sigma _{}}\left(R_{\mathrm{}}\frac{L_s}{L_{}}\right)^2,$$
(33)
The parameter $`\gamma `$ is plotted as a function of the degree of polarization, $`\zeta `$, in Fig. 3. The solid line is $`\gamma `$ for $`L/L_s=20`$ whereas the dashed line is the contribution of the first two terms in Eq. (28) only, neglecting the influence of density variations. It is worth noting that the analytic result for an F/N system is $`\gamma =0`$ and for a polydomain wire connected to an S reservoir it is $`\gamma =1/2`$. In these cases, $`\delta N=0`$. For moderate values of $`\zeta `$, we find that there is no large change in $`\gamma `$, $`\gamma 1/2`$, and the influence of density variations is small. However, for very large spin polarization, $`\zeta 0.8`$, $`\gamma `$ increases dramatically and the role of density variations is vital. Drawing on the interpretation of the results for the return probability mentioned in Appendix A, we naively estimate the resistance of a system with $`\zeta =1`$ by considering the majority (up) spin channel to behave as if connected to an insulating barrier and the minority channel to behave as if connected to a normal barrier. Such a procedure gives $`\gamma 1\sigma _{}/(\sigma _{}+\sigma _{})1`$. This rough estimation of the value of $`\gamma `$ at $`\zeta =1`$ appears to be in agreement with Fig. 3. However, some notes of caution about the results for large polarization need to made. Firstly, the graphs are shown as a function of $`\zeta `$ for fixed $`L/L_s=20`$, whereas the spin relaxation lengths, $`L_\alpha `$, in the individual spin channels depend on $`\zeta `$. In particular, $`L_{}L_s`$ for large $`\zeta `$. Secondly, the error in the numerical evaluation of the diffusion coefficient increases rapidly as $`\zeta 1`$. Finally, values of $`\zeta 0.8`$ are unlikely to be realised in real materials and, even if they were, the intrinsic magnetic field in such materials would destroy any quantum interference.
## IV Conclusion
The weak localization correction to the classical diffusion coefficient, $`\delta D`$, is dependent on polarization, with majority spins more likely to be reflected from the F/S interface than minority spins. Taking into account the change in the spin polarized particle distribution in the F wire arising from $`\delta D`$, we found that the weak localization correction to the contact resistance is related to the square of the resistance of a piece of F wire with length equal to the spin relaxation length (see Eq. (1)) with a numerical prefactor that is almost constant for all experimentally relevant values of spin polarization.
###### Acknowledgements.
The authors are grateful to D. E. Khmelnitskii, B. Pannetier, V. T. Petrashov, and I. A. Sosnin for useful discussions. This work was supported by EPSRC.
## A Correction to the diffusion coefficient
This appendix describes the evaluation of the weak localization correction to the classical diffusion coefficient, $`\delta D_\alpha `$, Eq. (14). In order to calculate the Cooperon, Eq. (15), we consider spinor eigenvectors with components $`\psi _{n\alpha }(x)`$ that obey the following diffusion equations:
$`\left(\nu _{}D_{}_x^2+\nu _{}\tau _{}^1\right)\psi _n(x)`$ $`=`$ $`\lambda _n\psi _n(x),`$ (A1)
$`\left(\nu _{}D_{}_x^2+\nu _{}\tau _{}^1\right)\psi _n(x)`$ $`=`$ $`\lambda _n\psi _n(x),`$ (A2)
such that the Cooperon $`C_{\alpha \alpha }`$ is given by
$$C_{\alpha \alpha }(x,x)=\underset{n0}{\overset{\mathrm{}}{}}\frac{\left|\psi _{n\alpha }(x)\right|^2}{\lambda _n}.$$
(A3)
The momenta in the spin channels, $`Q_{n\alpha }`$, are related, $`\nu _{}D_{}Q_n^2+\nu _{}\tau _{}^1=\nu _{}D_{}Q_n^2+\nu _{}\tau _{}^1\lambda _n.`$ At the left hand reservoir ($`x=L`$) the Cooperon is zero since a particle that escapes into it suffers dephasing, Eq. (16),
$$\psi _n(L)=\psi _n(L)=0.$$
(A4)
whereas when the reservoir on the right hand ($`x=0`$) is a superconductor, the boundary conditions Eqs. (17,18) give
$`\psi _n(0)`$ $`=`$ $`\mathrm{e}^{i\chi }\psi _n(0),`$ (A5)
$`\sigma _{}_x\psi _n(x)|_{x=0}`$ $`=`$ $`\sigma _{}\mathrm{e}^{i\chi }_x\psi _n(x)|_{x=0}.`$ (A6)
The spinor eigenvectors also obey a normalisation condition, $`_L^0(|\psi _n(x)|^2+|\psi _n(x)|^2)𝑑x=1`$. We find the general solutions to the diffusion equations and use the boundary conditions at the left hand lead ($`x=L`$), Eq. (A4), to give
$$\psi _{n\alpha }(x)=𝒩_{n\alpha }\mathrm{sin}\left[Q_{n\alpha }(x+L)\right]$$
(A7)
On substituting into the boundary conditions on the right hand side ($`x=0`$), we find
$`𝒩_n\mathrm{sin}\left[Q_nL\right]`$ $`=`$ $`\mathrm{e}^{i\chi }𝒩_n\mathrm{sin}\left[Q_nL\right],`$ (A8)
$`\sigma _{}Q_n𝒩_n\mathrm{cos}\left[Q_nL\right]`$ $`=`$ $`\mathrm{e}^{i\chi }\sigma _{}Q_n𝒩_n\mathrm{cos}\left[Q_nL\right].`$ (A9)
Eliminating the normalisation constants leaves an equation for determining the eigenvalues,
$`\sigma _{}Q_n`$ $`\mathrm{cos}\left[Q_nL\right]\mathrm{sin}\left[Q_nL\right]+`$ (A11)
$`\sigma _{}Q_n\mathrm{cos}\left[Q_nL\right]\mathrm{sin}\left[Q_nL\right]=0.`$
We solve this equation numerically to find the eigenvalues and determine their contribution to the Cooperon, Eq. (A3).
As an example we present results for the calculation of the spatially averaged return probability, $`\overline{\delta D}_\alpha =_L^0\delta D_\alpha (x)(dx/L),`$ for arbitrary spin polarization in the ferromagnet. For long wires, $`LL_S`$, we express it as
$$\overline{\delta D}_\alpha \delta D_\alpha ^\mathrm{B}\left(1\eta _\alpha \frac{L_\alpha }{L}\right),$$
(A12)
where $`\delta D_\alpha ^\mathrm{B}`$ is the bulk contribution, $`\delta D_\alpha ^\mathrm{B}=8L_\alpha /\pi \nu _\alpha L_{}^2`$. The numerical coefficient $`\eta _\alpha `$ depends on the state of the reservoir and the degree of polarization in the ferromagnet. We compare with analytic results obtained in certain limits . When the reservoir on the right hand side of the wire is in the normal state, the boundary conditions are the same as those in the left hand reservoir and $`\eta _\alpha =1`$. For a poly-domain ferromagnetic wire connected to an S reservoir, where the classical contact resistance Eq. (13) gives no effect, $`\eta _\alpha =1/2`$. The second term in the above equation represents Cooperon decay due to the probability of particle escape into the right hand reservoir. Multiple Andreev reflection at the boundary, illustrated in a sketch in Fig. 2b, causes a factor of two reduction in the case of a polydomain ferromagnetic wire connected to a superconducting reservoir. In the case of a ferromagnetic wire connected to an insulating reservoir, $`\eta _\alpha =0`$, since the probability of particle escape into the right hand reservoir is totally supressed.
In Fig. 4 we plot $`\eta _\alpha `$ for a ferromagnetic wire with arbitrary polarization connected to a superconducting reservoir. It is shown for each spin channel separately and we assume that the up-spins are the majority spins, $`\sigma _{}\sigma _{}`$. For zero spin polarization, $`\eta _{}=\eta _{}=1/2`$, as described above. On increasing the spin polarization, $`\eta _{}<1/2`$ whereas $`\eta _{}>1/2`$. Majority carriers at the superconducting interface cannot find minority carrier states to Andreev reflect into and are normally reflected which increases the probability to return whereas minority carriers at the interface have an enhanced probability to escape. In the limit of total spin polarization, $`\eta _{}0`$ and $`\eta _{}1`$ which means that the reservoir appears to be totally insulating to majority spins but normal as far as minority spins are concerned.
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# NEW DIRECTIONS FOR NEW DIMENSIONS: FROM STRINGS TO NEUTRINOS TO AXIONS TO… aafootnote a Invited plenary talk given at PASCOS ’99: 7th International Symposium on Particles, Strings, and Cosmology (held at Lake Tahoe, California, 10–16 December 1999). To appear in the Proceedings.
## 1 Introduction: Lowering the fundamental scales of physics
The possibility of large extra spacetime dimensions has recently received considerable attention. This is clearly an exciting prospect. One of the earliest proponents of TeV-scale extra dimensions was Antoniadis $`^\mathrm{?}`$, who attempted to use such extra dimensions to explain supersymmetry breaking. Later, Witten $`^\mathrm{?}`$ pointed out that extra large dimensions could lower the string scale below its usual value near 10<sup>18</sup> GeV, and subsequently Lykken $`^\mathrm{?}`$ proposed that Witten’s idea could be extended to lower the string scale all the way to the TeV-range. Finally, in March 1998, it was proposed that extra dimensions could also be used to lower the fundamental Planck scale $`^\mathrm{?}`$ as well as the fundamental GUT scale $`^\mathrm{?}`$. Thus, combining these different proposals, it becomes possible to contemplate a self-consistent scenario in which all high fundamental energy scales (GUT, Planck, and string scales) are eliminated in favor of large extra spacetime dimensions!
It is important to distinguish two different types of extra spacetime dimensions. First, there are so-called “universal” extra dimensions. These extra dimensions are experienced by all forces, both gauge and gravitational; in technical terminology, these extra dimensions are “in the brane”. Because they affect the gauge forces (as probed by accelerator experiments), such dimensions can be no larger than roughly an inverse TeV. By contrast, the second class of extra dimensions are felt only by gravity; they are perpendicular to the D-brane on which the gauge forces are localized, and may therefore be considered “off the brane”. The sizes of such extra dimensions are significantly less constrained, and may in fact be as large as a millimeter.
Both of these types of extra dimensions play a role in lowering the fundamental scales of physics. Indeed, as outlined above, there are three different proposals: extra dimensions to lower the GUT scale $`^\mathrm{?}`$, extra dimensions to lower the Planck scale $`^\mathrm{?}`$, and extra dimensions to lower the string scale $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$. We shall now briefly review these three proposals.
In the proposal of Ref. $`^\mathrm{?}`$ to lower the GUT scale, one introduces some number $`\delta `$ of “universal” extra spacetime dimensions “in the brane” \[so that the Standard Model resides on a $`D(3+\delta )`$ brane\], and imagines that these dimensions have a common radius $`R`$. Because these extra dimensions are felt by the gauge forces, they change the running of the three gauge couplings from logarithmic to power-law behavior:
$$\alpha _i^1(\mu )=\alpha _i^1(\mu _0)\frac{b_i\stackrel{~}{b}_i}{2\pi }\mathrm{ln}(R\mu )\frac{\stackrel{~}{b}_iX_\delta }{2\pi \delta }\left[\left(R\mu \right)^\delta 1\right].$$
(1)
The emergence of power-law behavior is expected simply from dimensional analysis, since the gauge couplings themselves become dimensionful in higher dimensions, and hence have a classical scaling in addition to their quantum-mechanical (logarithmic) running. This power-law behavior can also be realized via a Kaluza-Klein summation, as discussed in Refs. $`^{\mathrm{?},\mathrm{?}}`$. In Eq. (1), $`X_\delta `$ is a normalization constant and ($`b_i`$, $`\stackrel{~}{b}_i`$) represent the one-loop beta-functions appropriate for the zero-mode and excited Kaluza-Klein states respectively. The exact values of these beta-functions depend on details of the compactification, as discussed in Ref. $`^\mathrm{?}`$. However, as shown in Ref. $`^\mathrm{?}`$, the remarkable feature of this higher-dimensional running is that gauge coupling unification is still generally preserved, but with a lowered unification scale! As an interesting case, let us consider $`R^1=1`$ TeV and $`\delta =1`$. With one-loop beta-function coefficients $`(b_1,b_2,b_3)=(33/5,1,3)`$ and $`(\stackrel{~}{b}_1,\stackrel{~}{b}_2,\stackrel{~}{b}_3)=(3/5,3,6)`$, corresponding to a certain orbifold compactification discussed in Ref. $`^\mathrm{?}`$, we then find the unification shown in Fig. 1. An important point to notice is that no large hierarchy is needed between the scale of the extra dimensions and the lowered GUT scale.
This reduced-scale unification leads to many important quantitative questions. How predictive is this unification? How perturbative is it? How sensitive is it to unification-scale effects? What about higher-loop corrections? These issues are discussed in Ref. $`^\mathrm{?}`$. The upshot is that this sort of unification scenario is predictive, perturbative, and not unreasonably sensitive to unification-scale effects. There have also been many further extensions of these basic ideas $`^\mathrm{?}`$. These include the study of two- and higher-loop effects; the incorporation of extra matter beyond the MSSM in order to increase the numerical accuracy of the unification; alternative derivations of these RGE’s from a Wilsonian perspective; studies of regularization independence; the extension of these ideas to the power-law running of Yukawa couplings; the higher-dimensional evolution of soft supersymmetry-breaking mass parameters; multi-step higher-dimensional unification scenarios; and alternative embeddings of the Standard Model into higher dimensions. Alternative ideas pertaining to reduced-scale gauge unification have also been discussed in Refs. $`^\mathrm{?}`$.
Extra dimensions can also be used to lower the Planck scale, as pointed out in Ref. $`^\mathrm{?}`$. Indeed, in many respects this Planck-scale proposal and the above GUT-scale proposal are the gravitational/gauge counterparts of each other. Whereas the GUT proposal utilizes $`\delta `$ extra dimensions “in the brane” with radius $`R`$ to modify the running of the three gauge couplings, the Planck proposal of Ref. $`^\mathrm{?}`$ utilizes some number $`n`$ of extra dimensions of radius $`r`$ “off the brane” to modify the running of the effective dimensionless gravitational (Newton) coupling $`\stackrel{~}{G}_N(\mu )\mu ^2G_N`$. As expected, the presence of the extra dimensions enhances the power-law running of this gravitational coupling, changing the scaling behavior from $`\mu ^2`$ to $`\mu ^{2+n}`$. This in turn lowers the Planck scale \[i.e., the fundamental gravitational scale, defined as the scale where $`\stackrel{~}{G}_N(\mu )𝒪(1)`$\]. Unlike the GUT proposal, however, one typically requires a significant hierarchy between the scale of the extra dimensions and the lowered Planck scale. For example, in the case $`n=2`$ with a lowered Planck scale in the TeV-range, one finds $`r`$ millimeter $`(10^4`$ eV$`)^1`$. This hierarchy is exactly as large as the original hierarchy between the electroweak scale and the usual four-dimensional Planck scale.<sup>a</sup><sup>a</sup>a As pointed out in Ref. $`^\mathrm{?}`$, it may be possible to avoid the former hierarchy and nevertheless explain the latter hierarchy by virtue of a “warp” rescaling factor. Issues surrounding gauge coupling unification in this scenario are discussed in Ref. $`^\mathrm{?}`$.
Finally, combinations of both types of extra dimensions can be used to lower the string scale $`^{\mathrm{?},\mathrm{?}}`$. For Type I strings, the string scale ultimately depends on the six-volume $`V_6`$ of compactification from ten flat dimensions to four flat dimensions:
$$M_{\mathrm{string}}\sqrt{\frac{1}{\alpha _{\mathrm{GUT}}M_{\mathrm{Planck}}}}V_6^{1/4}.$$
(2)
Therefore, as discussed in Ref. $`^\mathrm{?}`$, if we seek to combine the above GUT and Planck scenarios together within string theory, we can write $`V_6=R^\delta r^{6\delta }`$ where $`(R,\delta )`$ describe the extra dimensions “in” the brane (to produce a lowered GUT scale) and $`(r,6\delta )`$ describe the extra dimensions “off” the brane (to produce a lowered Planck scale). If we demand that the lowered GUT scale coincide with $`M_{\mathrm{string}}`$, one can then solve to obtain a self-consistent solution. For example, let us consider the case with $`\delta =1`$ and $`R^10.5`$ TeV. (This extreme value may already be ruled out experimentally, but will serve our illustrative purposes.) This implies that $`M_{\mathrm{GUT}}^{}10`$ TeV, which in turn implies (after a T-duality transformation) that the remaining five extra dimensions must have radius $`r(10\mathrm{MeV})^1`$.
Thus, putting the pieces together in this example, we are led to a unified embedding into string theory, as illustrated in Fig. 2. Above the string scale $`M_{\mathrm{string}}=10`$ TeV, the physics is described in terms of a full Type I string theory. Below 10 TeV, by contrast, the physics is described by a series of effective field theories in which the gauge and gravitational forces feel different numbers of spacetime dimensions. Together, everything is balanced so as to produce a self-consistent simultaneous lowering of GUT, Planck, and string scales. Of course, other configurations are also possible.
## 2 Light neutrino masses without heavy mass scales: <br>A higher-dimensional seesaw mechanism
As we have seen, the lesson from the above developments has been that heavy mass scales in four dimensions can be replaced by lighter mass scales in higher dimensions. However, low-energy neutrino data seem to provide independent evidence for yet another heavy mass scale, namely the seesaw scale. The seesaw mechanism relies on the existence of a new heavy mass scale $`MM_{\mathrm{GUT}}`$ associated with a right-handed neutrino singlet field $`N`$. The question then emerges whether it is possible to generate light neutrino masses without the introduction of a heavy mass scale, perhaps by some intrinsically higher-dimensional mechanism.
To date, there have been essentially two ideas concerning how this might be accomplished within the large extra-dimension framework: one proposal $`^\mathrm{?}`$ utilizes a higher-dimensional seesaw mechanism, and the other $`^\mathrm{?}`$ utilizes a higher-dimensional volume factor. Both proposals originate with the same observation: because the right-handed neutrino is a Standard-Model gauge singlet, it need not be restricted to a “brane” with respect to the full higher-dimensional space. It is therefore possible for this field to experience extra spacetime dimensions and thereby accrue an infinite tower of Kaluza-Klein excitations. This then leads to a number of higher-dimensional mechanisms $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ for suppressing the resulting neutrino mass without a heavy mass scale. In the following, we shall concentrate on one of the mechanisms advanced in Ref. $`^\mathrm{?}`$, namely the intriguing possibility that in higher dimensions, neutrino oscillations need not imply the existence of neutrino masses at all! This would then eliminate the need for a high fundamental scale. Other mechanisms for explaining light but non-zero neutrino masses are also discussed in Refs. $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$.
We begin by assuming that the right-handed neutrino feels extra dimensions, while the left-handed neutrino $`\nu _L`$ does not. For concreteness, we consider a Dirac fermion $`\mathrm{\Psi }`$ in five dimensions, and work in the Weyl basis in which $`\mathrm{\Psi }`$ can be decomposed into two two-component spinors: $`\mathrm{\Psi }=(\psi _1,\overline{\psi }_2)^T`$. When the extra spacetime dimension is compactified on a $`\text{ZZ}_2`$ orbifold, it is natural for one of the two-component Weyl spinors, e.g., $`\psi _1`$, to be taken to be even under the $`\text{ZZ}_2`$ action $`yy`$, while the other spinor $`\psi _2`$ is taken to be odd. If the left-handed neutrino $`\nu _L`$ is restricted to a brane located at the orbifold fixed point $`y=0`$, then $`\psi _2`$ vanishes at this point and so the most natural coupling is between $`\nu _L`$ and $`\psi _1`$. For generality, we will also include a possible “bare” Majorana mass term for $`\mathrm{\Psi }`$ of the form $`\frac{1}{2}M_0\overline{\mathrm{\Psi }}^\mathrm{c}\mathrm{\Psi }`$. This then results in a Lagrangian of the form
$``$ $`=`$ $`{\displaystyle }d^4xdyM_s\{\overline{\psi }_1i\overline{\sigma }^\mu _\mu \psi _1+\overline{\psi }_2i\overline{\sigma }^\mu _\mu \psi _2+\frac{1}{2}M_0(\psi _1\psi _1+\psi _2\psi _2+\mathrm{h}.\mathrm{c}.)\}`$ (3)
$`+{\displaystyle }d^4x\{\overline{\nu }_Li\overline{\sigma }^\mu D_\mu \nu _L+(\widehat{m}\nu _L\psi _1|_{y=0}+\mathrm{h}.\mathrm{c}.)\}`$
where $`y`$ is the coordinate of the extra compactified spacetime dimension and where $`M_s`$ is the mass scale of the higher-dimensional fundamental theory (e.g., a reduced Type I string scale). Note that the last term represents the Dirac brane/bulk Yukawa coupling between $`\nu _L`$ and $`\psi _1`$.
Next, we compactify the Lagrangian (3) down to four dimensions by expanding the five-dimensional $`\mathrm{\Psi }`$ field in Kaluza-Klein modes. Imposing the orbifold relations $`\psi _{1,2}(y)=\pm \psi _{1,2}(y)`$ implies that our Kaluza-Klein decomposition takes the form $`\psi _1(x,y)=(2\pi R)^{1/2}_{n=0}^{\mathrm{}}\psi _1^{(n)}(x)\mathrm{cos}(ny/R)`$ and a similar result for $`\psi _2`$ with cosine replaced by sine. For convenience, we shall also define the linear combinations $`N^{(n)}(\psi _1^{(n)}+\psi _2^{(n)})/\sqrt{2}`$ and $`M^{(n)}(\psi _1^{(n)}\psi _2^{(n)})/\sqrt{2}`$ for all $`n>0`$. Inserting this decomposition into Eq. (3) and integrating over the compactified dimension, we then obtain an effective four-dimensional Lagrangian in which the Standard-Model neutrino $`\nu _L`$ mixes with the entire tower of Kaluza-Klein states of the higher-dimensional $`\mathrm{\Psi }`$ field with a mass mixing matrix of the form
$$=\left(\begin{array}{ccccccc}0& m& m& m& m& m& \mathrm{}\\ m& M_0& 0& 0& 0& 0& \mathrm{}\\ m& 0& M_0+1/R& 0& 0& 0& \mathrm{}\\ m& 0& 0& M_01/R& 0& 0& \mathrm{}\\ m& 0& 0& 0& M_0+2/R& 0& \mathrm{}\\ m& 0& 0& 0& 0& M_02/R& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(4)
In Eq. (4), we have defined the basis $`(\nu _L,\psi _1^{(0)},N^{(1)},M^{(1)},N^{(2)},M^{(2)},\mathrm{})`$. Note that $`m\widehat{m}/\sqrt{2\pi RM_s}`$ is the Dirac coupling suppressed by a volume factor corresponding to the extra spacetime dimension.
While in principle any value for $`M_0`$ is allowed (depending on the structure of the full effective Lagrangian derived from the particular underlying string model), the topological constraints associated with Scherk-Schwarz compactification naturally suggest two specific values: $`M_0=0`$ (corresponding to no breaking of lepton number), and $`M_0=(2R)^1`$ (corresponding to a global breaking of lepton number). Let us here consider the non-trivial possibility $`M_0=(2R)^1`$. Note that this value is fixed topologically, and hence does not require any fine-tuning. It is then possible to solve for the eigenvalues and eigenstates of the mass mixing matrix in Eq. (4). Remarkably, it turns out that for any value of $`mR`$, there exists an exactly zero eigenvalue, with a corresponding mass eigenstate given exactly by
$$|\stackrel{~}{\nu }_L=\frac{1}{\sqrt{1+\pi ^2m^2R^2}}\left\{|\nu _LmR\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k1/2}\left[|N^{(k1)}|M^{(k)}\right]\right\}$$
(5)
where we have defined $`N^{(0)}\psi _1^{(0)}`$. Even though this result is exact for all $`mR`$, in most realistic scenarios (see Ref. $`^\mathrm{?}`$), we have $`mR1`$. Thus, we see that even though this neutrino mass eigenstate contains a small, non-trivial admixture of Kaluza-Klein states, the dominant component of our massless neutrino eigenstate remains the left-handed gauge-eigenstate neutrino $`\nu _L`$, as required phenomenologically. Nevertheless, this particular admixture of excited Kaluza-Klein states has rendered the neutrino eigenstate exactly massless! In other words, the effects of the infinite tower of Kaluza-Klein states for the $`\mathrm{\Psi }`$ field have driven the neutrino mass exactly to zero.
One might still worry that a vanishing neutrino mass is unacceptable because of the recent evidence for neutrino oscillations. However, even though the neutrino mass is vanishing in this scenario, there continue to exist oscillations because of the non-trivial mixings between the left-handed neutrino and the infinite tower of Kaluza-Klein states. Specifically, upon diagonalizing the mass matrix (4) and calculating the resulting oscillation probabilities in the usual way, we find $`^\mathrm{?}`$ the result shown in Fig. 3. This figure thus provides explicit verification that neutrino oscillations do indeed occur, even though the physical neutrino is exactly massless. Of course, we have been discussing only the simple case of neutrino/anti-neutrino oscillations. However, this mechanism can easily be generalized to include the case of flavor oscillations as well.
At first sight, it may seem strange that we are able to have neutrino oscillations without neutrino masses. However, the crucial point is that once the infinite towers of Kaluza-Klein states are included, the mixing mass matrix (4) with $`M_0=(2R)^1`$ yields a zero eigenvalue without becoming diagonal. This is not possible in the usual four-dimensional neutrino/anti-neutrino scenarios, where the analogous mixing matrix takes the simple form $`\left(\begin{array}{cc}0& m\\ m& M_0\end{array}\right)`$. Indeed, the masses of the right-handed Kaluza-Klein states themselves are sufficient to generate the desired oscillations indirectly, even though these Kaluza-Klein states are in the “bulk” rather than on the brane. Thus, if such a scenario can be realized within the context of a fully realistic string model, then the recent observations of neutrino oscillations can be re-interpreted not as providing evidence for neutrino masses, but rather as providing evidence for extra spacetime dimensions!
## 3 Extra dimensions and “invisible” axions
Many of the above ideas are completely general, and apply to other bulk fields as well. Towards this end, let us now discuss how extra spacetime dimensions may contribute to the invisibility of the QCD axion. Like the graviton and right-handed neutrino, the QCD axion is also a Standard-Model singlet. The QCD axion is therefore free to propagate into the bulk.
Can this be used to lower the fundamental Peccei-Quinn (PQ) symmetry-breaking scale? This issue has been investigated in Refs. $`^{\mathrm{?},\mathrm{?}}`$. As explicitly shown in Ref. $`^\mathrm{?}`$ (and first proposed in Ref. $`^\mathrm{?}`$), it is indeed possible to exploit the volume factor of large extra dimensions in order to realize a large effective four-dimensional PQ scale from a smaller, higher-dimensional fundamental PQ scale. Thus, once again, no large fundamental energy scales are required.
However, as discussed in Ref. $`^\mathrm{?}`$, the presence of Kaluza-Klein axions can have important and unexpected effects on axion phenomenology. Just as in the neutrino case discussed above, the naïve four-dimensional axion mixes with the infinite tower of Kaluza-Klein axions, with a mass matrix given in Ref. $`^\mathrm{?}`$. This mixing has a number of interesting phenomenological consequences. First, as shown in Ref. $`^\mathrm{?}`$, under certain circumstances the mass of the axion essentially decouples from the PQ scale, and instead is set by the radius of the extra spacetime dimension! Thus, axions in the $`10^4`$ eV mass range are consistent with (sub-)millimeter extra dimensions. This decoupling implies that it may be possible to adjust the mass of the axion independently of its couplings to matter. This is not possible in four dimensions.
Second, as discussed in Ref. $`^\mathrm{?}`$, the usual four-dimensional axion is no longer a mass eigenstate because of the non-trivial axion mass mixing matrix. This implies that the four-dimensional axion should undergo laboratory oscillations which are entirely analogous to neutrino oscillations. Moreover, because the axion is now a bulk field, Standard-Model particles couple not only to the axion zero-mode, but rather to the entire linear superposition $`a^{}_na_n`$ (where $`a_n`$ are the axion Kaluza-Klein modes). Therefore, the quantity of phenomenological interest is the probability $`P_{a^{}a^{}}(t)`$ that $`a^{}`$ is preserved as a function of time. This probability $`^\mathrm{?}`$ is shown in Fig. 4. Unlike the neutrino case, we see that the probability drops rapidly from $`1`$ (at the initial time $`t=0`$) to extremely small values (expected to be $`10^{16}`$ when an appropriately truncated set of $`10^{16}`$ Kaluza-Klein states are included in $`a^{}`$). At no future time does this probability regenerate. Essentially, the axion state $`a^{}`$ has “decohered” and becomes invisible with respect to subsequent laboratory interactions. This decoherence is therefore a possible mechanism contributing to an invisible axion.
Finally, one can investigate the effects of Kaluza-Klein axions on cosmological relic axion oscillations. In this regard it is important to understand whether the coupled axion Kaluza-Klein states accelerate or retard the dissipation of the cosmological energy density associated with these oscillations. Remarkably, one finds $`^\mathrm{?}`$ that the net effect of these coupled Kaluza-Klein axions is to either preserve or enhance the rate of energy dissipation. This implies that the usual relic oscillation bounds are loosened in higher dimensions, which suggests that it may be possible to raise the effective PQ symmetry-breaking scale beyond its usual four-dimensional value. This could therefore potentially serve as another factor contributing to axion invisibility. Together, these results suggest that it may be possible to develop a new, higher-dimensional approach to axion phenomenology.
## 4 Conclusions
Only experiment will decide if large extra spacetime dimensions actually exist, and if the fundamental high-energy scales of physics are really as low as the TeV-range. Nevertheless, what is remarkable about the recent developments is that they illustrate that the fundamental energy scales are not immutable, and that the parameter space for physics beyond the Standard Model is significantly broader than had been previously thought. Moreover, it is equally remarkable and gratifying that ideas originally born in string theory are having such a profound effect on the answers to primarily phenomenological questions, and that these ideas may be potentially testable in the not-too-distant future. If nothing else, these may be the most valuable lessons that we may take with us from the brane world.
## Acknowledgments
I wish to thank the organizers of the PASCOS ’99 conference, and especially Jack Gunion, for the opportunity to speak at such a stimulating and wide-ranging conference.
## References
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# Warped Compactifications and AdS/CFT11footnote 1Talk presented at the TMR conference in Paris, September 99.
## 1 Introduction
The consideration of the near horizon geometry of branes on one hand, and the low energy dynamics on their worldvolume on the other hand has lead to conjectured duality relations between field theories and string theory (M theory) on certain backgrounds . The field theories under discussion are in various dimensions, can be conformal or not, and with or without supersymmetry. These properties are reflected by the type of string/M theory backgrounds of the dual description.
In this talk we discuss two classes of examples of warped products of AdS spaces in the context of the AdS/CFT correspondence. The warped product structure is a fibration of AdS over some manifold $``$ . It is the most general form of a metric that has the isometry of an AdS space .
The first class of examples appears in the construction of dual Type I’ string descriptions to five dimensional supersymmetric fixed points with $`E_{N_f+1}`$ global symmetry, $`E_{N_f+1}=(E_8,E_7,E_6,E_5=Spin(10),E_4=SU(5),E_3=SU(3)\times SU(2),E_2=SU(2)\times U(1),E_1=SU(2))`$. These fixed points are obtained in the limit of infinite bare coupling of $`𝒩=2`$ supersymmetric gauge theories with gauge group $`Sp(Q_4)`$, $`N_f<8`$ massless hypermultiplets in the fundamental representation and one massless hypermultiplet in the anti-symmetric representation . These theories were discussed in the context of the AdS/SCFT correspondence in . The dual background is obtained as the near horizon geometry of the D4-D8 brane system in massive Type IIA supergravity. The ten dimensional space is a fibration of $`AdS_6`$ over $`S^4`$ and has the isometry group $`SO(2,5)\times SO(4)`$. This space provides the spontaneous compactification of massive Type IIA supergravity in ten dimensions to the $`F(4)`$ gauged supergravity in six dimensions .
The second class of examples appears when considering the $`𝒩=2`$ superconformal theories defined on a $`3+1`$ dimensional hyperplane intersection of two sets of M5 branes. The $`𝒩=2`$ supersymmetry algebra in four dimensions contains a central extension term that corresponds to string charges in the adjoint representation of the $`SU(2)_R`$ part of the R-symmetry group. This implies that $`𝒩=2`$ supersymmetric gauge theories in four dimensions can have BPS string configurations at certain regions in their moduli space of vacua. In particular, at certain points in the moduli space of vacua these strings can become tensionless. A brane configuration that exhibits this phenomena consists of two sets of M5 branes intersecting on $`3+1`$ dimensional hyperplane. The theory on the intersection is $`𝒩=2`$ supersymmetric. One can stretch M2 branes between the two sets of M5 branes in a configuration that preserves half of the supersymmetry. This can be viewed as a BPS string of the four dimensional theory, and we will study such brane configurations using the AdS/CFT correspondence .
The talk is organized as follows. In the next two sections we will discuss the first class of examples. In section 2 we will discuss the D4-D8 brane system and its relation to the five dimensional fixed points. In section 3 we will construct the dual string description and use it to deduce some properties of the fixed points. In the following two sections we will discuss the second class of examples. In section 4 we will construct the dual string description of the four dimensional theory on the intersection and discuss the field theory on the intersection. In section 5 we will use the dual string description to deduce some properties of these strings. We will argue that from the four dimensional field theory viewpoint they are simply BPS string configurations on the Higgs branch.
## 2 The D4-D8 Brane System
We start with Type I string theory on $`R^9\times S^1`$ with $`N`$ coinciding D5 branes wrapping the circle. The six dimensional D5 brane worldvolume theory possesses $`𝒩=1`$ supersymmetry. It has an $`Sp(N)`$ gauge group, one hypermultiplet in the antisymmetric representation of $`Sp(N)`$ from the DD sector and 16 hypermultiplets in the fundamental representation from the DN sector. Performing T-duality on the circle results in Type I’ theory compactified on the interval $`S^1/Z_2`$ with two orientifolds (O8 planes) located at the fixed points. The D5 branes become D4 branes and there are 16 D8 branes located at points on the interval. They cancel the -16 units of D8 brane charge carried by the two O8 planes. The locations of the D8 branes correspond to masses for the hypermultiplets in the fundamental representation arising from the open strings between the D4 branes and the D8 branes. The hypermultiplet in the antisymmetric representation is massless.
Consider first $`N=1`$, namely one D4 brane. The worldvolume gauge group is $`Sp(1)SU(2)`$. The five-dimensional vector multiplet contains as bosonic fields the gauge field and one real scalar. The scalar parametrizes the location of the D4 branes in the interval, and the gauge group is broken to $`U(1)`$ unless the D4 brane is located at one of the fixed points. A hypermultiplet contains four real (two complex) scalars. The $`N_f`$ massless matter hypermultiplets in the fundamental and the antisymmetric hypermultiplet (which is a trivial representation for $`Sp(1)`$) parametrize the Higgs branch of the theory. It is the moduli space of $`SO(2N_f)`$ one-instanton. The theory has an $`SU(2)_R`$ R-symmetry. The two supercharges as well as the scalars in the hypermultiplet transform as a doublet under $`SU(2)_R`$. In addition, the theory has a global $`SU(2)\times SO(2N_f)\times U(1)_I`$ symmetry. The $`SU(2)`$ factor of the global symmetry group is associated with the massless antisymmetric hypermultiplet and is only present if $`N>1`$, the $`SO(2N_f)`$ group is associated with the $`N_f`$ massless hypermultiplets in the fundamental and the $`U(1)_I`$ part corresponds to the instanton number conservation.
Consider the D8 brane background metric. It takes the form
$`ds^2`$ $`=`$ $`H_8^{1/2}(dt^2+dx_1^2+\mathrm{}+dx_8^2)`$
$`+`$ $`H_8^{1/2}dz^2,`$
$`e^\varphi `$ $`=`$ $`H_8^{5/4}.`$ (1)
$`H_8`$ is a harmonic function on the interval parametrized by $`z`$. Therefore $`H_8`$ is a piecewise linear function in $`z`$ where the slope is constant between two D8 branes and decreases by one unit for each D8 brane crossed. Thus,
$$H_8(z)=c+16\frac{z}{l_s}\underset{i=1}{\overset{16}{}}\frac{|zz_i|}{l_s}\underset{i=1}{\overset{16}{}}\frac{|z+z_i|}{l_s},$$
(2)
where the $`z_i`$ denote the locations of the 16 D8 branes.
Denote the D4 brane worldvolume coordinates by $`t,x_1\mathrm{}x_4`$. The D4 brane is located at some point in $`x_5\mathrm{}x_8`$ and $`z`$. We can consider it to be a probe of the D8 branes background. The gauge coupling $`g`$ of the D4 brane worldvolume theory and the harmonic function $`H_8`$ are related as can be seen by expanding the DBI action of a D4 brane in the background (2). We get
$$g^2=\frac{l_s}{H_8},$$
(3)
where $`g_{cl}^2=\frac{c}{l_s}`$ corresponds to the classical gauge coupling. In the field theory limit we take $`l_s0`$ keeping the gauge coupling $`g`$ fixed, thus,
$$g^2=\text{fixed},\varphi =\frac{z}{l_s^2}=\text{fixed},l_s0.$$
(4)
In this limit we have
$$\frac{1}{g^2}=\frac{1}{g_{cl}^2}+16\varphi \underset{i=1}{\overset{16}{}}|\varphi m_i|\underset{i=1}{\overset{16}{}}|\varphi +m_i|,$$
(5)
where the masses $`m_i=\frac{z_i}{l_s^2}`$. Note that in the field theory limit we are studying the region near $`z=0`$. The coordinate $`\varphi `$ takes values in $`R^+`$ and parametrizes the field theory Coulomb branch.
Seiberg argued that the theory at the origin of the Coulomb branch obtained in the limit $`g_{cl}=\mathrm{}`$ with $`N_f<8`$ massless hypermultiplets is a non trivial fixed point. The restriction $`N_f<8`$ can be seen in the supergravity description as a requirement for the harmonic function $`H_8`$ in equation (2) to be positive when $`c=0`$ and $`z_i=0`$. At the fixed point the global symmetry is enhanced to $`SU(2)\times E_{N_f+1}`$. The Higgs branch is expected to become the moduli space of $`E_{N_f+1}`$ one-instanton.
The generalization to $`N=Q_4`$ D4 branes is straightforward . The gauge group is now $`Sp(Q_4)`$ and the global symmetry is as before. The Higgs branch is now the moduli space of $`SO(2N_f)`$ $`Q_4`$-instantons. At the fixed point the global symmetry is enhanced as before and the Higgs branch is expected to become the moduli space of $`E_{N_f+1}`$ $`Q_4`$-instanton. Our interest in the next section will be in finding dual string descriptions of these fixed points.
## 3 The Dual String Description
The low energy description of our system is given by Type I’ supergravity . The region between two D8 branes is, as discussed in , described by massive Type IIA supergravity . The configurations that we will study have $`N_f`$ D8 branes located at one O8 plane and $`16N_f`$ D8 branes at the other O8 plane. Therefore we are always between D8 branes and never encounter the situation where we cross D8 branes, and the massive Type IIA supergravity description is sufficient.
The bosonic part of the massive Type IIA action (in string frame) including a six-form gauge field strength which is the dual of the RR four-form field strength is
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _{10}^2}}{\displaystyle }d^{10}x\sqrt{g}(e^{2\mathrm{\Phi }}(R+4_\mu \mathrm{\Phi }^\mu \mathrm{\Phi })`$ (6)
$``$ $`{\displaystyle \frac{1}{26!}}|F_6|^2{\displaystyle \frac{1}{2}}m^2)`$
where the mass parameter is given by
$$m=\sqrt{2}(8N_f)\mu _8\kappa _{10}=\frac{8N_f}{2\pi l_s}.$$
(7)
We use the conventions $`\kappa _{10}=8\pi ^{7/2}l_s^4`$ and $`\mu _8=(2\pi )^{9/2}l_s^5`$ for the gravitational coupling and the D8 brane charge, respectively.
The Einstein equations derived from (6) read (in Einstein frame metric)
$`2R_{ij}`$ $`=`$ $`g_{ij}(R{\displaystyle \frac{1}{2}}|\mathrm{\Phi }|^2{\displaystyle \frac{e^{\mathrm{\Phi }/2}}{26!}}|F_6|^2{\displaystyle \frac{m^2}{2}}e^{5\mathrm{\Phi }/2})`$
$`+`$ $`_i\mathrm{\Phi }_j\mathrm{\Phi }+{\displaystyle \frac{e^{\mathrm{\Phi }/2}}{5!}}F_{im_1\mathrm{}m_5}F_j^{m_1\mathrm{}m_5},`$
$`0`$ $`=`$ $`^j_j\mathrm{\Phi }{\displaystyle \frac{5m^2}{4}}e^{5\mathrm{\Phi }/2}+{\displaystyle \frac{e^{\mathrm{\Phi }/2}}{46!}}|F_6|^2,`$
$`0`$ $`=`$ $`^i\left(e^{\mathrm{\Phi }/2}F_{im_1\mathrm{}m_5}\right).`$ (8)
Except for these formulas we will be using the string frame only.
For identifying the solutions of D4 branes localized on D8 branes it is more convenient to start with the conformally flat form of the D8 brane supergravity solution. It takes the form
$`ds^2`$ $`=`$ $`\mathrm{\Omega }(z)^2(dt^2+\mathrm{}+dx_4^2`$ (9)
$`+`$ $`d\stackrel{~}{r}^2+\stackrel{~}{r}^2d\mathrm{\Omega }_3^2+dz^2),`$
$`e^\mathrm{\Phi }`$ $`=`$ $`C\left({\displaystyle \frac{3}{2}}Cmz\right)^{\frac{5}{6}}\mathrm{\Omega }(z)=\left({\displaystyle \frac{3}{2}}Cmz\right)^{\frac{1}{6}}.`$
In these coordinates the harmonic function of $`Q_4`$ localized D4 branes in the near horizon limit derived from (3) reads
$$H_4=\frac{Q_4}{l_s^{10/3}(\stackrel{~}{r}^2+z^2)^{5/3}}.$$
(10)
This localized D4-D8 brane system solution, in a different coordinate system, has been constructed in . One way to determine the harmonic function (10) of the localized D4 branes is to solve the Laplace equation in the background of the D8 branes.
It is useful to make a change of coordinates $`z=r\mathrm{sin}\alpha ,\stackrel{~}{r}=r\mathrm{cos}\alpha ,0\alpha \pi /2`$. We get
$`ds^2`$ $`=`$ $`\mathrm{\Omega }^2(H_4^{\frac{1}{2}}(dt^2+\mathrm{}dx_4^2)`$
$`+`$ $`H_4^{\frac{1}{2}}(dr^2+r^2d\mathrm{\Omega }_4^2)),`$
$`e^\mathrm{\Phi }`$ $`=`$ $`C\left({\displaystyle \frac{3}{2}}Cmr\mathrm{sin}\alpha \right)^{\frac{5}{6}}H_4^{\frac{1}{4}},`$
$`\mathrm{\Omega }`$ $`=`$ $`\left({\displaystyle \frac{3}{2}}Cmr\mathrm{sin}\alpha \right)^{\frac{1}{6}},`$ (11)
$`F_{01234r}`$ $`=`$ $`{\displaystyle \frac{1}{C}}_r\left(H_4^1\right),`$
where $`C`$ is an arbitrary parameter of the solution and
$$d\mathrm{\Omega }_4^2=d\alpha ^2+(\mathrm{cos}\alpha )^2d\mathrm{\Omega }_3^2.$$
(12)
The background (3) is a solution of the massive Type IIA supergravity equations (3).
The metric (3) of the D4-D8 system can be simplified to
$`ds^2`$ $`=`$ $`\left({\displaystyle \frac{3}{2}}Cm\mathrm{sin}\alpha \right)^{\frac{1}{3}}(Q_4^{\frac{1}{2}}r^{\frac{4}{3}}dx_{}^2`$ (13)
$`+`$ $`Q_4^{\frac{1}{2}}{\displaystyle \frac{dr^2}{r^2}}+Q_4^{\frac{1}{2}}d\mathrm{\Omega }_4^2)`$
where $`dx_{}^2dt^2+\mathrm{}+dx_4^2`$. Define now the energy coordinate $`U`$ by $`r^2=l_s^5U^3`$. That this is the energy coordinate can be seen by calculating the energy of a fundamental string stretched in the $`r`$ direction or by using the DBI action as in the previous section. In the field theory limit, $`l_s0`$ with the energy $`U`$ fixed, we get the metric in the form of a warped product of $`AdS_6\times S^4`$
$`ds^2`$ $`=`$ $`l_s^2({\displaystyle \frac{3}{4\pi }}C(8N_f)\mathrm{sin}\alpha )^{\frac{1}{3}}(Q_4^{\frac{1}{2}}U^2dx_{}^2`$ (14)
$`+`$ $`Q_4^{\frac{1}{2}}{\displaystyle \frac{9dU^2}{4U^2}}+Q_4^{\frac{1}{2}}d\mathrm{\Omega }_4^2),`$
and the dilaton is given by
$$e^\mathrm{\Phi }=Q_4^{\frac{1}{4}}C\left(\frac{3}{4\pi }C(8N_f)\mathrm{sin}\alpha \right)^{\frac{5}{6}}.$$
(15)
The ten dimensional space described by (14) is a fibration of $`AdS_6`$ over $`S^4`$. It is the most general form of a metric that has the isometry of an $`AdS_6`$ space . The space has a boundary at $`\alpha =0`$ which corresponds to the location of the O8 plane ($`z=0`$). The boundary is of the form $`AdS_6\times S^3`$. In addition to the $`SO(2,5)`$ $`AdS_6`$ isometries, the ten dimensional space has also $`SO(4)`$ isometries associated with the spherical part of the metric (14). In general $`S^4`$ has the $`SO(5)`$ isometry group. However, this is reduced due to the warped product structure. As is easily seen from the form of the spherical part (12), only transformations excluding the $`\alpha `$ coordinate are isometries of (14). We are left with an $`SO(4)SU(2)\times SU(2)`$ isometry group.
The two different viewpoints of the D4-D8 brane system, the near horizon geometry of the brane system on one hand, and the low energy dynamics on the D4 branes worldvolume on the other hand suggest a duality relation. Namely, Type I’ string theory compactified on the background (14), (15) with a 4-form flux of $`Q_4`$ units on $`S^4`$ is dual to an $`𝒩=2`$ supersymmetric five dimensional fixed point. The fixed point is obtained in the limit of infinite coupling of $`Sp(Q_4)`$ gauge theory with $`N_f`$ hypermultiplets in the fundamental representation and one hypermultiplet in the antisymmetric representation, where $`m(8N_f)`$ as in (7). The $`SO(2,5)`$ symmetry of the compactification corresponds to the conformal symmetry of the field theory. The $`SU(2)\times SU(2)`$ symmetry of the compactification corresponds $`SU(2)_R`$ R-symmetry and to the $`SU(2)`$ global symmetry associated with the massless hypermultiplet in the antisymmetric representation.
At the boundary $`\alpha =0`$ the dilaton (15) blows up and Type I’ is strongly coupled. In the weakly coupled dual heterotic string description this is seen as an enhancement of the gauge symmetry to $`E_{N_f+1}`$. One can see this enhancement of the gauge symmetry in the Type I’ description by analysing the D0 brane dynamics near the orientifold plane . This means that we have $`E_{N_f+1}`$ vector fields that propagate on the $`AdS_6\times S^3`$ boundary, as in the Horava-Witten picture . The scalar curvature of the background (14), (15)
$$l_s^2(C(8N_f))^{\frac{1}{3}}Q_4^{\frac{1}{2}}(\mathrm{sin}\alpha )^{\frac{5}{3}},$$
(16)
blows up at the boundary as well. In the dual heterotic description the dilaton is small but the curvature is large, too. For large $`Q_4`$ there is a region, $`\mathrm{sin}\alpha Q_4^{\frac{3}{10}}`$, where both curvature (16) and dilaton (15) are small and thus we can trust supergravity.
The $`AdS_6`$ supergroup is $`F(4)`$. Its bosonic subgroup is $`SO(2,5)\times SU(2)`$. Romans constructed an $`𝒩=4`$ six dimensional gauged supergravity with gauge group $`SU(2)`$ that realizes $`F(4)`$ . It was conjectured in that it is related to a compactification of the ten dimensional massive Type IIA supergravity. Indeed, we find that the ten dimensional background space is the warped product of $`AdS_6`$ and $`S^4`$ (14) (with $`N_f=0`$). The reduction to six dimensions can be done in two steps. First we can integrate over the coordinate $`\alpha `$. This yields a nine dimensional space of the form $`AdS_6\times S^3`$. We can then reduce on $`S^3`$ to six dimensions, while gauging its isometry group. Roman’s construction is based on gauging an $`SU(2)`$ subgroup of the $`SO(4)`$ isometry group. Generally, it has been shown in that the compactification of ten dimensional massive Type IIA supergravity to six dimensional massive Type IIA supegravity on warped $`S^4`$, is a consistent non-linear Kaluza-Klein ansatz for the full bosonic sector of the theory.
The massive Type IIA supergravity action in the string frame goes like
$$l_s^8\sqrt{g}e^{2\mathrm{\Phi }}Q_4^{5/2},$$
(17)
suggesting that the number of degrees of freedom goes like $`Q_4^{5/2}`$ in the regime where it is an applicable description. Terms in the Type I’ action coming from the D8 brane DBI action turn out to be of the same order. Viewed from M theory point of view, we expect the corrections to the supergravity action to go like $`l_p^3l_s^3e^\mathrm{\Phi }1/Q_4`$, where $`l_p`$ is the eleven-dimensional Planck length. This seems to suggest that the field theory has a $`1/Q_4`$ expansion at large $`Q_4`$. For example, the one-loop correction of the form
$$l_s^8\sqrt{g}l_s^6^4Q_4^{1/2}$$
(18)
is suppressed by $`Q_4^2`$ compared to the tree-level action.
According to the $`AdS`$/CFT correspondence the spectrum of chiral primary operators of the fixed point theory can be derived from the spectrum of Kaluza-Klein excitations of massive Type IIA supergravity on the background (14). We will not carry out the detailed analysis here but make a few comments. The operators fall into representations of the $`F(4)`$ supergroup. As in the case of the six dimensional $`(0,1)`$ fixed point with $`E_8`$ global symmetry we expect $`E_{N_f+1}`$ neutral operators to match the Kaluza-Klein reduction of fields in the bulk geometry, and $`E_{N_f+1}`$ charged operators to match the Kaluza-Klein reduction of fields living on the boundary. Among the $`E_{N_f+1}`$ neutral operators we expect to have dimension $`3k/2`$ operators of the type $`Tr\varphi ^k`$ where $`\varphi `$ is a complex scalar in the hypermultiplet, which parametrize the Higgs branch of the theory. Like in we do not expect all these operators to be in short multiplets. We expect that those in long multiplets will generically receive $`1/Q_4`$ corrections to their anomalous dimensions. Unlike the hypermultiplet, the vector multiplet in five dimensions is not a representation of the superconformal group $`F(4)`$. Therefore, we do not expect Kaluza-Klein excitations corresponding to neutral operators of the type $`Tr\phi ^k`$ where $`\phi `$ is a real scalar in the vector multiplet, which parametrizes the Coulomb branch of the theory. Among the $`E_{N_f+1}`$ charged operators we expect to have the dimension four $`E_{N_f+1}`$ global symmetry currents that couple to the massless $`E_{N_f+1}`$ gauge fields on the boundary.
## 4 The M5-M5’ Brane System
We denote the eleven dimensional space-time coordinates by $`(x_{||},\stackrel{}{x},\stackrel{}{y},\stackrel{}{z})`$, where $`x_{||}`$ parametrize the $`(0,1,2,3)`$ coordinates, $`\stackrel{}{x}=(x_1,x_2)`$ the $`(4,5)`$ coordinates, $`\stackrel{}{y}=(y_1,y_2)`$ the $`(6,7)`$ coordinates and $`\stackrel{}{z}=(z_1,z_2,z_3)`$ the $`(8,9,10)`$ coordinates.
Consider two sets of fivebranes in M theory: $`N_1`$ coinciding M5 branes and $`N_2`$ coinciding M5’ branes. Their worldvolume coordinates are $`M5:(x_{||},\stackrel{}{y})`$ and $`M5^{}:(x_{||},\stackrel{}{x})`$. Such a configuraion preserves eight supercharges. The eleven-dimensional supergravity background takes the form
$`ds_{11}^2`$ $`=`$ $`(H_1H_2)^{2/3}[(H_1H_2)^1dx_{||}^2`$ (19)
$`+`$ $`H_2^1d\stackrel{}{x}^2+H_1^1d\stackrel{}{y}^2+d\stackrel{}{z}^2],`$
with the 4-form field strength $``$
$`=3(dH_1dy^1dy^2+dH_2dx^1dx^2),`$ (20)
where $``$ defines the dual form in the three dimensional space $`(z_1,z_2,z_3)`$. When the M5 and M5’ branes are only localized along the overall transverse directions $`\stackrel{}{z}`$ the harmonic functions are $`H_1=1+l_pN_1/2|\stackrel{}{z}|,H_2=1+l_pN_2/2|\stackrel{}{z}|`$, where $`l_p`$ is the eleven dimensional Planck length. This case has been discussed in . The near horizon geometry in this case does not have the AdS isometry group and thus cannot describe a dual SCFT. Moreover, there does not seem to be in this case a theory on the intersection which decouples from the bulk physics.
Consider the semi-localized case when the M5 branes are completely localized while the M5’ branes are only localized along the overall transverse directions. When the branes are at the origin of $`(\stackrel{}{x},\stackrel{}{z})`$ space the harmonic functions in the near core limit of the M5’ branes take the form
$$H_1=1+\frac{4\pi l_p^4N_1N_2}{(|\stackrel{}{x}|^2+2l_pN_2|\stackrel{}{z}|)^2},H_2=\frac{l_pN_2}{2|\stackrel{}{z}|}.$$
(21)
The numerical factors in (21) are determined by the requirment that the integral of $``$ yields the appropriate charges. This solution can also be obtained from the localized D2-D6 brane solution of by a chain of dualities.
It is useful to make a change of coordinates $`l_pz=(r^2\mathrm{sin}\alpha ^2)/2N_2,x=r\mathrm{cos}\alpha ,0\alpha \pi /2`$. In the near-horizon limit we want to keep the energy $`U=\frac{r}{l_p^2}`$ fixed. This implies, in particular, that membranes stretched between the M5 snd M5’ branes in the $`\stackrel{}{z}`$ direction have finite tension $`|\stackrel{}{z}|/l_p^3`$. Since we are smearing over the $`\stackrel{}{y}`$ directions we should keep $`\stackrel{}{y}/l_p`$ fixed. It is useful to make a change of variables $`\stackrel{}{y}/l_p\stackrel{}{y}`$. The near horizon metric is of the form of a $`warpedproduct`$ of $`AdS_5`$ and a six dimensional manifold $`_6`$
$`ds_{11}^2`$ $`=`$ $`l_p^2(4\pi N_1)^{1/3}(\mathrm{sin}^{2/3}\alpha )({\displaystyle \frac{U^2}{N_2}}dx_{||}^2`$ (22)
$`+`$ $`{\displaystyle \frac{4\pi N_1}{U^2}}dU^2+d_6^2),`$
with
$`d_6^2`$ $`=`$ $`4\pi N_1(d\alpha ^2+\mathrm{cos}^2\alpha d\theta ^2+{\displaystyle \frac{\mathrm{sin}^2\alpha }{4}}d\mathrm{\Omega }_2^2)`$ (23)
$`+`$ $`{\displaystyle \frac{N_2}{\mathrm{sin}^2\alpha }}(dy^2+y^2d\psi ^2).`$
The metric (22) has the $`AdS_5`$ isometry group . Therefore, in the spirit of , M theory on the background (22), (23) should be dual to a four dimensional $`𝒩=2`$ SCFT.
Note that the curvature of the metric diverges for small $`\alpha `$ as
$$\frac{1}{l_p^2N_1^{2/3}\mathrm{sin}^{8/3}\alpha }.$$
(24)
Away from $`\alpha =0`$, eleven-dimensional supergravity can be trusted for large $`N_1`$. The singularity at $`\alpha =0`$ is interpreted as a signal that some degrees of freedom have been effectively integrated and are needed in order to resolve the singularity. These presumably correspond to membranes that end on the M5’ branes.
The above M5-M5’ brane system can be understood as up lifting to eleven dimensions of an elliptic brane system of Type IIA. It consists of $`N_2`$ NS5-branes with worldvolume coordinates $`(0,1,2,3,4,5)`$ periodically arranged in the 6-direction and $`N_1`$ D4-branes with worldvolume coordinates $`(0,1,2,3,6)`$ stretched between them as in figure 1. When we lift this brane configuration to eleven dimensions, we delocalize in the eleven coordinate (which in our notation is 7) and since we have delocalization in coordinate 6 as well, we end up with the semi-localized M5-M5’ brane system.
The four dimensional theory at low energies on the D4-branes worldvolume is an $`SU(N_1)^{N_2}`$ gauge theory with matter in the bi-fundamentals. The metric (22) can be viewed as providing the eleven-dimensional supergravity description of an M5 brane with worldvolume $`R^4\times \mathrm{\Sigma }`$ , where $`\mathrm{\Sigma }`$ is the Seiberg-Witten holomorphic curve (Riemann surface) associated with the four dimensional SCFT at the origin of the moduli space of vacua . In this brane set-up the R-symmetry group $`SU(2)_R\times U(1)_R`$ is realized as the rotation group $`SU(2)_{8910}\times U(1)_{45}`$. The dimensionless gauge coupling of each $`SU(N_1)`$ part of the gauge group is $`g_{YM}^2g_sN_2`$.
The ten dimensional metric describing the elliptic Type IIA brane configuration is
$`ds_{10}^2`$ $`=`$ $`H_4^{1/2}dx_{||}^2+H_4^{1/2}d\stackrel{}{x}^2+H_4^{1/2}H_{NS}dy^2`$ (25)
$`+`$ $`H_4^{1/2}H_{NS}d\stackrel{}{z}^2,`$
where
$$H_4=1+\frac{4\pi g_sl_s^4N_1N_2}{[|\stackrel{}{x}|^2+2l_sN_2|\stackrel{}{z}|]^2},H_{NS}=\frac{l_sN_2}{2|\stackrel{}{z}|}.$$
(26)
Its near-horizon limit is
$$ds_{10}^2=l_s^2\left(\frac{U^2}{R^2}dx_{||}^2+R^2\frac{dU^2}{U^2}+d_5^2\right),$$
(27)
where
$`d_5^2`$ $`=`$ $`R^2(d\alpha ^2+\mathrm{cos}^2\alpha d\theta ^2+{\displaystyle \frac{\mathrm{sin}^2\alpha }{4}}d\mathrm{\Omega }_2^2)`$ (28)
$`+`$ $`{\displaystyle \frac{N_2^2}{R^2\mathrm{sin}^2\alpha }}dy^2,`$
and $`R^2=(4\pi g_sN_1N_2)^{1/2}`$. The curvature of the metric is $`\frac{1}{l_s^2R^2}`$ and it has a dilaton $`e^\mathrm{\Phi }=\frac{g_sN_2}{R}\mathrm{sin}\alpha ^1`$. This supergravity solution can be trusted when both the curvature and the dilaton are small. Away from $`\alpha =0`$ we need large $`N_2`$ and $`N_2N_1^{1/3}`$. When the latter condition is not satisfied the dilaton is large and we should consider the eleven-dimensional description.
The supergravity solution with both the M5 branes and the M5’ branes being completely localized has not been constructed yet. This fully localized system has been studied via a matrix description in . In the following we make a few comments on this branes configuration. It can be viewed as a single M5 brane with worldvolume $`R^4\times \mathrm{\Sigma }`$, where $`\mathrm{\Sigma }`$ is a singular Riemann surface defined by the complex equation
$$x^{N_2}y^{N_1}=0,$$
(29)
where $`x=x_1+ix_2,y=y_1+iy_2`$. Consider first the case $`N_1=N_2=1`$. The complex equation $`xy=0`$ can be written by a change of variables as $`w_1^2+w_2^2=0`$. This equation describes the singular Riemann surface at the point in the Seiberg-Witten solution where a monopole becomes massless. The massless spectrum consists of a $`U(1)`$ vector field and one hypermultiplet and the theory at this point is free in the IR. Indeed, the chiral ring associated with the singular variety $`xy=0`$ consists only of the identity operator. Since the field theory is free in the IR we expect the dual string description to be strongly coupled. That means that the dual background is highly curved and the classical supergravity analysis cannot be trusted. This theory has a dual string formulation . The more general $`N_1,N_2`$ case is harder to analyse. The reason being that the singular locus of the variety (29) is not isolated and the notion of a chiral ring which is used for an isolated singularity is not appropriate. We do expect a non trivial conformal field theory in this case, but it has not been identified yet.
## 5 BPS String Solitons
The centrally extended $`𝒩=2`$ supersymmetry algebra in four dimensions takes the form
$`\{Q_\alpha ^A,\overline{Q}_{\dot{\alpha }B}\}`$ $`=`$ $`\sigma _{\alpha \dot{\alpha }}^\mu P_\mu \delta _B^A+\sigma _{\alpha \dot{\alpha }}^\mu Z_{\mu B}^A,,`$
$`\{Q_\alpha ^A,Q_\beta ^B\}`$ $`=`$ $`\epsilon _{\alpha \beta }Z^{[AB]}+\sigma _{\alpha \beta }^{\mu \nu }Z_{\mu \nu }^{(AB)},`$ (30)
where $`A,B=1,2`$ and $`Z_{\mu A}^A=0`$. The $`SU(2)_R`$ part of the R-symmetry group acts on the indices $`A,B`$. In addition to the particle charge $`Z^{[AB]}`$, there are in (30) the string charges $`Z_{\mu B}^A`$ in the adjoint representation of $`SU(2)_R`$ and the membrane charges $`Z_{\mu \nu }^{(AB)}`$ in the 2-fold symmetric representation of $`SU(2)_R`$. Thus, $`𝒩=2`$ supersymmetric gauge theories in four dimensions can have BPS strings and BPS domain walls in addition to the well studied BPS particles.
The BPS particles and BPS strings are realized by stretching M2 brane between the M5 and M5’ branes. When the membrane worldvolume coordinates are $`(0,x,y)`$ where $`x`$ is one of the $`\stackrel{}{x}`$ components and $`y`$ is one of the $`\stackrel{}{y}`$ components we get a BPS particle of the four dimensional theory. It is charged under $`U(1)_R`$ that acts on $`\stackrel{}{x}`$ and is a singlet under $`SU(2)_R`$ that acts on $`\stackrel{}{z}`$. The BPS particles exist on the Coulomb branch of the gauge theory on which $`U(1)_R`$ acts, as in figure 3, and their mass is given by the Seiberg-Witten solution.
When the membrane worldvolume coordinates are $`(0,1,z)`$ where $`z`$ is one of the $`\stackrel{}{z}`$ components we get a BPS string of the four dimensional theory. It is not charged under $`U(1)_R`$ and it transforms in the adjoint $`SU(2)_R`$ since $`\stackrel{}{z}`$ transforms in this representation. The BPS strings exist on the Higgs branch of the gauge theory on which $`SU(2)_R`$ acts, as in figure 3. The tension of the BPS string is given by $`\frac{|\stackrel{}{z}|}{l_p^3}`$ where $`|\stackrel{}{z}|`$ is the distance between the M5 and M5’ branes and it is finite in the field theory limit. It is easy to see that we do not have in this set-up BPS domain walls since stretching a membrane with worldvolume coordinates $`(0,1,2)`$ breaks the supersymmetry completely. One expects BPS domain wall configurations when the moduli space of vacua has disconnected components.
In the Type IIA picture the BPS string is constructed by stretching a D2 brane between the D4 branes and NS5-brane. The seperation between these two types of branes in the $`\stackrel{}{z}`$ direction has two interpretations depending on whether the four dimensional gauge group has a $`U(1)`$ part or not . If is does then the separation is interpreted as the $`SU(2)_R`$ triplet FI parameters $`\stackrel{}{\zeta }`$. The BPS string tension is proportional to $`|\stackrel{}{\zeta }|`$. If the gauge group does not have a $`U(1)`$ part the seperation between these two types of branes in the $`\stackrel{}{z}`$ direction is interpreted as giving a vev to an $`SU(2)_R`$ triplet component of the meson $`\stackrel{~}{Q}Q`$ which transforms under $`SU(2)_R`$ as $`\mathrm{𝟐}\times \mathrm{𝟐}=\mathrm{𝟑}\mathrm{𝟏}`$. The BPS string tension is proportional to this vacuum expectation value. At the origin of the moduli space, where we have SCFT, the string becomes tensionless.
We can use the AdS/CFT correspondence in order to compute the self energy of a string and a potential between two such strings of opposite orientation . This is obtained, in the supergravity approximation, by a minimization of the M2 brane action
$$S=\frac{1}{(2\pi )^2l_p^3}𝑑\tau 𝑑\omega 𝑑\sigma \sqrt{detG_{\mu \nu }_ax^\mu _bX^\nu }$$
(31)
in the background (22), (23). Consider a static configuration $`x^0=\tau ,x^1=\omega `$ and $`x^i=x^i(\sigma )`$. Then, the energy per unit length of the string is given by
$$E=\frac{U\mathrm{sin}\alpha }{4\pi ^2N_2}\sqrt{dU^2+U^2d\alpha ^2}.$$
(32)
Thus, the string self-energy
$$E=(\alpha )\frac{U^2}{N_2}(\alpha )\frac{N_1}{(\delta x_{||})^2}.$$
(33)
The function $``$ parametrizes the dependence of the self-energy on the coordinate $`\alpha `$. From the field theory point of view, this parametrizes a dependence on the moduli that parametrize the space of vacua. $`\delta x_{||}`$ is a cut-off. In (33) we used the holographic relation between a distance $`\delta x_{||}`$ in field theory and the coordinate $`U`$ which in our case reads
$$\delta x_{||}\frac{(N_2N_1)^{1/2}}{U}.$$
(34)
This corresponds to the familiar relation
$$\delta x_{||}\frac{(g_{YM}^2N_1)^{1/2}}{U}.$$
(35)
Similarly, the potential energy per units length between two strings of opposite orientation seperated by distance $`L`$ is
$$E\frac{N_1}{L^2}.$$
(36)
here $`L`$ is distance between two strings. The results (33),(36) do not depend on $`N_2`$ which means that they do not depend on the gauge coupling of the four dimensional field theory $`g_{YM}^2N_2`$. This is not unexpected. The gauge coupling can be viewed as a vacuum expectation value of a scalar in the vector multiplet and does not appear in the hypermultiplet metric. Since the BPS strings exist on the Higgs branch it is natural that their potential and self-energy do not depend on the gauge coupling. However, (33) and (36) are only the large $`N_1`$ results and will presumably have $`1/N_1`$ corrections.
As noted above, the theory on the intersection of two sets of M5 branes containes BPS particles, which arise from M2 branes stretched in the directions $`(0,x,y)`$. Minimization of the M2 brane action in this case yields the potential between these two such objects of opposite charge separated by distance $`L`$
$$V\frac{(N_1N_2)^{1/2}}{L}.$$
(37)
Again, this is expected since for the four dimensional field theory it reads as $`V\frac{(g_{YM}^2N_1)^{1/2}}{L}`$.
The 11-dimentional supergravity action goes like
$$l_p^9\sqrt{g}N_1^2N_2$$
(38)
suggesting that the number of degrees of freedom goes like $`N_1^2N_2`$. This is also deduced from the two-point function of the stress energy tensor. This is, of course, expected for $`SU(N_1)^{N_2}`$ gauge theory. It is curious to note that there is similar growth ($`N^3`$) of the entropy for a system of $`N`$ parallel M5 branes .
The one loop correction has the following form $`l_p^9\sqrt{g}l_p^6^4N_2`$ which is suppressed by $`N_1^2`$ compared to the tree level action suggesting that the field theory has a $`1/N_1`$ expansion, which again is in agreement with the field theory expectation.
The Type IIA background (28) is T-dual to Type IIB on $`AdS_5\times S^5/Z_{N_2}`$ . The latter is the dual description of the $`Z_{N_2}`$ orbifold of $`𝒩=4`$ theory . For instance the number of degrees of freedom can be understood as $`c(𝒩=4)/N_2N_1^2N_2`$. We can identify the spectrum of chiral primary operators with the supergravity Kaluza-Klein excitations as analysed in .
###### Acknowledgments.
We would like to thank M. Alishahiha and A. Brandhuber for collaboration on the work presented in this talk.
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# The Bispectrum: From Theory to Observations
## 1 Introduction
It is widely accepted that the large-scale structure of the universe routinely observed in galaxy surveys is the result of gravitational instability that operates on (initially small) fluctuations in the density distribution. Essentially, initially dense (underdense) regions contract (expand) as a consequence of the attractive nature of gravity. As a result, most of the volume of the universe is in underdense regions whereas the rest contains dense structures such as clusters, galaxies, etc. This asymmetry between underdense and overdense regions implies non-Gaussianity in the galaxy distribution. Understanding the dynamics of gravitational instability thus gives us a powerful tool that can be used to extract useful information from the non-Gaussian clustering pattern of galaxies (Peebles 1980; Fry 1984; Goroff et al. 1986; Bernardeau 1992; Bouchet et al. 1995; Scoccimarro & Frieman 1996).
A quantitative description of galaxy clustering is obtained by applying statistical methods, in particular, multi-point correlation functions (Peebles 1980). In a Gaussian field, the statistical properties are fully characterized by its two-point correlation function or power spectrum, all higher-order (connected) correlation functions being zero. Information on non-Gaussianity is thus accessible by calculating higher-order correlations functions. An alternative to multi-point correlation functions is to consider higher-order moments (i.e. smoothed multi-point correlation functions collapsed at a single point); these characterize the one-point distribution of the galaxy distribution and can be measured with larger signal to noise than the corresponding multi-point functions. However, the latter contain more information, in particular, being evaluated at three or more points they are sensitive to the shapes of large-scale structures, which is the key to disentangle degeneracies present otherwise.
The bispectrum, the three-point correlation function in Fourier space, is the lowest order statistic sensitive to the shape of structures generated by gravitational instability. It can be used to probe the nature of primordial fluctuations (Gaussian versus non-Gaussian), galaxy biasing (i.e. the relation between the galaxy distribution and the underlying dark matter distribution), and help to break the degeneracy between linear bias and the matter density parameter $`\mathrm{\Omega }_m`$ present in power spectrum measurements in galaxy redshift surveys (Fry 1994; Hivon et al. 1995; Matarrese, Verde & Heavens 1997; Scoccimarro et al. 1998).
Pioneering work on three-point statistics from the Zwicky (Peebles & Groth 1975) and Lick (Groth & Peebles 1977; Fry & Seldner 1982) angular catalogs showed that the bispectrum, $`B(k_1,k_2,k_3)`$, scaled in terms of the power spectrum, $`P(k)`$, as $`B(k)P(k)^2`$. This motivatied the so-called hierarchical model and was also in agreement with the scaling at small scales inferred from the BBGKY equations of motion assuming stable clustering and self-similarity (Davis & Peebles 1977; Peebles 1980). At these small scales deep into the non-linear regime, however, limited understanding of the non-linear physics including the possibility of complicated galaxy biasing prevents using higher-order statistics to extract quantitative knowledge.
It was not until somewhat later, with the systematic development of non-linear perturbation theory (hereafter PT; Fry 1984), that it was possible to derive quantitatively how non-Gaussianity is generated by gravity. In particular, the scaling $`B(k)P(k)^2`$ (and in general, the $`n`$-point function $`\xi _n`$ scales as $`\xi _n\xi _2^{n1}`$) is recovered at large scales, being simply a result of quadratic non-linearities in the equations of motion of gravitational instability. An important signature predicted by PT is that the bispectrum at large scales should be a rather strong function of triangle configuration, with the reduced bispectrum $`Q_{123}`$
$$Q_{123}=\frac{B_{123}}{P_1P_2+P_2P_3+P_3P_1},$$
(1)
where $`B_{123}B(k_1,k_2,k_3)`$ and $`P_iP(k_i)`$, showing higher amplitude for colinear configurations than equilateral configurations, corresponding to the fact that gravity generates anisotropic large-scale structures. These results have been confirmed since then by N-body simulations (Fry, Melott & Shandarin 1993; Scoccimarro et al. 1998).
These predictions hold for the density field, in order to compare with observations one must necessarily deal with the possibility of galaxy biasing. At scales larger than those relevant to galaxy formation, one can make the simplifying assumption that biasing is a local function of the underlying smoothed density field, and thus expandable perturbatively at large scales as (Fry & Gaztañaga 1993; see also Coles 1993)
$$\stackrel{~}{\delta }_g(\text{x})=\underset{n=0}{}b_n\stackrel{~}{\delta }^n(\text{x}),$$
(2)
where $`b_n`$ are the bias parameters, and a tilde denotes the smoothing operation. Thus the galaxy bispectrum is calculable in terms of the density bispectrum at leading order in PT
$$Q_g=\frac{Q}{b_1}+\frac{b_2}{b_1^2}$$
(3)
in terms of linear ($`b_1`$) and non-linear ($`b_2`$) biasing parameters. Basically, a large linear bias decreases the dependence of the bispectrum on triangle, whereas antibias ($`b_1<1`$) enhances the configuration dependence of the galaxy bispectrum relative to that of the density. Since the density bispectrum depends to a good approximation only on spectral index, comparison of the density and galaxy bispectrum can give constraints on the biasing parameters independent of $`\mathrm{\Omega }_m`$ (Fry 1994). Note that $`Q_g`$ and $`Q`$ are functions that depend on triangle configuration; this is what allows to break the degeneracy between $`b_1`$ and $`b_2`$ present in analogous measurements of one-point moments using the skewness. Note that Eq. (3) assumes Gaussian initial conditions, otherwise there are additional contributions (Scoccimarro 2000).
A first attempt to use this procedure in the Lick catalog data, concluded that the lack of configuration dependence on the observed bispectrum required a large bias, $`b_13`$, assuming the absence of systematic errors (Fry 1994). However, other effects such as projection (Fry & Thomas 1999; Frieman & Gaztañaga 1999; Buchalter, Kamionkowski & Jaffe 2000) and most importantly non-linear evolution (Scoccimarro et al. 1998) also wash out the configuration dependence of the bispectrum, so application of this idea to observational data had to await until the availability of larger galaxy surveys. Recent determination of the angular three-point correlation function from the APM survey (Frieman & Gaztañaga 1999) shows a configuration dependence consistent with APM galaxies being unbiased ($`b_11`$, $`b_20`$), but with large error bars. This result is however consistent with the determination of one-point higher-order moments in the APM (Gaztañaga 1994), which show remarkable agreement with the predictions of PT from Gaussian initial conditions and no bias (Juszkiewicz, Bouchet & Colombi 1993; Bernardeau 1994a,1994b).
Redshift surveys map the three-dimensional distribution of galaxies in a large volume, and are thus ideally suited to use the bispectrum as a probe of galaxy biasing and non-Gaussian initial conditions. In this case, the observed distribution is actually distorted by peculiar velocities along the line of sight which contribute to redshifts. However, these redshift-space distortions are now well understood, at least in the plane-parallel approximation (Hivon et al. 1995, Verde et al. 1998; Scoccimarro, Couchman & Frieman 1999). Determination of three-point statistics from redshift surveys has been carried out for the CfA survey and a sample of redshifts in the Pisces-Perseus supercluster (Baumgart & Fry 1991) and the LCRS survey (Jing & Börner 1998). In both cases, however, there is a rather limited range of scales in the weakly non-linear regime to obtain quantitative constraints on biasing and non-Gaussian initial conditions. On the other hand, the next generation of large redshift surveys such as 2dF and SDSS will provide excellent conditions for the determination of the bispectrum at large scales.
There are however a number of issues which are potentially important for current and future surveys and must be addressed to reliably use the bispectrum to extract useful cosmological information. In this paper we consider effects that can potentially influence theoretical predictions of the bispectrum in redshift surveys: stochastic biasing, radial nature of redshift distortions, finite survey volume and sparse sampling. The main result of this study is a likelihood analysis that allows extracting accurate constraints on galaxy biasing and non-Gaussian initial conditions. In this work, we consider IRAS mock surveys, which map a large enough volume in the weakly non-linear regime. Application of our results to the actual data is considered in a companion paper (Scoccimarro et al. 2000).
This paper is organized as follows. In Section 2 we briefly review the behavior of the bispectrum at large scales from PT. In Section 3 we discuss the effects of galaxy biasing, including stochasticity. In Section 4 we review second-order Lagrangian PT, and discuss its regime of validity by comparing to N-body simulations. Section 5 presents results on the effects of redshift distortions in the plane-parallel and radial case. Section 6 discusses the generation of mock catalogs and the optimal weighting for higher-order statistics. In Section 7 we discuss the effects of finite volume of the survey and sampling on the bispectrum. Section 8 develops the likelihood analysis using bispectrum eigenmodes. Finally, Section 9 offers a summary of the results and our conclusions.
## 2 The Bispectrum induced by Gravity
If primordial density perturbations are Gaussian, they are fully characterized by their power spectrum
$$\delta (\text{k}_1)\delta (\text{k}_2)=[\delta _\mathrm{D}]_{12}P(k_1),$$
(4)
where $`[\delta _\mathrm{D}]_{1\mathrm{}N}(2\pi )^3\delta _D(\text{k}_1+\mathrm{}+\text{k}_N)`$, with $`\delta _D(\text{x})`$ the Dirac delta distribution. In this section we assume that the density field obeys statistical isotropy and homogeneity; these assumptions break down in presence of redshift-distortions, as will be discussed below. Even though the initial conditions are Gaussian, gravitational clustering is a non-linear process that induces non-Gaussianity through mode-mode coupling. In particular, the three-point cumulant in Fourier space, the bispectrum
$$\delta (\text{k}_1)\delta (\text{k}_2)\delta (\text{k}_3)_c=[\delta _\mathrm{D}]_{123}B_{123},$$
(5)
becomes non-zero, the leading order contribution arising from second-order PT (Fry (1984)). As discussed in the Introduction, since $`B(k)P(k)^2`$, it is convenient to work with the reduced bispectrum $`Q_{123}`$ (Fry & Seldner (1982)) defined as in Eq. (1), which has the desirable property that it is time and approximately scale independent to lowest order (tree-level) in non-linear PT, i.e.
$$B_{123}=2F_2(\text{k}_1,\text{k}_2)P_1P_2+\mathrm{cyc}.,$$
(6)
with
$$\delta _2(\text{k})=d^3qF_2(\text{q},\text{k}\text{q})\delta _1(\text{q})\delta _1(\text{k}\text{q}),$$
(7)
$$F_2(\text{k}_1,\text{k}_2)=\frac{5}{7}+\frac{\text{k}_1\text{k}_2}{2k_1k_2}\left(\frac{k_1}{k_2}+\frac{k_2}{k_1}\right)+\frac{2}{7}\left(\frac{\text{k}_1\text{k}_2}{k_1k_2}\right)^2,$$
(8)
where $`\delta _2(\text{k})`$ denotes the solution to the equations of motion of gravitational instability to second order in PT (Fry (1984); Goroff et al. (1986)). The result in Eq. (8) is for an Einstein-de Sitter model, but an important property of the kernel $`F_2`$ is its very weak dependence on density parameter $`\mathrm{\Omega }_m`$. A good approximation is to replace $`5/7(1+\kappa )/2`$ and $`2/7(1\kappa )/2`$, where for an open Universe $`\kappa \frac{3}{7}\mathrm{\Omega }_m^{2/63}`$, and for a flat Universe with cosmological constant, $`\kappa \frac{3}{7}\mathrm{\Omega }_m^{1/143}`$ (Bouchet et al. (1992, 1995)). For scale-free initial conditions ($`P(k)k^n`$), the density field reduced bispectrum at large scales thus becomes a function of only the spectral index $`n`$, and the shape of the triangle through e.g. the ratio between two sides $`k_1/k_2`$ and the angle between them.
So far we have discussed results for Gaussian initial conditions. When this assumption is relaxed, higher-order correlation functions are non-zero from the begining and the evolution of three-point statistics beyond linear PT is non-trivial (Fry & Scherrer (1994); Scoccimarro 2000). To illustrate this consider the tree-level evolution of the bispectrum from arbitrary non-Gaussian initial conditions
$`B_{123}^{(0)}`$ $`=`$ $`B_{123}^I+B_{123}^G+{\displaystyle d^3qF_2(\text{k}_1+\text{k}_2\text{q},\text{q})T_4^I(\text{k}_1,\text{k}_2,\text{k}_1+\text{k}_2\text{q},\text{q})},`$ (9)
where $`B_{123}^I`$ denotes the contribution of the initial bispectrum, scaled to the present time using linear PT, $`B_{123}^G`$ represents the usual gravitationally induced bispectrum, and the last term represents the contribution coming from the initial trispectrum linearly evolved to the present, $`T_4^I`$ given by
$$\delta ^I(\text{k}_1)\delta ^I(\text{k}_2)\delta ^I(\text{k}_3)\delta ^I(\text{k}_4)_c[\delta _\mathrm{D}]_{1234}T_4^I.$$
(10)
In this work we will consider one particular model of non-Gaussian initial conditions, the $`\chi ^2`$ model (e.g. Kofman et al. 1989; Antoniadis et al. 1997, Linde & Muhanov 1997; Peebles 1997).
## 3 Galaxy Biasing: Effects of Stochasticity
As discussed in the Introduction, when interpreting clustering statistics measured in galaxy surveys, one must necessarily deal with the issue of galaxy biasing. Equation (2) assumes not only that the bias is local (which seems a reasonable assumption at large scales) but also deterministic; that is, the galaxy distribution is completely determined by the underlying mass distribution. In practice, however, it is likely that galaxy formation depends on other variables besides the density field, and that consequently the relation between $`\stackrel{~}{\delta }_g(\text{x})`$ and $`\stackrel{~}{\delta }(\text{x})`$ is not deterministic but rather stochastic,
$$\stackrel{~}{\delta }_g(\text{x})=\underset{n=0}{}b_n\stackrel{~}{\delta }^n(\text{x})+\epsilon _\delta (\text{x}),$$
(11)
where the random field $`\epsilon _\delta (\text{x})`$ denotes the scatter in the biasing relation at a given $`\delta `$ due to the fact that $`\stackrel{~}{\delta }(\text{x})`$ does not completely determine $`\stackrel{~}{\delta }_g(\text{x})`$. Recent work has focussed on the effects of stochasticity on one-point statistics, the two-point correlation function and power spectrum (Scherrer & Weinberg 1998; Dekel & Lahav 1999). Clearly, if the field $`\epsilon (\text{x})`$ is arbitrary, the effects on clustering statistics could be rather strong. However, at large enough scales, we expect the field $`\epsilon _\delta (\text{x})`$ to be weakly correlated; that is, at large smoothing scales the scatter about the mean bias relation should be local. This means that $`\epsilon _\delta (\text{x}_1)\epsilon _\delta (\text{x}_2)\sigma _\epsilon ^2\mathrm{\Theta }(a|\text{x}_1\text{x}_2|)`$, where $`\mathrm{\Theta }(x)=1`$ if $`x>0`$ and zero otherwise, $`a`$ is the correlation “range” of the scatter and $`\sigma _\epsilon ^2`$ its strength. In this case, the power spectrum reads
$$P_g(k)=b_1^2P(k)+\sigma _\epsilon ^2V_aW_{\mathrm{TH}}(ka),$$
(12)
where $`V_a=4\pi a^3/3`$, and $`W_{\mathrm{TH}}(x)=3(\mathrm{sin}xx\mathrm{cos}x)/x^3`$ is the Fourier transform of the top-hat window. Here we used that, by definition, $`\epsilon \delta ^n=0`$. So, at large scales, $`ka1`$, stochastic bias leads to a constant offset to the power spectrum (Scherrer & Weinberg 1998; Dekel & Lahav 1999), similarly to Poisson fluctuations due to shot noise. However, in principle, the scatter at a given $`\delta `$ could be non-Poissonian (e.g. see Sheth & Lemson 1999). As long as the second term in Eq. (12) is small enough, the effect of local scatter should not very important. In this paper we do not explore general models of stochastic bias (see e.g. Scherrer & Weinberg 1998; Dekel & Lahav 1999; Blanton et al. 1999; Matsubara 1999), but rather test whether stochasticity affects the determination of bias parameters from bispectrum measurements that assume Eq. (3) in simple models where the scatter is local but not negligible when the “galaxy” density field is smoothed with a Gaussian filter of radius $`R_s10`$ Mpc/h. A similar approach in the context of higher-order one-point moments is given by Narayanan, Berlind, & Weinberg (2000), who also consider the effect of some non-local models of galaxy biasing.
We follow Cole et al. (1998) and generate non-linearly stochastic biased density fields by selecting dark matter particles to be “galaxies” with a probability function (see also Mann, Peacock, & Heavens 1998)
$`P(\nu )`$ $``$ $`\mathrm{exp}(a\nu +b\nu ^{3/2})(\nu 0)`$ (13)
$`P(\nu )`$ $``$ $`\mathrm{exp}(a\nu )(\nu <0),`$ (14)
where $`\nu (\text{x})\stackrel{~}{\delta }(\text{x})/\sigma `$ is the local normalized smoothed density field. The model has two free parameters $`a`$ and $`b`$ that can be adjusted to generate stochastic biased fields that are preferentially weighted towards low or high density regions. The mass density field is generated using 2LPT (see description in Section 4.2 below) in a 300 Mpc/h box, with a $`128^3`$ grid. The smoothed density field in Eqs.(13-refproba) then corresponds to about 3 Mpc/h smoothing. For a given particle, we choose the nearest grid point to evaluate $`\nu (\text{x})`$, thus the biasing scheme is effectively non-local at scales of about 3 Mpc/h.
Figure 1 shows four examples of such a procedure. The top panels correspond to bias towards high density regions, with $`b_1=b_2=1.54`$ (left panel) and $`b_1=1.45`$, $`b_2=1.06`$ (right panel), whereas the bottom panels show examples of bias that preferentially select low density regions with $`b_1=0.95`$, $`b_2=0.06`$ (left panel) and $`b_1=0.77`$, $`b_2=0.14`$ (right panel). All except the bottom right panel are from $`\mathrm{\Lambda }`$CDM realizations ($`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$) with $`\sigma _8=0.70`$ for the dark matter. The bottom right panel corresponds to a $`\tau `$CDM realization ($`\mathrm{\Omega }_m=1`$) with same normalization and power spectrum shape. The two right panels have roughly the same redshift distortions parameter $`\beta \mathrm{\Omega }_m^0.6/b_10.65`$, and thus will be nearly equal in terms of their redshift-space power spectrum. These bias parameter values in Figure 1 have been obtained by measuring the bispectrum for the “galaxy” distribution and comparing to the mass bispectrum using Eq. (3), and are plotted in each panel as the continuous solid line. The actual mean of the relation and its 90% confidence region for smoothed fields with Gaussian filter radius $`R_s=10`$Mpc/h show that even when the $`\delta \delta _g`$ relation has considerable scatter, the bispectrum recovers the mean of the biasing relation quite accurately. Below, we shall discuss whether this result is also valid when dealing with noisy data such as in a galaxy survey, including the effects of survey geometry and sampling.
## 4 Numerical Implementation of PT: Comparison with N-body Simulations
### 4.1 The Need for Numerical PT
In order to understand how geometry and sampling affects statistics such as the bispectrum, we need to simulate galaxy catalogs and measure their statistical properties to compare with the ideal case known from analytical PT calculations and N-body simulations. As we shall see, it will be become very valuable to make a large number (a few hundred to a thousand) of mock catalogs for different cosmological models and survey selection functions. As a result, it would be very expensive to proceed with N-body simulations, even with PM simulations. Since we are interested in the large-scale properties of clustering, where PT applies, it is natural to resort to a numerical implementation of PT.
Lagrangian PT is the natural choice. In this formulation, particles are displaced from their initial positions by a displacement field that is found order by order perturbatively from the equations of motion. The linear solution is the well-known Zel’dovich (1970) approximation (ZA), and higher-order corrections have been well studied in the literature (Moutarde et al. (1991); Buchert et al. (1994); Bouchet et al. (1995)). Although the ZA is the least expensive of them, it does not reproduce very accurately the statistical properties of clustering even at large scales (Grinstein & Wise (1987); Juszkiewicz et al. (1993); Bernardeau 1994b ; Catelan & Moscardini (1994); Juszkiewicz et al. (1995)). In particular, we are interested in studying the bispectrum, and the ZA overestimates its configuration dependence and underestimates its overall amplitude. The situation becomes even worse as higher-order statistics are considered. To illustrate this, Table 1 shows the ratio of the $`S_p`$ parameters in the ZA to their exact values at large scales as a function of spectral index $`n`$. The $`S_p`$ parameters characterize the one-point PDF of the smoothed density field at scale $`R`$ (Goroff et al. (1986))
$$S_p\frac{\stackrel{~}{\delta }^p_c}{\stackrel{~}{\delta }^2^{p1}},$$
(15)
and Table 1 shows $`r_pS_p^{\mathrm{ZA}}/S_p`$, where $`S_p`$ is the exact value calculated in full PT (Juszkiewicz et al. (1993); Bernardeau 1994a ; Bernardeau 1994b ). We see that for $`p=3`$ the skewness is poorly reproduced as the spectral index becomes positive, in the relevant range for CDM models at large scales. As a result, using the ZA to generate mock catalogs to study the statistical properties of the bispectrum could result in serious quantitative error.
The next level of accuracy is second-order Lagrangian PT (2LPT), which reproduces the three-point statistics (skewness and bispectrum) exactly, and gives reasonable values for higher-order statistics as well, in remarkable improvement over the ZA (Juszkiewicz et al. (1995)). Table 2 shows the analogous results for $`r_p`$ in 2LPT (Scoccimarro (1998)). Given the small additional computational expense of 2LPT over ZA (see discussion below), it is a very convenient scheme to implement for mock catalog production. It is still orders of magnitude faster than a PM simulation. The question then arises, is it worth to consider going beyond 2LPT? At third-order, 3LPT recovers four-point statistics exactly (see Table 3), but it involves an appreciable increase in complexity compared to 2LPT, requiring to solve three additional Poisson’s equations (Buchert et al. (1994)). For this reason, we consider 2LPT the optimal choice between speed and accuracy. We now briefly discuss 2LPT and its implementation.
### 4.2 Second-Order Lagrangian PT (2LPT)
In Lagrangian PT, the object of interest is the displacement field $`𝚿(\text{q})`$ which maps the initial particle positions q into the final Eulerian particle positions x,
$$\text{x}=\text{q}+𝚿(\text{q}).$$
(16)
The equation of motion for particle trajectories $`\text{x}(\tau )`$ is
$$\frac{d^2\text{x}}{d\tau ^2}+(\tau )\frac{d\text{x}}{d\tau }=\mathrm{\Phi },$$
(17)
where $`\mathrm{\Phi }`$ denotes the gravitational potential, and $``$ the gradient operator in Eulerian coordinates x. Taking the divergence of this equation we obtain
$$J(\text{q},\tau )\left[\frac{d^2\text{x}}{d\tau ^2}+(\tau )\frac{d\text{x}}{d\tau }\right]=\frac{3}{2}\mathrm{\Omega }_m^2(J1),$$
(18)
where we have used Poisson equation together with the fact that $`1+\delta (\text{x})=J^1`$, and the Jacobian $`J(\text{q},\tau )`$ is the determinant
$$J(\text{q},\tau )\mathrm{Det}\left(\delta _{ij}+\mathrm{\Psi }_{i,j}\right),$$
(19)
where $`\mathrm{\Psi }_{i,j}\mathrm{\Psi }_i/\text{q}_j`$. Equation (18) can be fully rewritten in terms of Lagrangian coordinates by using that $`_i=(\delta _{ij}+\mathrm{\Psi }_{i,j})^1_{q_j}`$, where $`_q/\text{q}`$ denotes the gradient operator in Lagrangian coordinates. The resulting non-linear equation for $`𝚿(\text{q})`$ is then solved perturbatively, expanding about its linear solution, the Zel’dovich (1970) approximation
$$_q𝚿^{(1)}=D_1(\tau )\delta (\text{q}),$$
(20)
which incorporates the kinematic aspect of the collapse of fluid elements in Lagrangian space. Here $`\delta (\text{q})`$ denotes the (Gaussian) density field imposed by the initial conditions and $`D_1(\tau )`$ is the linear growth factor. The solution to second order describes the correction to the ZA displacement due to gravitational tidal effects and reads
$$_q𝚿^{(2)}=\frac{1}{2}D_2(\tau )\underset{ij}{}(\mathrm{\Psi }_{i,i}^{(1)}\mathrm{\Psi }_{j,j}^{(1)}\mathrm{\Psi }_{i,j}^{(1)}\mathrm{\Psi }_{j,i}^{(1)}),$$
(21)
(e.g., Bouchet et al. 1995) where $`D_2(\tau )`$ denotes the second-order growth factor, which for $`0.1\mathrm{\Omega }_m3`$ ($`\mathrm{\Lambda }=0`$) obeys
$$D_2(\tau )\frac{3}{7}D_1^2(\tau )\mathrm{\Omega }_m^{2/63}\frac{3}{7}D_1^2(\tau )$$
(22)
to better than 0.5% and 7%, respectively (Bouchet et al. 1995), whereas for flat models with non-zero cosmological constant $`\mathrm{\Lambda }`$ we have for $`0.01\mathrm{\Omega }_m1`$
$$D_2(\tau )\frac{3}{7}D_1^2(\tau )\mathrm{\Omega }_m^{1/143}\frac{3}{7}D_1^2(\tau ),$$
(23)
to better than 0.6% and 2.6%, respectively (Bouchet et al. 1995). Since Lagrangian solutions up to second-order are curl-free, it is convenient to define Lagrangian potentials $`\varphi ^{(1)}`$ and $`\varphi ^{(2)}`$ so that in 2LPT
$$\text{x}(\text{q})=\text{q}D_1_q\varphi ^{(1)}+D_2_q\varphi ^{(2)},$$
(24)
and the velocity field then reads ($`t`$ denotes cosmic time)
$$𝐯\frac{d\text{x}}{dt}=D_1f_1H_q\varphi ^{(1)}+D_2f_2H_q\varphi ^{(2)},$$
(25)
where $`H`$ is the Hubble constant, and the logarithmic derivatives of the growth factors $`f_i(d\mathrm{ln}D_i)/(d\mathrm{ln}a)`$ can be approximated for open models with $`0.1\mathrm{\Omega }_m1`$ by
$$f_1\mathrm{\Omega }_m^{3/5},f_22\mathrm{\Omega }_m^{4/7},$$
(26)
to better than 2% (Peebles 1976) and 5% (Bouchet et al. 1995), respectively. For flat models with non-zero cosmological constant $`\mathrm{\Lambda }`$ we have for $`0.01\mathrm{\Omega }_m1`$
$$f_1\mathrm{\Omega }_m^{5/9},f_22\mathrm{\Omega }_m^{6/11},$$
(27)
to better than 10% and 12%, respectively (Bouchet et al. 1995). The accuracy of these two fits improves significantly for $`\mathrm{\Omega }_m0.1`$, in the range relevant for the present purposes. The time-independent potentials in Eqs. (24) and (25) obey the following Poisson equations (Buchert et al. 1994)
$`_q^2\varphi ^{(1)}(\text{q})`$ $`=`$ $`\delta (\text{q}),`$ (28)
$`_q^2\varphi ^{(2)}(\text{q})`$ $`=`$ $`{\displaystyle \underset{i>j}{}}[\varphi _{,ii}^{(1)}(\text{q})\varphi _{,jj}^{(1)}(\text{q})(\varphi _{,ij}^{(1)}(\text{q}))^2],`$ (29)
### 4.3 Regime of Validity
In order to assess the validity of 2LPT, we have run N-body simulations using the Hydra adaptive P<sup>3</sup>M code (Couchman, Thomas & Pearce (1995)). We have run 4 realizations of $`\mathrm{\Lambda }`$CDM ($`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`\sigma _8=0.90`$) and 2 realizations of SCDM ($`\sigma _8=0.61`$), using a 300 Mpc/h box, with $`128^3`$ particles and softening length $`ϵ=0.25`$ Mpc/h. The 2LPT realizations were made using in most cases $`256^3`$ particles in a 600 Mpc/h box (identical results are obtained using 300 Mpc/h boxes).
We now use the results of these simulations to check the regime of validity of 2LPT. We shall concentrate on power spectrum and bispectrum statistics, in real and redshift space. When dealing with Lagrangian PT, a signature of its breakdown is the amount of shell crossing; the fact that the Jacobian in Eq. (19) becomes zero at the regions where particles (initially at different positions) cross each other. This effect is well-known in the ZA, and leads to the “thickening of pancakes”, that is, to significant smoothing of high density regions. A way of characterizing this effect is to recall that the density field in the ZA is given by
$$1+\delta (\text{x})=\frac{1}{J(\text{q},\tau )}=\frac{1}{_{i=1}^3(1+\lambda _i(\text{q},\tau ))},$$
(30)
where $`\lambda _i`$ denotes the eigenvalues of the first-order deformation tensor $`\mathrm{\Psi }_{i,j}^{(1)}`$ at point q and time $`\tau `$ \[see Eq. 19\]. The condition for shell-crossing can then be explicitly written as $`1+\lambda _i(\text{q},\tau )=0`$, for any $`i=1,2,3`$. At these points, one expects not only a breakdown of the ZA, but also of higher-order Lagrangian PT, in particular 2LPT. An alternative indicator of the breakdown of 2LPT, as in any perturbative approach, is to compare the magnitude of the second-order to the first-order displacement field. One expects 2LPT will break down when the second-order contribution becomes significant compared to the ZA.
Table 4 shows these indicators in 2LPT realizations as described above for $`\mathrm{\Lambda }`$CDM with normalization $`\sigma _8=0.7`$. In order to investigate the dependence of the breakdown of 2LPT with the amount of small-scale power, we include a cutoff in the linear power spectrum so that $`P_{lin}(k)=0`$ for $`k>k_c`$. As more high-frequency waves are included in the representation of the linear density field, one expects 2LPT to break down at larger scales. On the other hand, not enough waves lead to lack of mode-coupling, and therefore inaccurate representation of the density field one is trying to simulate. Thus there is an “optimum” cutoff; here we considered models with $`k_c=0.3,0.4,0.5,\mathrm{}`$ h/Mpc.
In Table 4 we show much shell crossing is developed as a function of $`k_c`$, in terms of the quantity $`x_{sc}100\times N_{sc}/N_{grid}^3`$, where $`N_{sc}`$ is the number of grid points that obeyed the shell-crossing condition ($`|\lambda _i|1`$ for at least one $`i=1,2,3`$), and $`N_{grid}`$ is the total number of grid points (where particles are initially located, $`N_{par}=N_{grid}^3`$). Obviously, $`x_{sc}`$ is an increasing function of $`k_c`$. The remaining column in Table 4 shows the ratio of mean second to first-order displacement field magnitudes. Again, this ratio increases with $`k_c`$, as expected. We see that even though this ratio is always less than unity, significant shell crossing develops (e.g. for $`k_c>0.5`$), so breakdown of 2LPT is nonetheless expected to happen. In the following we use 2LPT realizations with $`k_c=0.5`$ h/Mpc, which we found reproduces clustering statistics in redshift space very well.
In Figure 2 we show a comparison of 2LPT (solid lines) with N-body simulations (symbols) in redshift space. The top panel shows the power spectrum (in terms of $`\mathrm{\Delta }(k)=4\pi k^3P(k)`$) as a function of scale, the dotted line corresponds to linear PT, Eq. 33 below (Kaiser 1987). At large scales, $`k<0.2`$ h/Mpc, 2LPT reproduces the suppresion of power compared to linear PT; at smaller scales the supression is overestimated, as we approach the cutoff scale of the power spectrum and also shell crossing starts playing a role. The bottom panel shows a similar comparison for the equilateral bispectrum, the dotted line now corresponds to the redshift space prediction at large scales from tree-level PT (Scoccimarro et al. 1999). We see that although tree-level PT breaks down at about $`k0.2`$ h/Mpc, 2LPT continues to hold at smaller scales, $`k0.4`$ h/Mpc. Figure 3 shows a comparison for the redshift space bispectrum for configurations in which $`k_2=2k_1=0.21`$ h/Mpc as a function of the angle $`\theta `$ between $`\text{k}_1`$ and $`\text{k}_2`$. The dotted line shows the prediction of tree-level PT (Hivon et al. 1995). As noted before (Scoccimarro et al 1999), the predictions of tree-level PT in redshift space break down at relatively large scales, due to the non-perturbative nature of the redshift-space mapping. Fortunately, since 2LPT performs the mapping exactly, the agreement with the numerical simulations is very good down to scales $`k0.4`$ h/Mpc. For our purposes, we are interested in studying statistics at large scales, $`k0.2`$ h/Mpc (with $`k_{nl}>0.2`$ h/Mpc), so the use of 2LPT is well justified.
## 5 Redshift Distortions
### 5.1 Plane-Parallel Approximation
In redshift space, the radial coordinate s of a galaxy is given by its observed radial velocity, a combination of its Hubble flow plus “distortions” due to peculiar velocities. The mapping from real-space position x to redshift space is given by:
$$\text{s}=\text{x}fu_z(\text{x})\widehat{z},$$
(31)
where $`f(\mathrm{\Omega }_m)\mathrm{\Omega }_m^{0.6}`$ is the logarithmic growth rate of linear perturbations, and $`\text{u}(\text{x})\text{v}(\text{x})/(f)`$, where $`\text{v}(\text{x})`$ is the peculiar velocity field, and $`(\tau )(1/a)(da/d\tau )=Ha`$ is the conformal Hubble parameter (with FRW scale factor $`a(\tau )`$ and conformal time $`\tau `$). In Eq. (31), we have assumed the “plane-parallel” approximation, so that the line-of-sight is taken as a fixed direction, denoted by $`\widehat{z}`$. Using this mapping, the Fourier transform of the density field contrast in redshift space reads (Scoccimarro et al. (1999))
$$\delta _s(\text{k})=\frac{d^3x}{(2\pi )^3}\mathrm{e}^{i\text{k}\text{x}}\mathrm{e}^{ifk_zu_z(\text{x})}\left[\delta (\text{x})+f_zu_z(\text{x})\right].$$
(32)
This equation describes the fully non-linear density field in redshift space in the plane-parallel approximation. In linear perturbation theory, the exponential factor becomes unity, and we recover the well known formula (Kaiser (1987))
$$\delta _s(\text{k})=\delta (\text{k})(1+f\mu ^2).$$
(33)
Given the general formula, Eq. (32), and the local biasing scheme in Eq. (2) (with $`\epsilon (\text{x})=0`$), it is easy to work out the perturbative solutions order by order (Scoccimarro et al. (1999)). Let’s write the Fourier components of the density field in redshift space as
$`\delta _s(\text{k})`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}D_1^n{\displaystyle [\delta _\mathrm{D}]_nZ_n(\text{k}_1,\mathrm{},\text{k}_n)\underset{i=1}{\overset{n}{}}\delta _1(\text{k}_i)d^3k_i},`$
where $`Z_n`$ are dimensionless, symmetric, scalar functions of their arguments, $`D_1`$ is the linear growth factor of the density field, and $`[\delta _\mathrm{D}]_n\delta _D(\text{k}\text{k}_1\mathrm{}\text{k}_n)`$. In particular the first and second-order PT kernels are (Hivon et al. (1995); Verde et al. (1998); Scoccimarro et al. (1999))
$`Z_1(\text{k})`$ $`=`$ $`(b_1+f\mu ^2)`$ (35)
$`Z_2(\text{k}_1,\text{k}_2)`$ $`=`$ $`{\displaystyle \frac{b_2}{2}}+b_1F_2(\text{k}_1,\text{k}_2)+f\mu ^2G_2(\text{k}_1,\text{k}_2)+{\displaystyle \frac{f\mu k}{2}}\left[{\displaystyle \frac{\mu _1}{k_1}}(b_1+f\mu _2^2)+{\displaystyle \frac{\mu _2}{k_2}}(b_1+f\mu _1^2)\right],`$ (36)
where we denote $`\mu \text{k}\widehat{z}/k`$, with $`\text{k}\text{k}_1+\mathrm{}+\text{k}_n`$, and $`\mu _i\text{k}_i\widehat{z}/k_i`$. As in Eq. (7), $`G_2`$ denotes the second-order kernels for the velocity-divergence field, given by Eq. (8) after replacing $`5/73/7`$ and $`2/74/7`$. We can now write the power spectrum and bispectrum in redshift-space,
$$P_s(\text{k})=Z_1(\text{k})^2P(k),$$
(37)
$$B_{123}=2Z_2(\text{k}_1,\text{k}_2)Z_1(\text{k}_1)Z_1(\text{k}_2)P(k_1)P(k_2)+\mathrm{cyc}.$$
(38)
Note that both statistics do now depend on the direction of the wave vectors with respect to the line of sight, since the redshift-space mapping in the plane-parallel approximation breaks statistical isotropy. One can proceed to decompose these statistics in multipole moments; however, in this paper we will be dealing exclusively with their monopoles. The resulting reduced bispectrum $`Q`$, defined as in Eq. (1) in terms of monopole quantities, now becomes a function of $`\beta f/b_1`$, $`b_1`$ and $`b_2`$. Even though Eq. (3) is not valid anymore in redshift-space, it is still a reasonable approximation (Scoccimarro et al. (1999)). As a result, even including redshift-distortions, the reduced bispectrum $`Q`$ is weakly sensitive to $`\mathrm{\Omega }_m`$ (or $`\beta `$) (Hivon et al. (1995); Verde et al. (1998); Scoccimarro et al. (1999)).
### 5.2 Effect of Radial Redshift-Distortions
A standard procedure with redshift distortions is to treat them in the plane-parallel approximation (Kaiser (1987)). When dealing with surveys with substantial sky coverage, this approximation is expected to break down. In an all-sky survey (as IRAS ) the radial character of redshift distortions mean that clustering statistics are statistically isotropic but inhomogeneous (rather than statistically homogeneous but anisotropic, as in the plane-parallel case). The breakdown of statistical homogeneity means that Fourier modes are not independent anymore, in particular, the power spectrum is not diagonal, e.g. different band powers are correlated, even in the linear regime (Zaroubi & Hoffman (1996)).
We will show, however, that for the monopole of the power spectrum and bispectrum in the all-sky case, radial and plane-parallel distortions agree with each other extremely well. This can be understood from the fact that the monopole is not sensitive to the orientation of structures but to the overall spherical average of power. On the other hand, higher-order multipoles (and thus the full unaveraged over angles statistics) such as the quadrupole are much more sensitive to the radial character of the distortions; in fact, in the limit of full sky, the anisotropy (with respect to a fixed direction) vanishes.
Figure 4 illustrates this for the monopole of the power spectrum. Solid lines denote the undistorted non-linear power spectrum, whereas dotted (dashed) lines denote the power spectrum monopole under radial (plane-parallel) redshift-space mapping. The radial results overlap with the plane-parallel ones except for a very small difference at large scales, where the radial results are slightly smaller. This is indeed expected at the 10% level for the scales plotted (Heavens & Taylor (1995)). A similar result holds for the scale dependence of the bispectrum. In the top left panel of Fig. 5 we show the ratio of the bispectrum under radial ($`Q_R`$) to plane-parallel ($`Q_P`$) redshift-space mapping as a function of the angle $`\theta `$ between $`\text{k}_1`$ and $`\text{k}_2`$ for $`k_2=2k_1=0.105`$ h/Mpc (triangles) and $`k_2=2k_1=0.21`$ h/Mpc (squares). In order to get this result we have averaged over 200 realizations of $`\mathrm{\Lambda }`$CDM. There is an apparent trend at the few percent level that indicates that under the radial mapping the configuration dependence is very slightly suppressed, but the effect is very small even at these large scales. In the following we use radial distortions, although for our purposes plane-parallel approximation would have been a good approximation. We use Fourier modes even though they are not the natural set of modes anymore in the presence of radial redshift-distortions (see e.g. Heavens & Taylor (1995); Hamilton & Culhane 1996).
### 5.3 Dependence of $`Q`$ on Cosmology and Power Spectrum
As discussed before, $`Q`$ depends mainly on the spectral index, biasing, Gaussianity of initial conditions, and it is quite insensitive to the matter density $`\mathrm{\Omega }_m`$ and should be independent to leading order on the power spectrum normalization $`\sigma _8`$. We now quantify these dependences in the presence of non-linear redshift distortions. Figure 5 shows the results of averaging 200 realizations of 2LPT, as described before, for different models. The top right panel shows the dependence of $`Q`$ on the matter density parameter $`\mathrm{\Omega }_m`$, for power spectrum shape given by $`\mathrm{\Gamma }=0.21`$. The scatter plot corresponds to measurements of $`Q`$ in all triangles (4741 in total) between scales of $`k=0.05`$ h/Mpc and $`k=0.2`$ h/Mpc. Equilateral triangles correspond to the small amplitude region (lower left corner), whereas colinear triangles ocupy the upper right corner. As predicted by tree-level PT, the $`\mathrm{\Omega }_m=1`$ case shows a higher colinear amplitude and lower equilateral amplitude than the $`\mathrm{\Omega }_m=0.3`$ case, although the precise value of these amplitudes does not quite agree with tree-level PT in the presence of non-linear redshift distortions, as we saw in Fig. 3. The effect of a change in $`\mathrm{\Omega }_m`$ is of the order of $`10\%`$ (Hivon et al. 1995).
The bottom left panel shows the effect of changing $`\sigma _8`$ for $`\mathrm{\Lambda }`$CDM models. As we see, non-linear redshift distortions introduce a very small dependence, with larger $`\sigma _8`$ leading to a slightly smaller configuration dependence of $`Q`$, but the effect is negligible for the analysis of IRAS surveys given the error bars involved, as we shall see. On the other hand, as expected, the dependence on the shape of the power spectrum, parameterized by $`\mathrm{\Gamma }`$ in CDM models is quite strong due to a change of the spectral index with scale. We see from the bottom right panel that for $`\mathrm{\Omega }_m=1`$, a $`\mathrm{\Gamma }=0.5`$ has a much weaker configuration dependence than a $`\mathrm{\Omega }_m=0.21`$, as expected since the spectral index of the latter is more negative at a given scale. The solid line in this panel shows the relation $`Q(\mathrm{\Gamma }=0.21)=1.3Q(\mathrm{\Gamma }=0.5)0.05`$, which fits the results quite well. As a result of this, a $`\mathrm{\Gamma }=0.5`$ model when compared with a given galaxy bispectrum will lead to smaller values for the bias parameters $`b_1`$ and $`b_2`$ than a $`\mathrm{\Gamma }=0.21`$ model, assuming no other constraint is imposed.
## 6 Mock Catalogs: Selection Functions and Optimal Weighing
We generate mock catalogs of different surveys by the following procedure. We use 2LPT to generate density and velocity fields for a given cosmological model and normalization $`\sigma _8`$. In most cases we use $`128^3`$ particles in a 600 Mpc/h box, for biased galaxy distributions we use a parent $`200^3`$ 2LPT realization from which an $`128^3`$ distribution is generated (see below). An observer is then picked at random with velocity consistent with that of the local group with respect to the CMB frame. The full (radial) redshift-space mapping is performed using the velocity field, to obtain the redshift-space distribution. Given the particle distribution in redshift-space, “galaxies” are selected using the desired IRAS survey selection function, replicating the box if necessary, including the galactic cut, and rejecting galaxies closer than 20 Mpc/h from the observer. The selection functions are given by:
$$n(r)=\frac{A}{r^2}\frac{\left(r/r_0\right)^a}{1+\left(r/r_0\right)^b},$$
(39)
with parameters given in Table 5. These were kindly provided by H. Feldman, and obtained by fitting the actual redshift distribution observed in the different IRAS surveys.
In order to correct for the selection function and the geometry of the survey, we use the standard FKP method (Feldman et al. (1994); hereafter FKP). For each realization of the mock survey, a corresponding random catalogue is generated using the same geometry and selection function, but about ten times more dense. We then measure statistics (power spectrum and bispectrum monopoles) for the auxiliary $`F(\text{k})`$ field,
$$F(\text{k})=d^3xw_k(\text{x})\left[n_g(\text{x})\alpha n_r(\text{x})\right]\mathrm{exp}(i\text{k}\text{x}),$$
(40)
where $`n_g(\text{x})`$ and $`n_r(\text{x})`$ are the galaxy and random catalog density fields, and $`\alpha =N_g/N_r`$ scales the overall random density to the survey density. Note that the use of a random catalog is not necessary, one could just as well use the survey selection function multiplied by the survey mask to substract the expected shot noise; or even better, substract the “actual” shot noise using the data itself (Hamilton 1997). The weight function $`w_k(\text{x})1/\left[1+\overline{n}P(k)\right]`$ is chosen as to minimize the variance of the power spectrum estimator at scales small compared to the survey size (FKP). In fact, under the same assumptions plus Gaussian errors, the same weight function minimizes the variance of the higher-order correlation functions as well. Indeed, the variance $`(\mathrm{\Delta }\widehat{T}_m)^2`$ in the estimator of the $`m`$point spectrum $`\widehat{T}_m`$ in this limit is
$$(\mathrm{\Delta }\widehat{T}_m)^2\frac{d^3xn^{2m}(\text{x})_{i=1}^m\left[P(k_i)+1/n\right]w_{k_i}^2(\text{x})}{\left[d^3xn^m(\text{x})_{i=1}^mw_{k_i}^2(\text{x})\right]^2},$$
(41)
and taking variations with respect to the weight function $`w_{k_j}(\text{r})`$ it follows that
$$\frac{w_{k_j}(\text{r})[1+n(\text{r})P_j]_{ij}^m[1+n(\text{r})P_i]w_{k_i}^2(\text{r})}{_{ij}^mw_{k_i}(\text{r})}=\frac{d^3xn^m(\text{x})_{i=1}^m[1+n(\text{x})P_i]w_{k_i}^2(\text{x})}{d^3xn^m(\text{x})_{i=1}^mw_{k_i}(\text{x})},$$
(42)
whose solution is the FKP weight function $`w_k(\text{x})=1/(1+n(\text{x})P(k))`$ (this result was derived independently by M. Zaldarriaga, private communication 2000). In practice however, when the FKP approximation breaks down, a more complicated procedure must be used to determine the optimal weight $`w_k(\text{x})`$, see e.g. Colombi, Szapudi & Szalay (1998) for one-point higher-order statistics. From Eq. (40) it follows that
$$P_F(k)=d^3k^{}P_G(|\text{k}\text{k}^{}|)P(\text{k}^{})+(1+\alpha )d^3xw_k(\text{x})\overline{n}(\text{x})^2,$$
(43)
where $`\overline{n}(\text{x})=n_g(\text{x})`$, $`P_F(k)`$ and $`P_G(k)`$ are the power spectra of the fields $`F`$ and $`G(\text{x})w_k(\text{x})\overline{n}(\text{x})`$. For IRAS surveys, $`P_G(k)`$ decays quite strongly with k, as seen in Fig. 6 for QDOT-PSCz (dashed), 1.2Jy (solid) and 2Jy (dot-dashed), where we assumed for definiteness $`w_k(\text{x})=1/(1+n(\text{x})P_0)`$ with $`P_0=2000`$. As a first approximation we can thus treat $`P_G(k)`$ as a delta function, and we obtain the power spectrum and bispectrum estimators (Matarrese et al. 1997)
$$\widehat{P}(k)=\frac{P_F(k)}{I_{22}}(1+\alpha )\frac{I_{12}}{I_{22}},,$$
(44)
$$\widehat{B}_{123}=\frac{B_F}{I_{33}}(\widehat{P}_1+\widehat{P}_2+\widehat{P}_3)\frac{I_{23}}{I_{33}}(1\alpha ^2)\frac{I_{13}}{I_{33}},,$$
(45)
where $`I_{ab}d^3x\overline{n}^a(\text{x})w_{k_1}(\text{x})\mathrm{}w_{k_b}(\text{x})`$ and $`\alpha `$ denotes the additional shot noise correction due to the use of a random catalog. In the volume limited case, $`I_{ab}=\overline{n}^a`$ and Eqs. (44-45) reduce to the standard shot noise correction results for $`\alpha =0`$, $`\widehat{P}=P\overline{n}^1`$, $`\widehat{B}=B\overline{n}^1_i\widehat{P}_i\overline{n}^2`$ (Peebles 1980). The estimator of the reduced bispectrum is thus
$$\widehat{Q}_{123}=\frac{\widehat{B}_{123}}{\widehat{P}_1\widehat{P}_2+\widehat{P}_2\widehat{P}_3+\widehat{P}_3\widehat{P}_1}.$$
(46)
Note that this involves a non-linear combination of estimators, so it is not necessarily and unbiased estimator of $`Q`$, as we shall discuss in the next section.
## 7 Effects of Survey Geometry and Sampling
### 7.1 Power Spectrum
Using Eq. (44) for the power spectrum, we find the recovered power spectrum shown in symbols in Fig. 6. The mean is obtained from 785 1.2Jy mock surveys, and the error bars correspond to the errors scaled to a single realization. The solid line shows the actual power spectrum in redshift space obtained from 200 2LPT realizations of $`\mathrm{\Lambda }`$CDM, $`\sigma _8=0.7`$ (error bars on these measurements are supressed for clarity, they are completely negligible). We see from Fig. 6 that the convolution with the window of the survey $`P_G(k)`$ in Eq. (43) affects the recovery of the density power spectrum at scales $`k<0.05`$ h/Mpc. For $`k>0.05`$ h/Mpc, which are the scales we use in this work, the recovered power spectrum in the narrow-window approximation agrees very well with the correct answer.
Figure 7 shows a comparison of the error bars as a function of scale in the ideal case (2LPT) with those in different IRAS surveys. Note that the horizontal scale is now linear, to show better the range of scales we actually use in this paper, and where the narrow-window approximation works. The 2LPT realizations are the same as in the previous figure, having $`256^3`$ particles in a 1200 Mpc/h box. The Gaussian prediction shown in Fig. 7 is obtained by the simple relation $`\mathrm{\Delta }P(k)/P(k)=2/N(k)`$, where $`N(k)`$ is the number of modes within a given shell in Fourier space centered at $`k`$. From this comparison, we see that even at $`k=0.3`$ h/Mpc, the Gaussian approximation to the errors works very well, consistent with the results from N-body simulations in Meiksin & White (1999) and Scoccimarro, Zaldarriaga & Hui (1999).
Also shown in Fig. 7 are the corresponding errors for QDOT, 1.2Jy and PSCz surveys (those of 2Jy are intermediate between QDOT and 1.2Jy and are supressed for clarity). All the solid lines correspond to using the FKP weighing $`P_0=2000`$, constant for all wavenumbers. The dashed lines show the resulting error bars in PSCz mock surveys when $`P_0=8000`$. As expected, the error bars in the surveys are considerably larger than in the 2LPT case, due to the smaller volume and the larger shot noise. In fact, when the latter dominates, errors stop decreasing as $`k`$ is increased, which happens first for QDOT, then 1.2Jy and last for PSCz, as expected. The errors in the PSCz case are smaller for $`P_0=8000`$ at low $`k`$ than for $`P_0=2000`$ as expected, if one wants to measure longer wavelengths it is better to weight distant galaxies more (higher $`P_0`$), as dictated by the FKP prescription.
The survey geometry and sparse sampling (shot noise) not only increases the power spectrum error bars at a given $`k`$, but also leads to correlations between different band powers. To quantify this, Fig. 8 shows the power spectrum correlation coefficient, obtained from the power spectrum covariance matrix, $`C_{ij}(\widehat{P}_iP_i)(\widehat{P}_jP_j)`$, by $`r_{ij}C_{ij}/\sqrt{C_{ii}C_{jj}}`$, for fixed $`k_i`$ as a function of $`k_j`$. The figure shows $`k_i=0.053,0.134,0.201`$ h/Mpc for different IRAS surveys and the “ideal” 2LPT case. In the latter (bottom left panel) we see that $`r_{ij}`$ is practically diagonal, no significant correlations are induced at these scales by non-linear evolution or the radial nature of redshift distortions. On the other hand, the 2Jy shows the slowest decay of $`r_{ij}`$, as expected due to its small volume. At the large $`k`$ end, $`r_{ij}`$ becomes roughly constant, when shot noise is expected to dominate the covariance between modes. All the mock survey results in this figure assume FKP weighing $`P_0=2000`$, except for the dashed line in the PSCz case (bottom right panel), which corresponds to $`P_0=8000`$. The larger value of $`P_0`$ in this case means that we are effectively increasing the volume of the survey, and thus the amount of cross-correlation between different band powers decreases. Thus, given the choice between $`P_0=2000`$ and $`P_0=8000`$ at $`k0.1`$ h/Mpc where both weights give similar error bars (see Fig. 7), the choice $`P_0=8000`$ should be preferred to decrease cross-correlations (which the FKP method does not attempt to minimize). In general, cross-correlation between different band powers are not negligible and must be taken into account when extracting cosmological information from power spectrum measurements (Eisenstein & Zaldarriaga 1999; Hamilton 2000; Hamilton & Tegmark 2000).
When estimating the power spectrum using likelihood analysis, an assumption must be made about the likelihood function. At large scales, the density field can be assumed to be Gaussian, and the power spectrum appears in the covariance matrix of the Fourier modes. At scales where non-linear effects become important, $`k0.20.3`$ h/Mpc, the density field cannot be assumed to be Gaussian anymore. A complementary approach at these scales is to consider that the power spectrum distribution about its mean is Gaussian; if there are enough independent modes contributing to a given shell in Fourier space, by the central limit theorem the distribution of the power spectrum should approach Gaussianity. However, due to the finite volume of the survey and shot noise, different modes are correlated, so it is not clear a priori at what scales Gaussianity applies. Figure 9 shows the power spectrum probability distribution function (PDF) obtained from 785 realizations of the 1.2Jy survey. The PDF is plotted as a function of the standarized variable $`\delta P/\mathrm{\Delta }P`$, where $`\delta P\widehat{P}P`$, $`\widehat{P}=P`$, and $`(\mathrm{\Delta }P)^2(\widehat{P}P)^2`$. The four curves show the average power spectrum PDF for $`k<0.067`$, $`0.067<k<0.134`$, $`0.134<k<0.201`$, and $`0.201<k<0.268`$ (all in units of h/Mpc). The solid line denotes a Gaussian distribution, the actual distribution of power is approximately a $`\chi ^2`$ distribution (with a number of degrees given by the effective number of modes one would obtain from Fig. 7 by assuming $`\mathrm{\Delta }P/P2/N_{\mathrm{eff}}(k)`$) and becomes closer to Gaussian at small scales, as expected. We see though that at $`k=0.268`$ h/Mpc there are still noticeable deviations from Gaussianity.
### 7.2 Bispectrum
#### 7.2.1 Choice of Triangle Variables
We now turn to a discussion of the bispectrum in mock surveys, focusing on the same issues illustrated by the power spectrum results discussed so far. For each realization of a given survey, the bispectrum is measured for all triangles with sides between $`k=0.05`$ h/Mpc and $`k=0.2`$ h/Mpc. The lower bound is imposed by the scale where the narrow-window approximation works well, whereas the upper bound is where shot noise can be best corrected for in a single realization, still well within the limits of applicability of 2LPT. The triangles are binned in terms of their sides, $`k_1k_2k_3`$, in steps of $`k_f`$, i.e. $`k_i=10,11,12,\mathrm{}`$, where $`k_f`$ is the “fundamental” mode of the Fourier space grid we use, $`k_f=2\pi /12000.005`$ h/Mpc. There are 4741 such triangles in total, each of these triangles in turn contains from $`10^4`$ to $`4\times 10^6`$ “elementary” triangles (those involving the product of 3 Fourier coefficients).
Since the window function of the survey has a width of order $`\mathrm{\Delta }k/k_f5`$ (see Fig. 6), triangles are very correlated if their sides differ by less than this width, as a result the number of “independent” triangles is of course much smaller than 4741. We therefore use coarser bins, defined in terms of shape $`s`$, ratio $`r`$, and scale $`k`$ parameters (Peebles 1980),
$$s\frac{k_1k_2}{k_3},r\frac{k_2}{k_3},kk_3,$$
(47)
where the “shape” parameter $`s`$ obeys $`0s1`$ ($`k_1k_2k_3`$), the ratio $`1r4`$, and the overall scale of the triangle satisfies $`10k/k_f40`$. The shape parameter is $`s=0`$ for isosceles triangles (and if $`r=1`$ these are equilateral triangles); for $`s=1`$ we have colinear triangles. Using 10 bins in each variable plus the closed triangle constraint yields 203 triangles.
One advantage of using these variables is that the main dependence of $`Q`$ is on the shape parameter $`s`$, with a small dependence on the ratio parameter $`r`$, and for a scale-free initial power spectrum no dependence on $`k`$. Therefore, if we plot $`Q`$ as a function of a single variable (“triangle”) which depends on ordered triplets $`s,k,r`$ as shown in the bottom panel of Fig. 10, $`Q`$ is roughly an increasing function of triangle number, which runs from 1 to 203. The order within the triplet is chosen so $`s`$ has the slowest variation, which means that low triangle number corresponds to nearly equilateral triangles, whereas high triangle number corresponds to nearly colinear triangles. For example, triangle number $`110`$ corresponds to equilateral triangles of scales $`k/k_f=10,\mathrm{},40`$. The top panel in Fig. 10 shows the reduced bispectrum $`Q`$ as a function of triangle number for 2LPT realizations in redshift-space of $`\mathrm{\Lambda }`$CDM, with $`\sigma _8=0.7`$. This scheme of using a single (somewhat complicated) variable that parametrizes $`Q`$ will be useful in the following since it allows to deal with all the data at once (i.e. plot all the triangles of different shapes and scales in terms of a single variable).
#### 7.2.2 Estimator Bias
The effects of finite volume on higher-order statistics has been well studied for one-point moments, in both N-body simulations (Colombi, Bouchet & Schaeffer 1994, 1995; Colombi, Bouchet & Hernquist 1996; Munshi et al. 1999) and galaxy surveys (Szapudi & Colombi 1996; Colombi, Szapudi & Szalay 1998; Kim & Strauss 1998; Hui & Gaztañaga 1999; Szapudi, Colombi & Bernardeau 1999). Here we consider the same problem for three-point statistics, i.e. the bispectrum.
Figure 11 shows a comparison of the amplitude of the reduced bispectrum $`Q`$ as a function of triangle for different surveys and the underlying bispectrum (labeled 2LPT), to illustrate the effects of finite survey volume and sparse sampling. The top panel shows a comparison of 2LPT with 2Jy mock surveys, we see that the effects of finite volume (2Jy being the most shallow of the IRAS surveys) are rather strong: colinear configurations can be underestimated by factors as large as $`70\%`$. On the other hand, for the deeper PSCz survey (bottom panel) the finite volume effect is as expected substantially smaller, at about $`20\%`$. The middle panel compares the QDOT and PSCz surveys, the former being a sparse sample (one in six galaxies) of the latter. As we see, apart from a significant increase in the error bars, sparse sampling does not significantly affect the amplitude of the reduced bispectrum $`Q`$ (provided of course that the discretness corrections in Eqs. (44-45) can be performed accurately), except perhaps by factors of order $`10\%`$ at colinear configurations.
The finite volume effect is particularly worrisome for the determination of galaxy bias parameters from surveys, as it has a similar systematic effect on $`Q`$ as galaxy biasing with $`b_1>1`$ and $`b_2>0`$. For example, if this effect is ignored, one would conclude erroneously from the unbiased realizations (i.e. $`b_1=1`$, $`b_2=0`$) in Fig. 11 that $`b_1=2.23,1.35`$ and $`b_2=2.18,0.42`$ for 2Jy and PSCz respectively. This “estimator bias”, the fact that the estimator of $`Q`$ is not unbiased, arises because $`\widehat{Q}_{123}`$ in Eq. (46) is a non-linear combination of estimators, which are nonetheless unbiased (provided that we use the mean value of the selection function, rather than its actual value to substract the mean density, and thus avoid the integral constraint bias). There are two sources of estimator bias in Eq. (46) corresponding to the two non-linear operations involved; the quadratic combination in the denominator of Eq. (46) and the ratio between numerator and denominator. Detailed analysis shows that both contributions turn out to be important, in particular, the quadratic combination overestimates the underlying value (obtained from 2LPT) for colinear configurations (thus leading to an understimate of $`Q`$ for these triangles), whereas taking the ratio also leads to an underestimate of $`Q`$ for colinear configurations. In addition, there is a rather small overestimate of $`Q`$ for equilateral configurations (only apparent in the least noisy PSCz mock surveys, bottom panel of Fig. 11). This however can be traced to a small estimator bias in the bispectrum for equilateral triangles at large scales, which indicates that the narrow window approximation used to deconvolve with the survey window bispectrum in Eq. (45) is beginning to break down. However, the effect of this particular estimation bias is in practice negligible when compared to the other estimation biases involved.
#### 7.2.3 Error Bars and Triangle Cross-Correlations
Figure 12 shows the relative error on the reduced bispectrum, $`\mathrm{\Delta }Q/Q`$ as a function of triangle for different surveys, where $`(\mathrm{\Delta }Q)^2(Q\overline{Q})^2`$. The top panel shows the ideal 2LPT case, and the errors expected in a single realization of the 2Jy survey. Similarly, the middle and bottom panels show the corresponding results for 1.2Jy and PSCz surveys. We see that in general errors decrease as configurations become more colinear (simply because the number of triangles increases), with spikes due to shot noise at small scales. For example, the spike at the tenth triangle is due to equilateral triangles at small scales ($`k=0.2`$ h/Mpc) becoming dominated by shot noise. For QDOT (not shown), $`\mathrm{\Delta }Q/Q>1`$ for all triangles, with most triangles having significantly larger errors. In Fig. 13, we show similar results for $`\chi ^2`$ initial conditions. The 2LPT results are from Scoccimarro (2000), and the 1.2Jy results show the rather strong finite volume effects (top panel) and the increase error bars with respect to the Gaussian initial conditions case (compare with middle panel in Fig. 12). These results are noiser than in the Gaussian case shown in Figs. 11-12 because only 200 (instead of $`10^3`$) mock catalogs are used to compute $`Q`$ and $`\mathrm{\Delta }Q`$.
Figure 14 shows the bispectrum correlation coefficient, obtained from the covariance matrix as in the power spectrum case. This measures correlations between different triangles and shows that in the ideal case (2LPT, top panel) different triangles are practically independent, whereas when survey geometry and sampling are included (1.2Jy, bottom panel), correlations between different triangles are significant. The peaks in the correlation coefficient can be matched with triangles that have their sides of similar length, as expected from the power spectrum correlations shown in Fig. 8.
#### 7.2.4 The Distribution of $`Q`$
The main ingredient to implement a likelihood analysis is to understand the probability distribution (PDF) of the reduced bispectrum $`Q`$. In order to do that, we use the 2LPT Monte Carlo realizations to construct the PDF’s for different mock surveys. As found for the power spectrum (Fig. 9), we expect that bispectrum correlations due to survey geometry and sampling as shown in Fig. 14 invalidate straightforward application of the central limit theorem, and thus convergence to a Gaussian distribution for $`Q`$.
Figure 15 shows the results. We found that the distributions of $`Q`$ for all the triangles we consider were very similar when plotted in terms of the scaled variable $`\delta Q/\mathrm{\Delta }Q`$, where $`\delta QQ\overline{Q}`$. We have thus averaged all these PDF’s to increase the signal to noise. In reality one expects these PDF’s to approach Gaussianity when the error is small (i.e. survey geometry and sampling are not significant), and thus a variation of PDF on scale; however for the range of scales we consider the errors don’t change significantly (see Fig. 12), so to first approximation the PDF’s do not change.
The dotted lines in the top panel of Fig. 15 shows indeed that in the ideal case (2LPT) the PDF of $`Q`$ approaches Gaussianity (solid smooth line). On the other hand, mock surveys show clear deviations from Gaussianity, with positive skewness, exponential tails for 1.2Jy and 2Jy surveys, and in the sparse sampling case (QDOT) power-law tails and rather strong kurtosis. For $`\chi ^2`$ non-Gaussian initial conditions (bottom panel), the non-Gaussianity is even more pronounced, as expected. This clearly shows that likelihood analysis for current surveys based on a Gaussian likelihood for the bispectrum are not justified. In particular, the skewness of the PDF can lead to a “statistical” bias in the estimation of galaxy bias: since the most likely value for $`Q`$ is to underestimate the mean, use of a Gaussian likelihood will infer bias parameters that overestimate $`b_1`$ and understimate $`b_2`$.
## 8 Likelihood Analysis
### 8.1 $`Q`$ Eigenmodes
A general treatment of likelihood non-Gaussianity and cross-correlations between triangles is very complicated. Here we will make the assumption that the eigenmodes of $`Q`$, defined by diagonalizing its covariance matrix, are effectively independent (which would be true if their joint PDF were Gaussian), and thus write the full likelihood as a product of eigenmode-likelihoods (which are computed by 2LPT realizations of a given survey, and are generally non-Gaussian). In general, diagonalizing the covariance matrix does not guarantee that the eigenmodes are independent (third and higher-order correlations could be still non-zero), but as we shall see this simplification seems to work very well in practice.
Given a set of measured reduced bispectrum amplitudes $`\{Q_m\}`$, $`m=1,\mathrm{},N_T`$, where $`N_T`$ is the number of closed triangles in the survey, we diagonalize their covariance matrix so the $`Q`$-eigenmodes $`\widehat{q}_n`$,
$$\widehat{q}_n=\underset{m=1}{\overset{N_T}{}}\gamma _{mn}\frac{Q_m\overline{Q}_m}{\mathrm{\Delta }Q_m},$$
(48)
satisfy
$$\widehat{q}_n\widehat{q}_m=\lambda _n^2\delta _{nm},$$
(49)
where $`\overline{Q}Q`$ and $`(\mathrm{\Delta }Q)^2Q^2\overline{Q}^2`$. These $`Q`$-eigenmodes have “signal to noise” ratio $`S/N`$,
$$\left(\frac{S}{N}\right)_n\frac{1}{\lambda _n}\left|\underset{m=1}{\overset{N_T}{}}\gamma _{mn}\frac{\overline{Q}_m}{\mathrm{\Delta }Q_m}\right|.$$
(50)
The physical interpretation of $`Q`$-eigenmodes becomes clear when ordered in terms of their signal to noise, as shown in the top panel in Fig. 16. The best eigenmode (highest signal to noise, $`S/N3`$), say $`n=1`$, corresponds to all weights $`\gamma _{m1}>0`$; that is, it represents the overall amplitude of the bispectrum (Fig. 16, middle panel). The next $`Q`$-eigenmode, $`n=2`$ with $`S/N1`$, has $`\gamma _{m2}>0`$ for nearly colinear triangles and $`\gamma _{m2}<0`$ for nearly equilateral triangles (Fig. 16, bottom panel); that is, it represents the configuration dependence of the bispectrum. Higher-order eigenmodes contain further information such as variations of amplitude and shape with scale, but generally have $`S/N<1`$ in IRAS catalogs (although there are many of them). The $`S/N`$ of eigenmodes for a particular survey is sensitive mostly to the biasing scheme and the amplitude of fluctuations $`\sigma _8`$, stronger clustering leading to larger $`S/N`$ overall.
### 8.2 Recovering Bias Parameters
As discussed above, we can write down the likelihood as a function of the bias parameters,
$$(\alpha _1,\alpha _2)\underset{i=1}{\overset{N_T}{}}P_i[\nu _i(\alpha _1,\alpha _2)],$$
(51)
where $`\alpha _11/b_1`$, $`\alpha _2b_2/b_1^2`$, and
$$\nu _i(\alpha _1,\alpha _2)\frac{1}{\lambda _i}\underset{j=1}{\overset{N_T}{}}\gamma _{ji}\frac{Q_j(\alpha _1\overline{Q}_j+\alpha _2)}{\mathrm{\Delta }Q_j},$$
(52)
where the $`Q_j`$ ($`j=1,\mathrm{},N_T`$) are the data, and the mean $`\overline{Q}_j`$, standard deviation $`\mathrm{\Delta }Q_j`$, and the non-Gaussian PDF’s $`P_i(\nu _i)`$ are extracted from the mock catalogs. Essentially, the $`Q`$-eigenmode PDF’s look similar to those shown in Fig. 15 for $`n=1`$ (for which $`\gamma _{m1}>0`$), and a symmetrized version of it for other eigenmodes where $`\gamma _{mn}`$ takes positive and negative values equally frequent. We assume that to first approximation, $`\mathrm{\Delta }Q_j`$ does not depend on biasing, which would be true if biasing is linear and deterministic. Thus, in the following all results shown ignore a possible dependence of $`\mathrm{\Delta }Q_j`$ on non-linear biasing.
Figure 17 shows an example of the use of Eq. (51) to recover bias parameters. We use parent 2LPT realizations of $`200^3`$ particles in a 600 Mpc/h box, which are then used to select about $`2\times 10^6`$ “galaxies” according to the prescription in Eq. (14). These realizations are then used to generate mock catalogs of the 1.2Jy survey. We generated 200 mock catalogs for each of the biasing schemes shown in Fig. 1, and then measure $`Q`$ in redshift space. Three of the models are generated from $`\mathrm{\Lambda }`$CDM ($`\sigma _8=0.7`$, $`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$) and a fourth model with same power spectrum shape and normalization but $`\mathrm{\Omega }_m=1`$ (top right panel). The models in the right panels have about the same $`\beta \mathrm{\Omega }_m^{0.6}/b_10.65`$. We now describe the results of analyzing these biased samples with the likelihood in Eq. (51) where all the quantities ($`P_i,\overline{Q}_j,\mathrm{\Delta }Q_j`$) are evaluated from the unbiased ($`b_1=1`$, $`b_2=0`$) realizations of the surveys described above to test the assumptions behind our likelihood analysis.
The plots in Fig. 17 show typical outcomes of the maximum likelihood procedure; triangles denote the “correct” biasing parameters shown in Fig. 1 (obtained by measuring the bispectrum in the parent 2LPT realization). The best fit parameters and error bars from these particular single realizations of biased catalogs are shown in Table 6. Note that, as expected, the derived error bars increase for the antibiased samples. It is also important to note that although the bispectrum can break the degeneracy between models with the same $`\beta `$ but different bias and $`\mathrm{\Omega }_m`$, for the 1.2Jy survey the errors are large enough that the two models considered overlap to within $`68\%`$. These results are only meant to illustrate the expected error bars and do not constitute a test of the likelihood analysis (since one can always find a single realization out of 200 which shows reasonable reconstruction of the biasing parameters). To test the accuracy of the likelihood analysis, we checked whether $`68\%`$ likelihood ellipses enclosed the correct answer in $`68\%`$ of the realizations. We found that in practice this number is about $`70\%`$; if we instead use a Gaussian likelihood (including correlations) we found that this number drops to about $`45\%`$. Correlations between triangles account for a significant area of the ellipses in Fig. 17, neglecting correlations leads to error ellipses that contract by a factor of order 10 in area.
The use of maximum likelihood (ML) to estimate the best fit bias parameters is not guaranteed to give statistically unbiased results, since the likelihood is not Gaussian. To test this, Fig. 18 shows the ML estimates as a function of the cumulative number of mock catalog realizations used, with panels ordered as in Fig. 1. The deviations observed at low numbers of catalogs is due to the effect discussed above, i.e. one expects individual realizations which are off the underlying value (error bars for each curve in Fig. 18 are avoided for reasons of clarity, but can be estimated from the ones in Table 6 by scaling by $`1/\sqrt{N_{\mathrm{cat}}}`$). We see that as the number of catalogs $`N_{\mathrm{cat}}`$ becomes large (and the error bars become small), the results generally approach the expected asymptotic value of $`1/b_1`$, $`b_2/b_1^2`$ (shown as horizontal marks in Fig. 18; see also Fig. 1 and Table 6). However, there are detectable (small) estimation biases, in particular for the non-linear bias parameter. For the lower left panel, the resulting parameters actually make the agreement in Fig. 1 better. Since these estimation biases are about $`34`$ times smaller than the expected error bars for a single survey, we can safely ignore this problem, as it appears not to be systematic. For larger surveys, where the errors are expected to be smaller it is also expected that the likelihood will be better approximated by a Gaussian, thus the estimation bias from ML is also likely to be smaller, so it is not clear this is a serious problem for constraining bias parameters, but it is worth keeping in mind.
### 8.3 Comparison with Previous Work
A likelihood analysis for obtaining bias parameters from bispectrum measurements has been pioneered by Matarrese et al. (1997; MVH), and extended to redshift-space by Verde et al. (1998). Their approach differs from ours in several ways, as they intended their treatment to next-generation redshift surveys where some of issues faced in this work will not be as serious (although this will have to be checked in each particular case). First, they assume a Gaussian likelihood for the bispectrum, which as discussed above breaks down in the case of IRAS surveys. This approximation will considerably improve for surveys such as 2dF and SDSS.
Second, we consider all triangle shapes with their full covariance matrix, rather than using only equilateral and colinear triangles to make the covariance approximately diagonal. In fact, for IRAS surveys Fig. 14 shows that triangles of similar shapes and scales are strongly correlated by the window of the survey and sparse sampling. In addition, discarding triangle shapes throws away information, which is undesirable. This is of course not a limitation of MVH’s method; on the other hand, analytic calculation of the covariance matrix between general triangle shapes including the effects of survey geometry and sampling quickly becomes complicated if one intends to go beyond the Gaussian approximation. In our case this is handled automatically by the numerical 2LPT realizations.
Our treatment of perturbation theory (PT) leads to some minor improvements. By being Lagrangian, our calculation includes loop corrections beyond leading order, although only approximately. Their treatment using second-order Eulerian PT is only consistent for the bispectrum, not its covariance (which has contributions from the sixth-point function and thus requires fifth-order Eulerian PT); however the leading order contribution is Gaussian, so to leading order (most of the contribution at the scales of interest) both treatments should agree. 2LPT is not exact beyond second-order in Eulerian space, however the results of Table 2 show that it gives a very good approximation to the higher-order moments.
The most important difference is the treatment of redshift distortions, our numerical calculations allow us to perform the redshift-space mapping exactly (rather than perturbatively), which makes quantitative difference even at large scales (e.g. solid versus dotted lines in Fig. 3). Moreover, by using non-linear dynamics rather than a phenomenological model to treat non-linear distortions, we avoid introducing auxiliary quantities such as the velocity dispersion parameter $`\sigma _v`$ (Verde et al. 1998; Scoccimarro et al. 1999).
Our calculations of the finite volume effects are in reasonable agreement with the results found for one-point higher-order statistics (Szapudi et al. 1999a, 1999b). In particular, we find that the PDF of the bispectrum is generally non-Gaussian with positive skewness, whereas Szapudi et al. (1999b) find that the PDF of the skewness parameter $`S_3`$ can be approximated by a suitably scaled lognormal. We also find that in general the estimator bias for $`Q`$ is smaller than the error $`\mathrm{\Delta }Q`$, although in the case of the 2Jy survey both quantities are comparable. It is thus important to correct for estimator bias (Hui & Gaztañaga 1999), even when it is smaller than the error, since the systematic shift affects the performance of the likelihood estimation; i.e. without this correction $`68\%`$ likelihood contours would not enclose the correct answer $`68\%`$ of the time. Our approach can be also applied to one-point statistics, and would be interesting to compare with the results by Szapudi et al. (1999a,1999b) since they involve different approximations.
## 9 Summary and Conclusions
We studied the use of the bispectrum to recover bias parameters and constrain non-Gaussian initial conditions in realistic redshift surveys. We considered the effects of stochastic non-linear biasing, radial redshift-distortions, survey geometry and sampling on the shape dependence of the bispectrum.
We found that bias stochasticity does not seem to affect the use of the bispectrum to recover the mean biasing relation between galaxies and mass, at least for models in which the scatter is uncorrelated at large scales. Clearly more work is necessary along these lines, we have just considered a very simple extension to the local deterministic bias model. We also found that the radial nature of redshift distortions changes the bispectrum monopole compared to the plane-parallel values by only a few percent, well below current error bars.
On the other hand, survey geometry leads to finite volume effects which must be taken into account in current surveys before comparison with theoretical predictions can be made. Similarly, sparse sampling and survey geometry correlate different triangles leading to a breakdown of the Gaussian likelihood approximation. We developed a likelihood analysis using bispectrum eigenmodes, calculated by Monte Carlo realizations of mock surveys generated with second-order Lagrangian perturbation theory (2LPT) and checked against N-body simulations. This can be used to derive quantitative constraints from current and future redshift surveys. In a companion paper (Scoccimarro et al. 2000) we apply the results obtained here to the analysis of the bispectrum in IRAS surveys.
Our treatment can be improved in several ways for applications to surveys that demand more accuracy. In particular, we relied on the simple relation in Eq. (3) to relate the galaxy bispectrum to the mass bispectrum; in practice, however, redshift-distortions and biasing do not conmute, and this relation is more complicated (Scoccimarro, Couchman & Frieman 1999). Related to this issue is a formulation, so far lacking, of a joint likelihood for the power spectrum and bispectrum, which involves dependence on $`\mathrm{\Omega }_m`$ and biasing parameters. In addition, we have only considered closed triangles in Fourier space. In practice, due to breakdown of translation invariance by the survey geometry, selection function, and radial redshift distortions, open triangles also contain useful information. In order to do this, one possibility is to construct a general cubic estimator for the bispectrum, analogous to the quadratic estimator for the power spectrum (e.g. Tegmark et al. 1998), which seems a rather difficult task for non-Gaussian fields. However, future large redshift surveys such as 2dF and SDSS will be the ideal datasets to take full advantage of the cosmological information available from higher-order statistics, and thus the effort in developing more accurate statistical methods will likely be rewarded.
###### Acknowledgements.
I thank my collaborators on the IRAS bispectrum, H. Feldman, J.A. Frieman and J.N. Fry, for comments and discussions. I also would like to thank S. Colombi, E. Gaztañaga and L. Hui for conversations on finite volume effects, R. Crittenden for help with improving the speed of my bispectrum code, R. Sheth for discussions about galaxy biasing, and C. Murali, S. Prunet and especially M. Zaldarriaga for many helpful discussions regarding maximum likelihood estimation. The N-body simulations generated for this work were produced using the Hydra $`N`$-body code (Couchman, Thomas, & Pearce 1995). I thank H. Couchman and R. Thacker with help regarding the use of Hydra. This work was supported by endowment funds from the Institute for Advanced Study.
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# References
Status of the neutrino decay solution to the solar neutrino problem Sandhya Choubey<sup>a</sup> <sup>*</sup><sup>*</sup>*sandhya@tnp.saha.ernet.in, Srubabati Goswami<sup>a</sup> sruba@tnp.saha.ernet.in, Debasish Majumdar<sup>b</sup> debasish@tnp.saha.ernet.in
$`^b`$Present Address: Saha Institute of Nuclear Physics, Calcutta, INDIA.
<sup>a</sup>Saha Institute of Nuclear Physics,
1/AF, Bidhannagar, Calcutta 700 064, INDIA.
<sup>b</sup>92 Acharya Prafulla Chandra Road, Calcutta 700 009, INDIA.
Abstract
We re-examine the neutrino decay solution to the solar neutrino problem in the light of the SuperKamiokande (SK) data. For the decay solution the SK spectrum data by its own can provide a fit comparable to the fit obtained from the MSW solution. However when one combines the results from the total rates of the $`{}_{}{}^{37}Cl`$ and $`{}_{}{}^{71}Ga`$ experiments the fit becomes much poorer.
PACS numbers: 14.60.Pq, 26.60.+t, 96.60.Tv, 13.35.Hb
keywords: neutrino mass, solar neutrinos, decay, mixing
In this paper we examine the status of the neutrino decay solution to the solar neutrino problem in the light of the SK data <sup>1</sup><sup>1</sup>1The possibility of neutrino decay as a solution to the atmospheric neutrino problem has been considered recently in . Neutrino decay as a solution to the solar neutrino problem has been considered earlier (pre-SK) . However since the Ga data implies that the low energy pp neutrinos are less suppressed compared to the high energy $`{}_{}{}^{8}B`$ suppression seen in Cl or the Kamiokande experiments, this solution was ruled out at 99% C.L. since the decay term $`\mathrm{exp}(\alpha L/E)`$, (where $`\alpha =m_2/\tau _o`$, $`m_2`$ being the mass of the unstable state and $`\tau _o`$ its rest frame lifetime), suppresses the low energy neutrino flux more than the high energy flux. However the SK spectrum data shows more events for the high energy bins. Though the statistics in these high energy bins need improvement, we explore the status of the neutrino decay scenario in the context of this behaviour of the SK spectrum data.
Radiative decays of neutrinos are severely constrained and thus we are interested in the non-radiative decay modes. Two classes of models have been considered in the literature in this connection.
1. In neutrinos are considered to be Dirac particles. They consider the decay mode $`\nu _2\overline{\nu }_{1R}+\varphi `$, where $`\overline{\nu }_{1R}`$ is a right handed singlet and $`\varphi `$ is an iso-singlet scalar. Thus all the final state particles are sterile.
2. In neutrinos are assumed to be Majorana particles and the decay mode is $`\nu _2\overline{\nu }_1+J`$, where $`\overline{\nu }_1`$ interacts as a $`\overline{\nu }_e`$ with a probability $`|U_{e1}|^2`$ and J is a Majoron.
We work with just two flavors for simplicity and assume that the state $`\nu _2`$ is unstable, which decays with a rest frame lifetime $`\tau _o`$. The other neutrino mass states have lifetimes much greater than the sun-earth transit time and hence can be taken as stable. In presence of decay,
$`P_{\nu _e\nu _e}`$ $`=`$ $`(1|U_{e2}|^2)^2+|U_{e2}|^4\mathrm{exp}(4\pi L/\lambda _d)`$ (1)
$`+2|U_{e2}|^2(1|U_{e2}|^2)\mathrm{exp}(2\pi L/\lambda _d)\mathrm{cos}(2\pi L/\lambda _{osc})`$
$`P_{\nu _e\nu _x}`$ $`=`$ $`|U_{e2}|^2(1|U_{e2}|^2)\{1+\mathrm{exp}(4\pi L/\lambda _d)`$ (2)
$`2\mathrm{exp}(2\pi L/\lambda _d)\mathrm{cos}(2\pi L/\lambda _{osc})\}`$
where $`x`$ can be either $`\mu `$ or $`\tau `$ in the two flavour framework. $`\lambda _d`$ is the decay length defined as,
$$\lambda _d=2.5\times 10^3km\frac{E}{MeV}\frac{eV^2}{\alpha }$$
(3)
$`\lambda _{osc}`$ is the oscillation wavelength defined as
$$\lambda _{osc}=2.5\times 10^3km\frac{E}{MeV}\frac{eV^2}{\mathrm{\Delta }m^2}$$
(4)
$`L=R(t)(1r/R(t))`$, r being the distance of the point of neutrino production from the center of the sun and $`R(t)`$ is the sun-earth distance given by,
$$R(t)=R_0\left[1ϵ\mathrm{cos}(2\pi \frac{t}{T})\right]$$
(5)
Here, $`R_0=1.49\times 10^{13}`$ cm is the mean sun-earth distance and $`ϵ=0.0167`$ is the ellipticity of the earth’s orbit, $`t`$ is the time of the year and $`T`$ is 1 year. The $`\mathrm{\Delta }m^2`$ dependent oscillations are important around $`10^{10}10^{11}`$ eV<sup>2</sup>. We assume $`\mathrm{\Delta }m^2`$ to be much higher so that the cosine term averages out to zero. We further assume $`U_{e1}=\mathrm{cos}\theta `$ and $`U_{e2}=\mathrm{sin}\theta `$ so that the survival probability is
$$P_{\nu _e\nu _e}=\mathrm{cos}^4\theta +\mathrm{sin}^4\theta \mathrm{exp}(4\pi L/\lambda _d)$$
(6)
In Fig. 1 we plot $`P_{\nu _e\nu _e}`$ as a function of $`\alpha `$ for $`\mathrm{sin}^2\theta =0.51`$ for two illustrative values of $`\nu _e`$ energies 7 and 13 MeV. For $`\alpha <10^{13}eV^2`$, the $`\mathrm{exp}(4\pi L/\lambda _d)1`$ signifying very little decay. Beyond $`\alpha =10^{10}eV^2`$ the $`\mathrm{exp}(4\pi L/\lambda _d)0`$ which is the case where we have very fast decay. In both these regions the probabilities are energy independent. In the region where $`\alpha `$ is around $`10^{11}eV^2`$ the probability does depend on the energy.
The details of the solar neutrino code employed for performing the $`\chi ^2`$-analysis is described in . We use the BP98 solar model as the reference SSM. We perform $`\chi ^2`$-analyses
* using the total rates from $`{}_{}{}^{37}Cl`$, $`{}_{}{}^{71}Ga`$ and SK experiments. We incorporate the theory errors and their correlations.
* using the 825 day SK spectrum data including the uncorrelated as well as the energy-bin correlated errors.
* of the combined rates and SK spectrum data
For the last two cases we vary $`X_B`$ – the normalization of the $`{}_{}{}^{8}B`$ flux with respect to the SSM value as free parameter and determine its best-fit value.
From the fact that the neutrinos from SN1987A have not decayed on their way one gets a lower bound on the electron neutrino lifetime as $`\tau _0>5.7\times 10^5(m_{\nu _e}`$/eV) sec. However if one includes neutrino mixing then shorter lifetimes are allowed provided $`|U_{e2}|^2<0.81`$ . To be consistent with this, in our analysis we vary $`\mathrm{sin}^2\theta `$ in the range 0 to 0.8.
The data used for the total rates is presented in Table 1. Since SK has much higher statistics, we have not included the Kamiokande results. The results obtained for Model 1 are summarised in Table 2. For the Model 2 $`\nu _2`$ decays to $`\overline{\nu }_1`$, which interacts as $`\overline{\nu _e}`$ with a probabitity $`U_{e1}^2`$ and $`\overline{\nu }_x`$ with a probability $`U_{e2}^2`$. We present in Table 3 the results for Model 2, where we have taken into account the $`\overline{\nu _e}e`$ and $`\overline{\nu }_xe`$ scattering in addition to the $`\nu _ee`$ scattering in SK, while for the radiochemical experiments there is no change. Since the $`\overline{nu_1}`$ is degraded in energy the effect is not significant.
| Experiment | Chlorine | Gallium | Super-Kamiokande |
| --- | --- | --- | --- |
| $`\frac{\mathrm{Observed}\mathrm{Rate}}{\mathrm{BP98}\mathrm{Prediction}}`$ | $`0.33\pm 0.029`$ | $`0.562\pm 0.043`$ | $`0.471\pm 0.015`$ |
Table 1: The ratio of the observed solar neutrino rates to the corresponding BP98 SSM predictions used in this analysis. The results are from refs. and . For Gallium, the weighted average of the SAGE and Gallex results has been used.
| Parameters | Rates | Spectrum | Rates+Spectrum |
| --- | --- | --- | --- |
| | (d.o.f = 1) | (d.o.f = 15) | (d.o.f = 18) |
| $`\alpha `$ (eV<sup>2</sup>) | 0 | 4.22 $`\times 10^{11}`$ | 3.29 $`\times 10^{12}`$ |
| $`\mathrm{sin}^2\theta `$ | 0.5 | 0.51 | 0.32 |
| $`X_B`$ | Fixed at SSM value | 1.5 | 0.72 |
| $`\chi _{min}^2`$ | 12.59 | 17.68 | 33.59 |
| g.o.f | 3.88 $`\times 10^2`$ % | 27.99% | 1.41% |
Table 2: The best-fit values of the parameters, the $`\chi _{\mathrm{min}}^2`$ and the goodness of fit (g.o.f) for Model 1.
| Parameters | Rates | Spectrum | Rates+Spectrum |
| --- | --- | --- | --- |
| | (d.o.f = 1) | (d.o.f = 15) | (d.o.f = 18) |
| $`\alpha `$ (eV<sup>2</sup>) | 0 | 3.3 $`\times 10^{11}`$ | 3.29 $`\times 10^{12}`$ |
| $`\mathrm{sin}^2\theta `$ | 0.5 | 0.53 | 0.32 |
| $`X_B`$ | Fixed at SSM value | 1.5 | 0.71 |
| $`\chi _{min}^2`$ | 11.98 | 17.62 | 33.12 |
| g.o.f | 5.38 $`\times 10^2`$ % | 28.32% | 1.61% |
Table 3: The best-fit values of the parameters, the $`\chi _{\mathrm{min}}^2`$ and the goodness of fit (g.o.f) for Model 2.
From the analysis of only rates data the best-fit comes in the region where From the analysis of only rates data the best-fit comes in the region where $`\alpha =0`$, implying that the neutrinos are stable. For the total rates, the low energy neutrino flux should be suppressed less which is in contradiction with the energy dependence of the exponential decay term. Thus the best fit comes in the region where the decay term goes to 1. Since the probability in this region is energy independent the quality of the fit is not good.
For the only spectrum analysis, the best-fit comes in the region where the $`\mathrm{exp}(\alpha L/E)`$ term is non-vanishing and the high energy neutrinos are suppressed less. Since high energy bins have more number of events, the fit is much better compared to the fit to the total rates and is comparable to the ones obtained for the MSW oscillation solution ($`\chi _{min}^2=17.62`$) . The vacuum oscillation solution gives a better fit . The best-fit values quoted in Table 2 are obtained with $`X_B`$ constrained to be $``$ 1.5. However if we remove the upper limit and allow $`X_B`$ to take arbitrary values a slightly better fit is obtained for very high values of $`X_B`$:
* $`\alpha =6.53\times 10^{11}`$ eV<sup>2</sup>, $`\mathrm{sin}^2\theta =0.8`$, $`X_B=6.42,\chi _{min}^2=17.22,g.o.f.=30.41\%`$
We note that the $`\alpha `$ does not change much. But at higher $`X_B`$ one needs a higher $`\theta `$. To understand these features, in fig. 2 we plot $`X_B`$ vs $`\alpha `$ for various fixed values of $`\mathrm{sin}^2\theta `$. For getting this curve we determine the $`X_B`$ that gives the minimum $`\chi ^2`$ at a particular $`\alpha `$ and then repeat this excersise for $`\alpha `$ varying in the range $`10^{14}10^8`$ eV<sup>2</sup>. The minimum $`\chi ^2`$ obtained at each point of the parameter space of fig. 2 is within the 90% C.L. limit of the global $`\chi ^2`$ minimum. There are three regions
1. For very low values of $`\alpha `$ (upto $`10^{13}`$ eV<sup>2</sup>) the exponential term is 1 and $`P_{\nu _e\nu _e}`$ = 1 - 0.5 $`\mathrm{sin}^22\theta `$ which is the average oscillation probability. This can vary from 1 to 0.5 depending on $`\theta `$. However in 11 out of 18 bins of the SK spectrum data, the rate<sub>obs</sub>/SSM $`<`$ 0.5 and this can be achieved only by keeping $`X_B<1.0`$.
2. In the range $`10^{13}10^{10}`$ eV<sup>2</sup> the exponential term contributes to the survival probability and it falls sharply with increasing $`\alpha `$. Therefore in this range $`X_B`$ rises sharply as $`\alpha `$ increases.
3. Beyond $`10^{10}`$ eV<sup>2</sup> the exponential term goes to zero and the probability is $`P_{\nu _e\nu _e}=cos^4\theta `$. In this zone one can achieve a probability $`<`$ 0.5 by adjusting $`\theta `$ only and $`X_B`$ does not play an important role. However if $`X_B`$ is allowed to vary then for smaller values of $`\theta `$ the $`X_B`$ needed is low and vice-versa.
In fig. 3 we plot the $`\chi ^2`$ for the SK spectrum data, as a function of one of the parameters, keeping the other two unconstrained. In fig. 3a the solid(dashed) line gives the variation of $`\chi ^2`$ with $`\alpha `$ keeping $`X_B`$ unconstrained (fixed at 1.0). If $`X_B`$ is 1 then lower values of $`\alpha `$ are not allowed as by varying $`\theta `$ alone one cannot get a probability less than 0.5. If $`X_B`$ is allowed to vary then low values of $`\alpha `$ are also admissible. Higher values of $`\alpha `$ are allowed for both the cases. The minimum of course comes in the region where $`\alpha 10^{11}`$ eV<sup>2</sup> for which the high energy neutrino events are less suppressed.
Fig. 3b gives the variation of $`\chi ^2`$ with $`\mathrm{sin}^2\theta `$ keeping $`X_B`$ unconstrained (solid line) and $`X_B=1`$ (dashed line). For both the curves $`\alpha `$ can take any value. For $`X_B`$ = 1 one cannot get a good fit in the low $`\alpha `$ region as discussed above and the fit goes to the high $`\alpha `$ region. In this zone the probability ($`\mathrm{cos}^4\theta `$) increases as we decrease $`\theta `$ and this puts a lower limit on the allowed value of $`\theta `$. This can be evaded if $`X_B`$ is allowed to vary, as by adjusting $`X_B`$ one can get a good fit even if $`\theta `$ is very low.
Fig. 3c gives the allowed range of $`X_B`$ keeping the other two parameters unconstrained. From the figure we see that values of $`X_B`$ below 0.4 are not allowed at 90% C.L.. As discussed, $`X_B`$ plays an important role in the low and intermediate $`\alpha `$ regions. In the latter zone the $`X_B`$ required is high, therefore the constraints on the low values of $`X_B`$ comes from the low $`\alpha `$ region. In this region as we decrease $`X_B`$ the number of events will decrease which can be adjusted by increasing the survival probability and the maximum value of $`P_{\nu _e\nu _e}=1.0`$ gives a lower limit on $`X_B`$. For very high values of $`X_B`$ the best-fit goes in the high $`\alpha `$ region where the probability is $`cos^4\theta `$. As we increase $`X_B`$ the number of events will increase which can be compensated by increasing $`\theta `$. However from SN1987A constraints $`\mathrm{sin}^2\theta `$ is restricted to be $`<`$ 0.8 and therefore beyond $`X_B=7.2`$ one does not get a good fit.
The rest frame lifetime obtained at the best-fit $`\alpha `$ from the spectrum data is $`\tau _o=(m_2/\mathrm{eV})1.54\times 10^5`$ sec. This is small enough for decay to happen before neutrino decoupling in the early universe. This can increase the effective number of light neutrinos, $`N_\nu `$, from 3. However depending on the data used the upper limit on $`N_\nu `$ can be 5 or 6 which is consistent with the model of neutrino decay used here.
The rest frame lifetime of $`\nu _2`$ is given by
$$\tau _0=\frac{16\pi }{g^2}\frac{m_2(1+m_1/m_2)^2}{\mathrm{\Delta }m^2\text{ }}$$
(7)
From the best-fit $`\alpha `$ $`10^{11}`$ eV<sup>2</sup> and assuming $`m_2>>m_1`$ one gets
$$g^2\mathrm{\Delta }m^2\text{ }16\pi 10^{11}$$
(8)
If we now incorporate the bound $`g^2<4.5\times 10^5`$ as obtained from K decay modes we obtain $`\mathrm{\Delta }m^2\text{ }>10^5`$ eV<sup>2</sup>, consistent with our assumption. At the best-fit $`\alpha `$ the decay length for atmospheric neutrinos is $`\lambda 2.5\times 10^{11}`$ km and hence they do not decay. However for $`\mathrm{\Delta }m^2\text{ }10^3`$ eV<sup>2</sup> and large mixing angles there will be substantial $`\nu _e\nu _x`$ conversion in conflict with the SK atmospheric neutrino and CHOOZ data .
In conclusion, neutrino decay solution to the solar neutrino problem is ruled out at 99.96% C.L. from the current data on total rates. For the SK spectrum data however, one can get much better fits, for $`\nu _2`$ lifetimes consistent with cosmological and supernova constraints. Although the best-fit for the spectrum data comes in the region where the probability is energy dependent, if $`X_B`$ is allowed to vary the $`\chi ^2`$ becomes flat over the entire range of parameter space. This implies that even energy independent suppression of the $`{}_{}{}^{8}B`$ flux is allowed. Even if $`X_B`$ is fixed at 1, the decay scenario can give $`\chi ^2`$ comparable to the best-fit in the energy independent high $`\alpha `$ regime. For the $`{}_{}{}^{8}B`$ flux normalisation factor a wide range 0.4 $`<X_B<7.2`$ is allowed at 90% C.L. just from the SK spectral data, if neutrino decay is operative. This is much broader than the range allowed by SSM uncertainties. For the rate+spectrum analysis the decay solution is disfavoured at more than 98% C.L. even if $`X_B`$ is allowed to take arbitrary values.
The authors thank A. Raychaudhuri for his involvement during the development of the solar neutrino code that has been used in this paper. They also wish to thank S. Pakvasa and K. Kar for discussions and the organisers of WHEPP-6 where this work was initiated. D.M. acknowledges financial support from the Eastern Center of Reasearch in Astrophysics, India.
Figure Captions
Fig. 1: The variation of the survival probability $`P_{\nu _e\nu _e}`$ (thick lines) and the exponential decay term $`\mathrm{exp}(\alpha `$ L/E) (thin lines) with $`\alpha `$, for two different values of neutrino energies. The solid curves correspond to neutrino energy = 7 MeV while the dotted curves are for neutrino energy = 13 MeV. For both the cases $`\mathrm{sin}^2\theta `$ is fixed at 0.51.
Fig. 2: The variation of $`X_B`$ with $`\alpha `$ for three different values of $`\mathrm{sin}^2\theta `$ shown in the plot. Each point on these curves is obtained by keeping $`\alpha `$ and $`\mathrm{sin}^2\theta `$ fixed and determining the $`X_B`$ corresponding to the minimum $`\chi ^2`$.
Fig. 3: The variation of $`\chi ^2`$ with (a) $`\alpha `$ for $`X_B`$ unconstrained (solid line) and fixed at 1.0 (dashed line), while $`\mathrm{sin}^2\theta `$ is kept unconstrained for both curves; (b) $`\mathrm{sin}^2\theta `$ keeping both $`\alpha `$ and $`X_B`$ unconstrained for the solid curve and with $`\alpha `$ unconstrained and $`X_B`$ fixed at 1.0 for the dashed curve; (c) $`X_B`$ keeping both $`\alpha `$ and $`\mathrm{sin}^2\theta `$ unconstrained. The dotted line shows the 90% C.L. limit for 3 parameters ($`\chi ^2=\chi _{min}^2+6.25`$) and the dash-dotted line gives the corresponding limit for 2 parameters ($`\chi ^2=\chi _{min}^2+4.61`$).
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# The Role of the QCD Vacuum in the Heavy-Quark Bound State Dynamics11footnote 1 Talk given by N.B. at the Fifth Workshop on Quantum Chromodynamics, Villefranche-sur-Mer, France, 3-7 January 2000.
## 1 Introduction
It is well known that the complex structure of the QCD vacuum can, in some regimes, be parametrized by local vacuum condensates, i.e. the expectation values of operators where all the perturbative vacuum fluctuations are taken out at the best of our ability . The first attempt to calculate the effect of non-perturbative vacuum condensates on the energy levels of heavy quarkonia was performed in . The leading contribution (in $`\alpha _\mathrm{s}`$ and $`\mathrm{\Lambda }_{\mathrm{QCD}}`$, the scale of non-perturbative physics) of the vacuum condensates to Coulombic $`n,l`$ heavy quarkonium states reads :
$$\delta E_{nl}^{\mathrm{V}\mathrm{L}}=m\frac{ϵ_nn^6\pi ^2G_2}{(mC_F\alpha _\mathrm{s})^4},$$
(1)
where $`ϵ_n`$ is a (known) number of order 1 and $`G_2(\alpha _\mathrm{s}/\pi )F_{\mu \nu }^a(0)F^{a\mu \nu }(0)`$ is the gluon condensate. The above expression for the corrections to Coulombic energy levels displays two relevant characteristics: 1) it is a correction of non-potential type, like the Lamb shift in QED; 2) it grows like $`n^6`$, thus being out of control for levels beyond the ground state. It was soon realized that the strong growth in $`n`$ could be corrected to some extent by considering non-local gluon condensates
$$G_2(x)\frac{\alpha _\mathrm{s}}{\pi }F_{\mu \nu }^a(x)\varphi (x,0)_{ab}^{\mathrm{adj}}F^{b\mu \nu }(0),$$
(2)
where $`\varphi (x,0)^{\mathrm{adj}}`$ is the Schwinger line (in the adjoint representation) connecting $`x`$ with $`0`$. These were understood as due to the presence of a fluctuating gluonic background with a characteristic time length $`T_g\mathrm{\Lambda }_{\mathrm{QCD}}^1`$. It is apparent that the local condensates stem from an expansion of the non-local ones in cases in which the correlation length is large with respect to the other physical scales of the system.
Heavy quarkonium is a non-relativistic bound system. Besides $`\mathrm{\Lambda }_{\mathrm{QCD}}`$, it is characterized by at least three hierarchically ordered scales, $`mv^2mvm`$, where $`m`$ is the mass of the heavy quark and $`v`$ its velocity. Therefore, it is only under the condition $`\mathrm{\Lambda }_{\mathrm{QCD}}mv^2`$ that the use of local condensates and, hence, the Voloshin–Leutwyler formula (1) can be justified. In the other regimes where heavy quarkonium states can sit, the non-perturbative dynamics will be encoded into more extended objects: non-local condensates and Wilson loop operators. This is the case for higher quarkonium levels ($`n>1`$).
In the following we will consider non-perturbative effects in heavy quarkonium systems in different kinematic regimes, in an effective field theory context. This approach has not only the considerable practical advantage of disentangling the different dynamical scales of the system, but also the conceptually relevant feature of disentangling perturbative effects from non-perturbative ones. This allows us to fully exploit the predictive power of QCD.
## 2 Local and non-local condensates: $`\mathrm{\Lambda }_{\mathrm{QCD}}mv`$
For the lower lying quarkonium states, it can be expected that the inverse of the typical size of the system is larger than $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. If this condition is fulfilled, both scales $`m`$ and $`mv`$ can be integrated out perturbatively from QCD, leading to an effective field theory where only the degrees of freedom of order $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ or $`mv^2`$ remain dynamical. This effective field theory is known as potential NRQCD, pNRQCD .
The infrared sensitivity of the quark–antiquark static potential at three loops signals that it may become sensitive to non-perturbative effects if the next relevant scale after $`mv`$ is $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. Indeed, in the situation $`mv\mathrm{\Lambda }_{\mathrm{QCD}}mv^2`$, the leading non-perturbative contribution (in $`\alpha _\mathrm{s}`$ and in the multipole expansion) to the static potential reads
$$V_0(r)^{\mathrm{non}\mathrm{pert}}=i\frac{g^2}{N_c}T_F\frac{r^2}{3}_0^{\mathrm{}}𝑑te^{itC_A\alpha _\mathrm{s}/(2r)}𝐄^a(t)\varphi (t,0)_{ab}^{\mathrm{adj}}𝐄^b(0)(\mu ),$$
(3)
where $`r`$ is the quark–antiquark distance. This term explicitly cancels, up to the considered order, the dependence of the perturbative static potential on the infrared scale $`\mu `$. It is interesting to note that the leading contribution in the $`\mathrm{\Lambda }_{\mathrm{QCD}}/mv^2`$ expansion of $`V^{\mathrm{non}\mathrm{pert}}`$ (obtained by putting the exponential equal to 1) cancels the order $`\mathrm{\Lambda }_{\mathrm{QCD}}^3r^2`$ renormalon that affects the static potential (the leading-order renormalon, of order $`\mathrm{\Lambda }_{\mathrm{QCD}}`$, cancels against the pole mass). Therefore, also in renormalon language, the above operator is the relevant non-perturbative contribution to the static potential in the considered kinematic situation.
If ΛQCD<
mv2subscriptΛQCD<
𝑚superscript𝑣2\Lambda_{\rm QCD}{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 6.0pt\vbox{\hbox{$\sim$}}}}\ }mv^{2} the static potential is purely perturbative and its explicit dependence on the infrared scale $`\mu `$ is reabsorbed in a physical observable by non-potential contributions. In the specific case of the quarkonium energy levels up to order $`\alpha _\mathrm{s}^5\mathrm{ln}\mu `$, these contributions are
$`\delta E_{n,l,j}=i{\displaystyle \frac{g^2}{3N_c}}T_F{\displaystyle _0^{\mathrm{}}}𝑑tn,l|𝐫e^{it(E_nH_o)}𝐫|n,l𝐄^a(t)\varphi (t,0)_{ab}^{\mathrm{adj}}𝐄^b(0)(\mu ),`$ (4)
$`E_nm{\displaystyle \frac{C_F^2\alpha _\mathrm{s}^2}{4n^2}},H_o{\displaystyle \frac{𝐩^2}{m}}+{\displaystyle \frac{1}{2N_c}}{\displaystyle \frac{\alpha _\mathrm{s}}{r}},`$
where $`|n,l`$ are the Coulomb wave functions. It is worth while to notice that, if $`\mathrm{\Lambda }_{\mathrm{QCD}}mv^2`$, the scale $`mv^2`$ can be integrated out perturbatively from the above formula and the non-perturbative contributions reduce to the Voloshin–Leutwyler formula (1), as can easily be seen by recognizing that:
$$m\frac{ϵ_nn^6}{(mC_F\alpha _\mathrm{s})^4}=\frac{1}{3N_c}T_Fn,l\left|𝐫\frac{1}{E_nH_o}𝐫\right|n,l.$$
The above outline leads to the following conclusion. For quarkonium of a typical size smaller than $`1/\mathrm{\Lambda }_{\mathrm{QCD}}`$, the most relevant operator of the non-perturbative dynamics is the bilocal gluon condensate $`𝐄^a(t)\varphi (t,0)_{ab}^{\mathrm{adj}}𝐄^b(0)`$, which belongs to the class of non-local gluon condensates considered in the introduction. In the following section we will discuss our present knowledge of it.
### 2.1 The non-local condensate $`F_{\mu \nu }^a(x)\varphi (x,0)_{ab}^{\mathrm{adj}}F_{\lambda \rho }^b(0)`$
The correlator $`F_{\mu \nu }^a(x)\varphi (x,0)_{ab}^{\mathrm{adj}}F_{\lambda \rho }^b(0)`$ is perturbatively known at the next-to-leading order in $`\alpha _\mathrm{s}`$ . However, here we are interested in its non-perturbative behaviour. Different parametrizations have been proposed . Because of its Lorentz structure, the correlator is in general described by two form factors. A convenient choice of these consists in the chromoelectric and chromomagnetic correlators:
$$𝐄^a(x)\varphi (x,0)_{ab}^{\mathrm{adj}}𝐄^b(0),𝐁^a(x)\varphi (x,0)_{ab}^{\mathrm{adj}}𝐁^b(0).$$
The strength of the correlators is of the order of the gluon condensate. In the long range ($`x^2\mathrm{}`$) they fall off exponentially (in the Euclidean space) with some typical correlation lengths. In the following we will concentrate on these correlation lengths.
The lattice calculation , using cooling techniques, obtains the same correlation length ($`T_g`$) for both form factors, and this is
$`T_g=0.34\pm 0.02\pm 0.03\text{ fm (4 flavours, }am=0.01\text{)},`$ (5)
$`T_g=0.22\pm 0.01\pm 0.02\text{ fm (quenched)}.`$ (6)
The less accurate, but traditional (quenched) lattice calculation done in obtains two different correlation lengths for the chromoelectric ($`T_g^E`$) and the chromomagnetic correlators ($`T_g^B`$):
$`T_g^ET_g^B0.1\text{}\mathrm{\hspace{0.17em}0.2}\text{ fm (quenched)}.`$ (7)
Finally, a recent sum-rule estimation obtains $`T_g^E<T_g^B`$. The sum rule turns out not to be stable for the chromoelectric correlator, while for the chromomagnetic correlation length it gives
$`T_{g}^{}{}_{}{}^{B}=0.13_{0.02}^{+0.05}\text{ fm (3 flavours)},`$ (8)
$`T_{g}^{}{}_{}{}^{B}=0.11_{0.02}^{+0.04}\text{ fm (quenched)}.`$ (9)
The correlation lengths $`T_g^E`$ and $`T_g^B`$ have a precise physical interpretation. Their inverses correspond to the masses of the lowest-lying vector and pseudovector static quark–gluon hybrids, respectively. This can be explicitly seen in the short-range limit, $`𝐱0`$, where the hybrids (in this case also called gluelumps) operators can be explicitly constructed . The suitable effective field theory is pNRQCD in the static limit . Gluelump operators are of the type $`\mathrm{Tr}\{\mathrm{O}H\}`$, where $`\mathrm{O}=O^aT^a`$ corresponds to a quark–antiquark state in the adjoint representation (the octet) and $`H=H^aT^a`$ is a gluonic operator. By matching the QCD static hybrid operators into pNRQCD, we get the static energies (also called potentials) of the gluelumps. At leading order in the multipole expansion, they read
$`V_H(r)=V_o(r)+{\displaystyle \frac{1}{T_g^H}},`$ (10)
$`H^a(t)\varphi (t,0)_{ab}^{\mathrm{adj}}H^b(0)^{\mathrm{non}\mathrm{pert}.}he^{it/T_g^H}+\mathrm{}.`$
Since hybrids are classified in QCD according to the representations of $`D_{\mathrm{},h}`$, while in pNRQCD, where we have integrated out the length $`r`$, their classification is done according to the representations of $`O(3)\times C`$, the static hybrid short-range spectrum is expected to be more degenerate than the long-range one . The lattice measure of the hybrid potentials done in confirms this feature. In it has been shown that the quantum numbers attribution of pNRQCD to the short-range operators, and the expected $`O(3)\times C`$ symmetry of the effective field theory match the lattice measurements. By using only $`𝐄`$ and $`𝐁`$ fields and keeping only the lowest-dimensional representation we may identify the operator $`H`$ for the short-range hybrids called $`\mathrm{\Sigma }_g^+^{}`$ (and $`\mathrm{\Pi }_g`$) with $`𝐫𝐄`$ (and $`𝐫\times 𝐄`$) and the operator $`H`$ for the short-range hybrids called $`\mathrm{\Sigma }_u^{}`$ (and $`\mathrm{\Pi }_u`$) with $`𝐫𝐁`$ (and $`𝐫\times 𝐁`$). Hence, the corresponding static energies for small $`r`$ are
$$V_{\mathrm{\Sigma }_g^+^{},\mathrm{\Pi }_g}(r)=V_o(r)+\frac{1}{T_g^E},V_{\mathrm{\Sigma }_u^{},\mathrm{\Pi }_u}(r)=V_o(r)+\frac{1}{T_g^B}.$$
The lattice measure of shows that, in the short range, $`V_{\mathrm{\Sigma }_g^+^{},\mathrm{\Pi }_g}(r)>V_{\mathrm{\Sigma }_u^{},\mathrm{\Pi }_u}(r)`$. This supports the sum-rule prediction that the pseudovector hybrid lies lower than the vector one, i.e. $`T_g^E<T_g^B`$.
## 3 Wilson loop operators: $`\mathrm{\Lambda }_{\mathrm{QCD}}mv`$
For higher quarkonium levels, $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ is expected to be comparable with $`mv`$. We cannot match into pNRQCD perturbatively, since the scale associated to the quarkonium size $`r`$ is already non-perturbative. The relevant non-perturbative dynamics is therefore contained in more extended objects than (local or non-local) gluon condensates: Wilson loops and field insertions on these.
In particular, disregarding effects due to scales lower than $`\mathrm{\Lambda }_{\mathrm{QCD}}`$, the static potential is given by
$$V_0(r)=\underset{T\mathrm{}}{lim}\frac{i}{T}\mathrm{ln}W_{\mathrm{}},$$
(11)
where $`W_{\mathrm{}}`$ is the static Wilson loop of size $`r\times T`$ and $``$ means an average over the gauge fields. Lattice studies tell us that at distances $`r1/\mathrm{\Lambda }_{\mathrm{QCD}}`$ the potential is no longer Coulombic but rises linearly ($`V_0(r)\sigma r`$). Higher-order corrections in the $`1/m`$ expansion have been calculated over the years and are given by field strength insertions on the Wilson loop. For instance the next-to-leading potential in the $`1/m`$ expansion is
$`{\displaystyle \frac{V_1}{m}}`$ $`=`$ $`{\displaystyle \frac{1}{m}}\underset{T\mathrm{}}{lim}({\displaystyle \frac{g^2}{4T}}{\displaystyle _{T/2}^{T/2}}dt{\displaystyle _{T/2}^{T/2}}dt^{}|tt^{}|[𝐄(t)𝐄(t^{})_{\mathrm{}}`$ (13)
$`𝐄(t)_{\mathrm{}}𝐄(t^{})_{\mathrm{}}]),`$
where $`_{\mathrm{}}`$ means a normalized gauge average in the presence of the static Wilson loop.
Wilson-loop operators of the above type may be interpreted as a superposition of states, describing gluonic excitations between static sources . They have been so far evaluated only inside QCD vacuum models or by lattice simulations (for some reviews, see and also ).
### 3.1 The stochastic expansion
We can now ask if some relation can be established between the Wilson loop (and field strength insertions on it) and the non-local gluon condensates discussed above. In fact we can express the Wilson loop as a formal expansion in terms of gluonic correlation functions by means of the so-called stochastic expansion
$`\mathrm{ln}W_{\mathrm{}}`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(ig)^n}{n!}}{\displaystyle _{S(\mathrm{})}}dS_{\mu _1\nu _1}(u_1)\mathrm{}dS_{\mu _n\nu _n}(u_n)\varphi (0,u_1)`$ (14)
$`\times F^{\mu _1\nu _1}(u_1)\varphi (u_1,0)\mathrm{}\varphi (0,u_n)F^{\mu _n\nu _n}(u_n)\varphi (u_n,0)_{\mathrm{cum}},`$
where $`S(\mathrm{})`$ denotes a surface whose contour the rectangular Wilson loop. The cumulants $`_{\mathrm{cum}}`$ are defined as
$`\varphi (0,u_1)F(u_1)\varphi (u_1,0)_{\mathrm{cum}}=\varphi (0,u_1)F(u_1)\varphi (u_1,0)=0,`$
$`\varphi (0,u_1)F(u_1)\varphi (u_1,u_2)F(u_2)\varphi (u_2,0)_{\mathrm{cum}}=`$
$`\varphi (0,u_1)F(u_1)\varphi (u_1,u_2)F(u_2)\varphi (u_2,0)`$
$`\varphi (0,u_1)F(u_1)\varphi (u_1,0)\varphi (0,u_2)F(u_2)\varphi (u_2,0)\mathrm{}=`$
$`F(u_1)\varphi (u_1,u_2)^{\mathrm{adj}}F(u_2),`$
$`\mathrm{}\mathrm{}`$
It is important to realize that the expansion of Eq. (14) is substantially different with respect to the expansions in $`1/m`$ and $`r`$, which led to the construction of the low-energy effective field theories discussed above. Those expansions were justified by the dynamics of the system under study, and each term of it is, indeed, suppressed by powers of the highest dynamical scale left divided by the scale (which is larger) that has been integrated out. Instead, from a power counting point of view, each term of the expansion (14) is of the same size (for instance each term of the series is in general expected to contribute to the string tension $`\sigma `$). As soon as we truncate the series, e.g. up to the bilocal cumulant (the first non-vanishing one), we introduce an uncontrolled approximation and define a model. This model is known as the model of the stochastic vacuum . Indeed, this model turns out to be quite successful in the study of processes that involve quark systems that may be described by almost static Wilson loops (for some reviews, see ). Let us only mention that the model predicts, in agreement with the lattice data, a long-range linear static potential with slope $`\sigma T_g^{E\mathrm{\hspace{0.17em}2}}G_2`$. It would be highly desirable to have a field theoretical justification of the expansion (14); however, such a justification is missing up to now (for a recent investigation on higher cumulants, see ).
## 4 Conclusions
In the previous sections we have shown how the non-perturbative QCD vacuum enters the dynamics of heavy quarkonium. As much as the non-perturbative scale $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ is bigger than the dynamical scales of the non-relativistic system ($`m`$, $`mv`$ and $`mv^2`$), as extended the relevant non-perturbative operators are. In the situation $`\mathrm{\Lambda }_{\mathrm{QCD}}mv^2`$ these operators reduce to local condensates, i.e. some numbers. This is the most favourable situation for a theoretical investigation. The bottomonium and charmonium ground states have been investigated in this framework (for a recent review, see ). In the limiting case of $`t\overline{t}`$ threshold production the non-perturbative corrections can be neglected and the investigation is completely accessible to perturbative QCD . In the situation mv2<
ΛQCDmvmuch-less-than𝑚superscript𝑣2<
subscriptΛQCD𝑚𝑣mv^{2}{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 6.0pt\vbox{\hbox{$\sim$}}}}\ }\Lambda_{\rm QCD}\ll mv the non-perturbative physics is encoded into non-local condensates. The dominant one is the bilocal gluon condensate. As we have discussed above it essentially depends on its strength, i.e. the gluon condensate, and on some correlation lengths, which can be related to the masses of the lowest-lying (static) quark–gluon hybrid resonances. Although this situation seems quite interesting, it has been poorly investigated in heavy quarkonium phenomenology, and mainly in the framework of models . The main reason is, in our opinion, that, while the quarkonium ground state seems to be accessible by a purely perturbative treatment (plus local condensates) and the potential describing higher excited quarkonium states seems entirely dominated by non-perturbative effects, it is, a priori, not clear to which states the situation mv2<
ΛQCDmvmuch-less-than𝑚superscript𝑣2<
subscriptΛQCD𝑚𝑣mv^{2}{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 6.0pt\vbox{\hbox{$\sim$}}}}\ }\Lambda_{\rm QCD}\ll mv is applicable. We will come back to this point at the end of this section. Finally, in the situation $`\mathrm{\Lambda }_{\mathrm{QCD}}mv`$ the relevant non-perturbative objects are Wilson-loop operators. A lattice study of them and the subsequent quarkonium spectroscopy have been done in .
The outlined study is rigorous and allows a systematic disentanglement of the high from the low energy scales of the heavy quarkonium system under study. In this specific sense also perturbative and non-perturbative effects turn out to be disentangled.
We conclude by mentioning a somehow delicate point, connected with the scales of heavy quarkonium systems: the difficulty to state a priori in which kinematic situation a particular quarkonium state is. This is mainly due to the fact that the scales $`mv`$ and $`mv^2`$ are not so well defined, nor so widely separated that different kinematic situations cannot overlap in a preliminary analysis (for a related discussion on the energy scales also in quarkonium physics, see ). Therefore, there are situations in which scales can only be fixed a posteriori, i.e. by assuming a particular situation and by checking that the final result is consistent with it and with the experimental data.
Acknowledgements
N.B. would like to thank for the invitation Jan Rafelski, chairman and organizer of the Stochastic Vacuum Session and the organizers of the conference, Herbert Fried, Bernd Müller and Yves Gabellini. N.B. is grateful to the theory group of the University of Milan for supporting her participation to the conference. N.B. and A.V. gratefully acknowledge the Alexander von Humboldt Foundation.
## References
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# Flow Between Two Sites on a Percolation Cluster
## I Introduction
Percolation theory is a paradigmatic model for connectivity, originally introduced as a mathematical subject in the late 1950s . Thereafter, percolation theory has been found useful to characterize many disordered systems .
The simplest percolation model is a lattice of bonds occupied with probability $`p`$. Neighboring bonds are considered to be connected if both are occupied. A set of sites connected by bonds is called a cluster. As $`p`$ increases, new clusters are formed and previously existing clusters not only grow, but become connected as more sites are occupied in the system. At a critical value of $`p`$, $`p_c`$ (known as the percolation threshold), one spanning cluster appears and provides overall connectivity. Just at the critical point, the incipient infinite percolation cluster is an example of a random fractal that is a useful model for real disordered systems. While the actual threshold value $`p_c`$ depends on the particular lattice chosen, the behavior of the properties measured near the percolation threshold is universal, depending only on the dimensionality of the system. For example, the mass $`M`$ of a cluster diverges at the percolation threshold as a power $`M|pp_c|^\gamma `$ where $`\gamma =43/18`$ ($`d=2`$) and $`1.795\pm 0.005`$ ($`d=3`$) regardless of the type of lattice . This universality concept is extremely useful as it means that we can understand a new system knowing that its critical exponents are the same as previously-studied systems. Also, off-lattice continuum percolation appears to have the same exponent values as those for lattice percolation .
A comprehensive set of exactly- and numerically-calculated critical exponents is now available to describe many of the features of percolation , and the percolation paradigm has been applied to problems of practical interest in heterogeneous chemistry , polymer science and transport phenomena in disordered systems . In studying transport phenomena in disordered systems, one must couple additional processes to the geometrical features of the percolation representation. Typical examples include the problems of diffusion, reaction and flow through porous media, as well as the metal-to-insulator transition in polymer composites and alloys. These are systems in which the interplay between structure and phenomenology must be investigated in detail, as it might be the relevant factor determining optimum material properties.
A useful concept—first put forward by Ambegaokar, Halperin and Langer —is that electrical transport in disordered media with a broad distribution of conductance values is dominated by those regions where the conductances are larger than some critical value $`\sigma _c`$. This critical value is the largest conductance such that the set of conductances above this threshold value still preserves the global connectivity of the system. In percolation terminology, this cluster would be analogous to the conducting spanning cluster. It is then possible to reduce the general problem of transport in a highly connected disordered media with a broad distribution of conductances to a percolation problem at criticality. Once the “critical cluster” is identified within the disordered geometry, one can estimate macroscopic transport properties of the system. This approximation, commonly called “the critical path method” , has been extended by Katz and Thompson to estimate transport properties (e.g., permeability and electrical conductivity) in disordered materials. Thus, although the majority of studies of fluid flow through disordered media are for systems well above the threshold (e.g., sandstone), if the distribution of permeabilities is sufficiently broad, we can still make use of percolation concepts to model the relevant geometry of well-connected disordered media. Therefore, percolation theory is certainly an appropriate description for a large number of disordered systems (see also Ref. and references therein). This is the basis of our study and also the basis of the vast literature on fluid flow in percolation clusters.
The aim of the present paper is to discuss the potential application of percolation theory as a convenient geometrical model for understanding numerous aspects of flow through porous rocks . Special emphasis will be given to the study of oil displacement, i.e., how hydrocarbons propagate through geological formations between a pair of wells in the oil field. This work could also be applied to the breakthrough time for contamination of a water supply, or the time for released radioactive material to get from a leaking nuclear repository into the biosphere.
Oil fields are extremely complex, containing geological heterogeneities on a wide range of length scales from centimeters to kilometers. These heterogeneities, caused by the sedimentary processes that deposited the rocks and the subsequent actions on the rock, such as fracturing by tectonic forces and mineral deposition from aquifer flow, have a significant impact on hydrocarbon recovery.
The most common method of oil recovery is by displacement. Either water or a miscible gas (carbon dioxide or methane) is injected in a well (or wells) to displace the oil to other wells. Ultimately the injected fluid will break through into a production well where it must be separated from the oil, which is a very costly process. Once the injected fluid has broken through, the rate of oil production declines as more injected fluid is produced. For economic purposes it is important to know when the injected fluid will break through and what the rate of decline of oil production will be so that the economic limit of production can be determined.
Because the sedimentary process that produces the porous rock through which the fluid flows is very chaotic, the rock is highly heterogeneous. However, in many cases the rock can be separated into two types—high permeability (“good”) and low or zero permeability (“bad”)—and for all practical purposes we can assume the flow takes place only in the good rock. The spatial distribution of the rock types is often close to random, in which case the classical percolation problem is a good approximation. The place of the occupancy probability $`p`$ is taken by the volume fraction of the good rock, called the net-to-gross ratio in the petroleum literature.
We have very little direct knowledge about the spatial distribution of rock properties in a reservoir. Direct measurements are limited to samples that represent a fraction of $`10^{13}`$ of the total reservoir volume. These samples are taken from the well locations. Elsewhere the properties have to be inferred from knowledge of the general geological environment and by analogy with other reservoirs or surface outcrops. Hence, there is a great deal of uncertainty in our prediction of the spatial distribution of these rock properties. This leads to an uncertainty in our ability to predict the flow performance, principally the breakthrough time and the production decline rate. We need to estimate the uncertainty accurately so that economic risk evaluations can be made.
The conventional approach to this problem is to build a detailed (numerical) model of the rock properties. These models will honor the one and two point statistics observed from the wells and analogue outcrops. The models must also agree with the observed data values at wells. These models are statistical in nature and conventionally one samples realizations from the models and performs numerical flow calculations on the realizations to give a Monte Carlo prediction of breakthrough and production decline. Unfortunately this process is so time consuming as to be impractical in many cases. Typically the flow simulation can take several hours on reasonable workstations. When hundreds of realizations are sampled to get good statistics, the total computing time becomes unwieldy. Thus there is a strong need to make this more efficient so as to come up with very quick, but accurate, predictions of recovery and the uncertainty due to the lack of knowledge of the underlying rock properties.
The purpose of this work is to accomplish the goals described above using methods derived from percolation theory. It is based on two key assumptions. One is that for many cases the permeability disorder can be approximated by either permeable or impermeable rock. For example the reservoir may have been deposited by meandering river belts in which case the good sand occurs as “packages” in an insulating background. The other assumption is that the flow paths are strongly controlled by the permeability disorder and not strongly modified by the flow dynamics themselves. Again there are many cases when this is reasonable, in particular if the viscosity ratio between the injected and displaced fluids is not too large or when the system is highly disordered.
Under these assumptions we can consider the underlying heterogeneity to be that of a percolating system (not necessarily at threshold). This has previously been done to study the fraction of sand connected to well pairs . We then look at the dynamic displacement on this cluster where the flow is controlled by Darcy’s law (analogous to Ohm’s law in electrical current flow). We assume that the injected fluid can be treated as a passive tracer (i.e., one that is not absorbed by the rock) that is convected along these flow paths. Then the breakthrough time is the same as the first passage time and is strongly controlled by the shortest path length between the wells; the post breakthrough production decline is controlled by the longer paths. Figure 1a illustrates a typical percolation backbone and the shortest path between two points. Note that the lines here represent not microscopic pores but rather sand bodies whose size are of the order of tens of meters.
The rest of the paper is organized as follows. In Sec. II, because of its importance to flow in the cluster, we study the distribution of shortest paths between two sites in a percolation cluster. In Sec. III we study the distribution functions for the dynamic quantities: minimal traveling time, and length of the path corresponding to minimal traveling time. Finally, in Sec. IV, we draw some conclusions and discuss possible future work.
## II Shortest path
This section deals with the distribution of the shortest path between two sites on a percolation cluster. Because of the qualitative resemblance between the shortest path and the minimal traveling time of a tracer particle, the first step in understanding fluid transport between two sites in a percolation system is to characterize the geometrical properties of the shortest connecting path. For example, if we assume that the traveling time along a path is proportional to the path length (i.e., all velocities are equal), then we can obtain a rough estimate for the traveling time from purely geometrical arguments.
### A Basic distribution functions
The shortest path or chemical distance, $`\mathrm{}`$, between two sites on a percolation cluster is defined as the shortest path connecting the two sites (Fig. 1) . The typical value $`\mathrm{}^{}`$ of the shortest path between two sites on a cluster scales with the geometrical distance, $`r`$, between these points as
$$\mathrm{}^{}r^{d_{\text{min}}},$$
(1)
where
$$d_{\text{min}}=\{\begin{array}{cc}1.13\pm 0.02\hfill & [d=2]\hfill \\ 1.374\pm 0.02\hfill & [d=3]\hfill \end{array}$$
(2)
is the fractal dimension of the shortest path .
Consider a hypercubic lattice of $`L^d`$ sites. All information about the distribution of shortest paths is contained in the joint probability density function $`P(r,\mathrm{})`$, i.e., the probability that two sites on the same spanning cluster are separated by geometrical distance $`r`$ and chemical path $`\mathrm{}`$. We sum over all chemical paths $`\mathrm{}`$to calculate the probability distribution that the Euclidean distance between two sites is $`r`$,
$$P(r)P(r,\mathrm{})𝑑\mathrm{}.$$
(3)
Similarly, we obtain the probability distribution that two sites are separated by the chemical distance $`\mathrm{}`$by summing over all possible geometrical distances,
$$P(\mathrm{})P(r,\mathrm{})𝑑r.$$
(4)
Given that the shortest distance between these sites is $`\mathrm{}`$, the conditional probability that the geometrical distance between two sites is $`r`$ is
$$P(r|\mathrm{})=\frac{P(r,\mathrm{})}{P(\mathrm{})}.$$
(5)
For isotropic media this function has been studied extensively and $`P(r|\mathrm{})`$ is of the form
$$P(r|\mathrm{})=A_{\mathrm{}}\left(\frac{r}{\mathrm{}^{\stackrel{~}{\nu }}}\right)^{g_r}\mathrm{exp}\left(a\left(\frac{r}{\mathrm{}^{\stackrel{~}{\nu }}}\right)^{\stackrel{~}{\delta }}\right),$$
(6)
where
$$\stackrel{~}{\delta }=\frac{1}{1\stackrel{~}{\nu }}=\frac{d_{\text{min}}}{d_{\text{min}}1},$$
(7)
and
$$\stackrel{~}{\nu }1/d_{\text{min}}.$$
(8)
For $`d=2`$, Ziff recently argued that
$$g_r1=25/24[d=2].$$
(9)
Our simulations confirm the analytical form of the $`P(r|\mathrm{})`$ as well as these values of $`\stackrel{~}{\nu }`$ and $`g_r`$ (see Fig. 2a).
The function of interest to us is the conditional probability for two sites to be separated by the shortest path $`\mathrm{}`$, given that the geometrical distance between these sites is $`r`$
$$P(\mathrm{}|r)=\frac{P(r,\mathrm{})}{P(r)}.$$
(10)
¿From (10) and (5), we see that $`P(r|\mathrm{})`$ and $`P(\mathrm{}|r)`$ are related
$$P(\mathrm{}|r)=P(r|\mathrm{})\frac{P(\mathrm{})}{P(r)}.$$
(11)
At the percolation threshold, it has been shown that, in analogy with (6),
$$P(\mathrm{}|r)\frac{1}{r^{d_{\text{min}}}}\left(\frac{\mathrm{}}{r^{d_{\text{min}}}}\right)^g_{\mathrm{}}\mathrm{exp}\left(a\left(\frac{\mathrm{}}{r^{d_{\text{min}}}}\right)^\varphi _{\mathrm{}}\right),$$
(12)
where
$$g_{\mathrm{}}1=\frac{(g_r1)+(2d_f)}{d_{\text{min}}},$$
(13)
$$\varphi _{\mathrm{}}=\stackrel{~}{\delta }\stackrel{~}{\nu }=\stackrel{~}{\nu }/(1\stackrel{~}{\nu })=\frac{1}{d_{\text{min}}1},$$
(14)
and
$$d_f=\{\begin{array}{cc}91/48\hfill & [d=2]\hfill \\ 2.524\pm 0.008\hfill & [d=3]\hfill \end{array}$$
(15)
is the fractal dimension of the incipient infinite cluster . Substituting (9) into (13), we find for $`d=2`$
$$g_{\mathrm{}}=2.01\pm 0.02[d=2].$$
(16)
The probability distribution of more practical interest is $`P^{}(\mathrm{}|r)`$, defined in the same way as $`P(\mathrm{}|r)`$ but for any two randomly-chosen points separated by geometrical distance $`r`$ and on the same cluster, but not necessarily on the incipient infinite cluster . As shown in Appendix A, $`P^{}(\mathrm{}|r)`$ has the same scaling form as in Eq. (12), but with $`g_{\mathrm{}}`$ replaced by
$$g_{\mathrm{}}^{}=g_{\mathrm{}}+\frac{dd_f}{d_{\text{min}}}.$$
(17)
Figure 2b illustrates the difference between $`g_{\mathrm{}}`$ and $`g_{\mathrm{}}^{}`$.
### B Distribution of shortest path
The complete scaling form of $`P^{}(\mathrm{}|r)`$, which accounts also for finite size effects and off-critical behavior, has been studied for $`d=2`$ and reported in . Specifically, the following Ansatz has been proposed
$`P^{}(\mathrm{}|r)`$ $``$ $`{\displaystyle \frac{1}{r^{d_{\text{min}}}}}\left({\displaystyle \frac{\mathrm{}}{r^{d_{\text{min}}}}}\right)^g_{\mathrm{}}^{}f_1\left({\displaystyle \frac{\mathrm{}}{r^{d_{\text{min}}}}}\right)`$ (18)
$`f_2\left({\displaystyle \frac{\mathrm{}}{L^{d_{\text{min}}}}}\right)f_3\left({\displaystyle \frac{\mathrm{}}{\xi ^{d_{\text{min}}}}}\right),`$ (19)
where $`\xi |pp_c|^\nu `$ is the pair connectedness length, and the scaling functions have the form
$$f_1(x)\mathrm{exp}(ax^\varphi ),$$
(20)
$$f_2(x)\mathrm{exp}(bx^\psi ),$$
(21)
and
$$f_3(x)\mathrm{exp}(cx).$$
(22)
The function $`f_1`$ accounts for the lower cut-off due to the constraint $`\mathrm{}>r`$, while $`f_2`$ and $`f_3`$ account for the upper cut-offs due to the finite size effect and due to the finite correlation length respectively. Either $`f_2`$ or $`f_3`$ becomes irrelevant, depending on the magnitudes of $`L`$ and $`\xi `$: for $`L<\xi `$, $`f_2`$ dominates the upper cut-off, otherwise $`f_3`$ dominates. We assume the independence of the finite size effect and the effect of the concentration of the occupied sites, so that Eq. (19) can be represented as a product of the terms which are responsible for the finite size effect ($`f_2`$) and the effect of the concentration ($`f_3`$). Simulations for $`d=2`$ have been used to test this assumption.
### C Shortest path in three dimensions
#### 1 Behavior at criticality
Here we extend the study of $`P^{}(\mathrm{}|r)`$ to $`d=3`$. We numerically test the scaling conjecture (19) exactly at the percolation threshold $`p=p_c`$—in which case $`\xi =\mathrm{}`$ so $`f_3=f(0)=\text{const}`$. We build clusters using the Leath algorithm . Since the Leath algorithm corresponds to the process of selecting a random point on the lattice, the probability $`P^{}(\mathrm{}|r)`$ is equal to the probability that a pair of randomly selected points has chemical distance $`\mathrm{}`$ and geometrical distance $`r`$, given that they belong to the same cluster, a cluster that is not necessarily the infinite cluster. Hence Eq. (19) reduces to
$`P^{}(\mathrm{}|r)`$ $``$ $`{\displaystyle \frac{1}{r^{d_{\text{min}}}}}\left({\displaystyle \frac{\mathrm{}}{r^{d_{\text{min}}}}}\right)^g_{\mathrm{}}^{}`$ (23)
$`f_1\left({\displaystyle \frac{\mathrm{}}{r^{d_{\text{min}}}}}\right)f_2\left({\displaystyle \frac{\mathrm{}}{L^{d_{\text{min}}}}}\right)[p=p_c].`$ (24)
Figure 3a shows that, in the range $`r^{d_{\text{min}}}<\mathrm{}<L^{d_{\text{min}}}`$, $`P^{}(\mathrm{}|r)`$ has power-law behavior with slope
$$g_{\mathrm{}}^{}=2.3\pm 0.1,[d=3]$$
(25)
and rapidly vanishes for $`\mathrm{}<r^{d_{\text{min}}}`$ and for $`\mathrm{}>L^{d_{\text{min}}}`$. To determine the functions $`f_1`$ and $`f_2`$, we compute the rescaled probability distribution
$$\mathrm{\Phi }\left(\frac{\mathrm{}}{r^{d_{\text{min}}}}\right)P^{}(\mathrm{}|r)\mathrm{}^g_{\mathrm{}}^{}r^{d_{\text{min}}(g_{\mathrm{}}1)},$$
(26)
and plot it against scaling variable $`x\mathrm{}/r^{d_{\text{min}}}`$ (see Fig. 3b) using the value $`d_{\text{min}}=1.374`$. According to Eq. (24)
$$\mathrm{\Phi }(x)=Af_1(x)f_2\left[x\left(\frac{r}{L}\right)^{d_{\text{min}}}\right].$$
(27)
Therefore, $`\mathrm{\Phi }(x)`$ should depend only on $`x`$ and the ratio $`r/L`$. Indeed, Fig. 3b shows excellent data collapse for $`L/r=8`$, with sharp cutoffs governed for $`x<1`$ by $`f_1(x)`$ and for $`x>(L/r)^{d_{\text{min}}}`$ by $`f_2[x(r/L)^{d_{\text{min}}}]`$.
In order to test the assumption that the functions $`f_1`$ and $`f_2`$ are stretched exponentials with exponents $`\varphi _{\mathrm{}}`$ and $`\psi _{\mathrm{}}`$, we plot
$$\mathrm{\Pi }(x)\mathrm{log}_{10}[A/\mathrm{\Phi }(x)]$$
(28)
versus $`x`$ in double logarithmic scale for various values of normalization constant $`A`$ (see Fig. 3c). If the stretched exponential conjecture is correct, $`\mathrm{\Pi }(x)`$ should have two straight line asymptotes for $`\mathrm{log}x+\mathrm{}`$ with the slope $`\psi _{\mathrm{}}`$ and for $`\mathrm{log}x\mathrm{}`$ with the slope $`\varphi _{\mathrm{}}`$. We find that the slopes $`\varphi _{\mathrm{}}`$ and $`\psi _{t_{\text{min}}}`$ of the straight line fits depend weakly on the value of $`A`$. Using $`A=0.08`$, we obtain the longest regimes of straight line behavior. For this value of $`A`$, we find $`\varphi _{\mathrm{}}2.1`$ and $`\psi _{\mathrm{}}2.5`$. Equation (14) yields a predicted value of $`\varphi _{\mathrm{}}=2.67`$ in good agreement with our simulation result.
#### 2 Off-critical behavior
For $`pp_c`$, we identify three regimes determined by the value of the connectedness length, $`\xi `$, in relation to the values of $`r`$ and $`L`$:
* $`\xi >L>r`$. In this regime, the fact that $`pp_c`$ cannot be detected because the connectedness length is larger than the other relevant variables.
* $`L>\xi >r`$. In this case, the upper cutoff of the distribution Eq. (19) is governed by $`f_3`$ and the functional form of the rescaled probability $`\mathrm{\Phi }`$ is given by
$$\mathrm{\Phi }(\mathrm{}/r^{d_{\text{min}}})f_1\left(\frac{\mathrm{}}{r^{d_{\text{min}}}}\right)f_3\left(\frac{\mathrm{}}{\xi ^{d_{\text{min}}}}\right).$$
(29)
For large $`\mathrm{}`$, we suggest an exponential decay of $`\mathrm{\Phi }`$
$$\mathrm{\Phi }(\mathrm{}/r^{d_{\text{min}}})\mathrm{exp}\left(c\frac{\mathrm{}}{\xi ^{d_{\text{min}}}}\right).$$
(30)
Indeed, for $`p<p_c`$, semi-logarithmic plots of $`\mathrm{log}\mathrm{\Phi }(\mathrm{}/r^{d_{\text{min}}})`$ versus $`\mathrm{}`$ shown in Fig. 4a can be approximated by straight lines with slopes which approach zero as $`pp_c`$. According to Eq. (30), these slopes $`k(p)`$ should be proportional to $`\xi ^{d_{\text{min}}}|pp_c|^{d_{\text{min}}\nu }|pp_c|^{1.19}`$. Figure 4b shows a double logarithmic plot of $`|k(p)|`$ versus $`|pp_c|`$ for $`p<p_c`$. This curve can be well approximated by a straight line with slope 1.22 in good agreement with the scaling conjecture (24). For $`p>p_c`$ a similar analysis should hold. However, limitations on the size of the system we can simulate make the analysis problematic. Figure 4c shows $`P^{}(\mathrm{})`$ for various values of $`p>p_c`$. Note that it is only for values of $`pp_c+0.03`$ that the distributions “cut-off” at smaller $`\mathrm{}`$ than the distribution for $`p=p_c`$. Thus it is only for values of $`pp_c0.03`$ that the large $`\mathrm{}`$ behavior of Eq. (19) is determined by the fact that the system is not at criticality (i.e., by $`f_3`$) as opposed to being determined by the finite size of the system (i.e., by $`f_2`$). Below $`p=p_c+0.03`$, $`\xi `$ is still greater than $`L`$. On the other hand, if $`p`$ is not close to $`p_c`$, the scaling form is not expected to hold. Thus, the results are inconclusive based on the sizes of the systems we can generate—we cannot determine the parameters that govern the large $`\mathrm{}`$ behavior of Eq. (19) above $`p_c`$.
* $`L>r>\xi `$. When the connectedness length $`\xi `$ is smaller than the distance $`r`$ between the wells, the system can be considered homogeneous . This can be seen in Fig. 5a in which we plot $`P(\mathrm{}|r)`$ for various values of $`r`$ at $`p=0.7`$ for 2d site percolation ($`p_c=0.593`$). As $`r`$ increases from below to above the connectedness length, the form of the distribution changes from the power law distribution of Eq. (19) to a Gaussian distribution with a pronounced peak, a characteristic of homogeneous systems. Furthermore, as shown in Fig. 5b, the fractal dimension of the shortest length crosses over from $`d_{\text{min}}=1.13`$ to $`d_{\text{min}}=1.0`$, characteristic of a homogeneous system . The convergence to a Gaussian can be expected due to the following considerations. The minimal path connecting the wells separated by distance $`r`$ passes through $`r/\xi `$ independent blobs. For each of these blobs, the probability distribution for the shortest path across the blob, $`\mathrm{}_b`$, is still given by Eq. (19), but with $`r`$ and $`L`$ replaced by $`\xi `$ and $`\mathrm{}`$ replaced by $`\mathrm{}_b`$. This distribution is characterized by $`\mathrm{}_b\xi ^{d_{\text{min}}}`$ and variance $`\sigma _b^2\mathrm{}_b^2\mathrm{}_b^2\xi ^{2d_{\text{min}}}`$. The total minimal path is the sum of $`n=r/\xi `$ independent variables $`\mathrm{}_b`$, hence it converges to a Gaussian with
$$\mathrm{}r\xi ^{d_{\text{min}}1}\text{and}\sigma ^2r\xi ^{2d_{\text{min}}1}.$$
(31)
Thus the slope of the graph, $`k(p)`$, of $`\mathrm{}`$ vs $`r`$ in Fig. 5c should decay as
$$k(p)|pp_c|^{\nu (d_{\text{min}}1)}=|pp_c|^{0.17}$$
(32)
and the slope of $`\sigma ^2`$ versus $`r`$ should decay as
$$|pp_c|^{\nu (2d_{\text{min}}1)}=|pp_c|^{1.7}.$$
(33)
Indeed, (see Fig 5d) we see that the slope of $`\sigma ^2`$ versus $`r`$ decays with $`p`$ more strongly than that of $`\mathrm{}`$ versus $`r`$. The numerical values of slopes from Figs. 5c and 5d are in good agreement with the theoretical predictions Eqs. (32,33). For $`d=3`$ we expect similar behavior.
### D Rectangular boundary conditions
Since realistic oil fields do not have square boundaries, it is reasonable to ask what is the effect of rectangular boundary conditions. Here we study the distributions of the shortest length and minimal time on a two-dimensional lattice of size $`L_x\times L_y`$, $`L_xL_y`$. We position the wells at points $`A`$ and $`B`$ separated by distance $`r=16`$ along $`x`$-axis. We study two cases (Fig. 6):
* $`L_x`$ is fixed ($`L_x=32`$), and we vary $`L_y`$ ($`L_y=64`$, 128, 256, 512, 1024);
* $`L_y`$ is fixed ($`L_y=32`$), and we vary $`L_x`$ ($`L_x=64`$, 128, 256, 512, 1024).
We find that (i) the shortest length and minimal time distributions are identical in all of the above cases and (ii) the scaling form for these anisotropic cases is the same as the isotropic case with $`L`$ replaced by the minimum of $`L_x`$ and $`L_y`$, i. e. the scaling form of the distribution Eq. (19) remains unchanged with the exception that $`L`$ is determined by
$$L=\mathrm{min}(L_x,L_y).$$
(34)
Equation (34) can be a result of competing exponentials $`f_{2x}`$ and $`f_{2y}`$. Since both $`f_{2x}`$ and $`f_{2y}`$ are rapidly decaying functions with $`L_x`$ and $`L_y`$ \[see Eq. (21)\], the finite-size cutoff of $`P(\mathrm{}|r)`$ is determined by the smaller of $`L_x`$ and $`L_y`$. The fact that the results are independent of the axis along which the wells are aligned can be understood by realizing that the probability that “oblong” clusters connect two points separated by a distance longer than the minimum of $`L_x`$ and $`L_y`$ is low. This finding has the implication that in anisotropic fields, a number of well pairs would be needed to optimize the recovery of oil in the field.
## III Minimal Traveling Time and Fastest Path
We turn next to dynamics, the study of flow on percolation clusters, which has close ties to such applications as hydrocarbon recovery and ground-water pollution . In this section, we discuss the properties of the flow on $`d=2`$ and $`d=3`$ bond percolation clusters. Specifically, we investigate the scaling properties of the distributions of minimal traveling time and the length of the path corresponding to the minimal traveling time (fastest path) of the tracer particles. Some of the results in $`d=2`$ were reported previously . Here we extend the work to $`d=3`$, and study the effects of finite system size and off-criticality for $`d=2`$ and $`d=3`$.
### A The model
We study incompressible flow between two sites $`A`$ and $`B`$ separated by Euclidean distance $`r`$ (see Fig. 1). To model the flow front, we use passive tracers—particles not absorbed by the surroundings, that move only by convection, ignoring molecular diffusion which is slow on the time scales of interest. The convection is governed by the flow field due to the pressure difference between sites connected by the bonds. We simulate the flow of a tracer particle starting from the injection point $`A`$ traveling through the medium along a path connected to the recovery point $`B`$. The dynamics of flow at a macroscopic level on the percolation cluster is determined by the local flow (local currents) on the individual bonds in the backbone of the cluster. The velocity of a tracer at each bond is determined by the pressure difference across that bond (Darcy’s law ):
$$v_{ij}=T(P_jP_i),$$
(35)
where $`P_i`$ and $`P_j`$ are the values of pressure at sites $`i`$ and $`j`$. The coefficient $`T`$, which is a function of permeability $`k`$, viscosity $`\eta `$, and the length of a bond $`L_b`$ ($`T=k/(\eta L_b)`$), is set to 1. We normalize the velocities assuming the total flow $`J`$ between $`A`$ and $`B`$ is fixed, independent of the distance between $`A`$ and $`B`$, and the realization of the porous media. This resembles more closely oil recovery processes where constant flow, as opposed to constant pressure, is maintained.
We obtain the pressure difference across each bond by solving Kirchhoff’s law
$$\underset{j}{}v_{ij}=0,$$
(36)
for each node $`i`$ in the cluster where the summation is over all bonds connected to that node. Fig. 1b shows the results of solving these equations for the cluster discussed in Section I (Fig. 1b). Magnitudes of currents on cluster backbone are depicted in gray scale with the lightest areas corresponding to the smallest currents and the darkest to the largest currents.
We define the traveling time, $`\stackrel{~}{t}`$, of a path $`𝒞`$ as the sum of the tracer’s traveling times $`t_{ij}`$ at each bond $`(ij)`$ joining sites $`i`$ and $`j`$ which are on the path,
$$\stackrel{~}{t}=\underset{(ij)𝒞}{}t_{ij}.$$
(37)
The traveling length, $`\stackrel{~}{\mathrm{}}`$, in turn, is the number of bonds present in path $`𝒞`$. Among the ensemble of all paths $`\{𝒞\}`$, we select the path $`𝒞^{}`$ that has the minimal traveling time, $`t_{\text{min}}`$,
$$t_{\text{min}}(𝒞^{})=\underset{\{𝐂\}}{\mathrm{min}}\stackrel{~}{t}(𝒞)$$
(38)
and we define the length of the fastest path, $`\mathrm{}_{\text{min}}`$, corresponding to the minimal traveling time, as the number of bonds present in path $`𝒞^{}`$. The first quantity $`t_{\text{min}}`$ is the breakthrough time of the gas/liquid that displaces the oil during recovery and has fundamental importance to the oil industry. The quantity $`\stackrel{~}{t}`$ determines post-breakthrough behavior. We also define the exponents $`d_x`$, where $`x`$ denotes $`\mathrm{}_{\text{min}}`$, $`t_{\text{min}}`$, $`\stackrel{~}{\mathrm{}}`$ or $`\stackrel{~}{t}`$ by
$$x^{}r^{d_x}$$
(39)
and where $`x^{}`$ is the characteristic (most probable) length or time of the corresponding distribution.
Using a “burning” algorithm , we then find the minimal time and the fastest path for the particle to travel between points $`A`$ and $`B`$. Figure 1c shows the propagation of the tracer particles through the same backbone shown in Fig. 1a and 1b. At $`t=t_{\text{min}}`$, the tracer particles spread over $`t=t_{\text{min}}J`$ bonds, which constitute a subset of the backbone with fractal dimension $`d_{\text{tm}}`$, which is larger than the fractal dimension of the minimal path but smaller than the fractal dimension of the entire backbone $`d_\text{B}`$. Hence
$$d_{\text{min}}<d_{\text{tm}}<d_\text{B}.$$
(40)
### B Minimal Traveling Time
We first study the minimal traveling time for $`d=2`$. In Fig. 7, a scatter plot of the minimal traveling time versus shortest path, we see that the minimal times are strongly correlated with the shortest paths in the realizations simulated, $`t_{\text{min}}\mathrm{}^z`$, where $`z1.17`$. Since $`\mathrm{}`$ scales as $`r^{d_{\text{min}}}`$ we propose that $`t_{\text{min}}`$ scales as $`r^{d_{\text{tm}}}`$ with $`d_{\text{tm}}=zd_{\text{min}}=1.33`$. This suggests that the same scaling form which applies to the distribution of shortest paths can also be applied to the distribution of minimal times, but with different exponents and amplitudes. Thus, we expect the Ansatz similar to Eq. (19)
$`P^{}(t_{\text{min}}|r)`$ $``$ $`{\displaystyle \frac{1}{r^{d_{\text{tm}}}}}\left({\displaystyle \frac{t_{\text{min}}}{r^{d_{\text{tm}}}}}\right)^{g_{\text{tm}}^{}}f_1\left({\displaystyle \frac{t_{\text{min}}}{r^{d_{\text{tm}}}}}\right)`$ (41)
$`f_2\left({\displaystyle \frac{t_{\text{min}}}{L^{d_{\text{tm}}}}}\right)f_3\left({\displaystyle \frac{t_{\text{min}}}{\xi ^{d_{\text{tm}}}}}\right),`$ (42)
where the scaling functions are $`f_1(x)=\mathrm{exp}(a_{\text{tm}}x^{\varphi _{\text{tm}}})`$, $`f_2(x)=\mathrm{exp}(b_{\text{tm}}x^{\psi _{\text{tm}}})`$ and $`f_3(x)=\mathrm{exp}(c_{\text{tm}}x^{\pi _{\text{tm}}})`$. Here $`\xi `$ is a characteristic length of the pair connectedness function and has a power-law dependence on the occupancy probability $`p`$ as
$$\xi |pp_c|^\nu .$$
(43)
The first function $`f_1`$ accounts for the lower cut-off due to the constraint $`\mathrm{}>r`$, while $`f_2`$ and $`f_3`$ account for the upper cut-offs due to the finite size effect and due to the finite connectedness length, respectively. Either $`f_2`$ and $`f_3`$ becomes irrelevant, depending on which of the two values $`L`$ or $`\xi `$ is greater. For $`L<\xi `$, $`f_2`$ dominates the upper cut-off, otherwise $`f_3`$ dominates. We have assumed independence of the finite size effect and off-criticality effect, so that Eq. (42) can be represented as a product of the terms which are responsible for the finite size effect ($`f_2`$) and the effect of the concentration ($`f_3`$).
We sample over $`10^6`$ different realizations with the two sites $`A`$ and $`B`$ fixed. For each realization, we calculate exactly the minimal traveling time and the path which corresponds to the minimal traveling time to obtain $`P(t_{\text{min}})`$ and $`P(\mathrm{}_{\text{min}})`$.
#### 1 Behavior at criticality
We first test numerically the scaling conjecture Eq. (42) at the percolation threshold $`p=p_c`$. In this case, $`\xi =\mathrm{}`$ and $`f_3`$ is a constant. Hence Eq. (42) reduces to
$`P^{}(t_{\text{min}}|r)`$ $``$ $`{\displaystyle \frac{1}{r^{d_{\text{tm}}}}}\left({\displaystyle \frac{t_{\text{min}}}{r^{d_{\text{tm}}}}}\right)^{g_{\text{tm}}^{}}`$ (44)
$`f_1\left({\displaystyle \frac{t_{\text{min}}}{r^{d_{\text{tm}}}}}\right)f_2\left({\displaystyle \frac{t_{\text{min}}}{L^{d_{\text{tm}}}}}\right)(p=p_c).`$ (45)
Figure 8a shows that $`P^{}(t_{\text{min}}|r)`$ has a power-law regime with slope
$$g_{\text{tm}}^{}=2.0\pm 0.1.$$
(46)
To determine the functions $`f_1`$ and $`f_2`$, we compute the rescaled probability distribution
$$\mathrm{\Phi }\left(\frac{t_{\text{min}}}{r^{d_{\text{tm}}}}\right)P^{}(t_{\text{min}}|r)(t_{\text{min}})^{g_{\text{tm}}^{}}r^{d_{\text{tm}}(g_{\text{tm}}1)},$$
(47)
and plot it against scaling variable $`x=t_{\text{min}}/r^{d_{\text{tm}}}`$ (see Fig. 8b). According to Eq. (45)
$$\mathrm{\Phi }(x)=Af_1(x)f_2\left[x\left(\frac{r}{L}\right)^{d_{\text{tm}}}\right].$$
(48)
Therefore, $`\mathrm{\Phi }(x)`$ should depend only on $`x`$ and the ratio $`r/L`$. Unlike the fractal dimension of the shortest path, $`d_{\text{min}}`$, there have been no calculations of the fractal dimension of the minimal traveling time, $`d_{\text{tm}}`$. We estimate $`d_{\text{tm}}`$ by finding the value which yields the best data collapse for Eq. (48). For $`d_{\text{tm}}=1.33`$, Fig. 8b shows excellent data collapse with sharp cutoffs governed for small $`x<1`$ by $`f_1(x)`$ and for large $`x>(L/r)^{d_{\text{tm}}}`$ by $`f_2[x(r/L)^{d_{\text{tm}}}]`$.
In order to test the assumption that the functions $`f_1`$ and $`f_2`$ are stretched exponentials with exponents $`\varphi _{\text{tm}}`$ and $`\psi _{\text{tm}}`$, we make a log-log plot of $`\mathrm{\Pi }(x)\mathrm{log}_{10}[A/\mathrm{\Phi }(x)]`$ versus $`x`$ for various values of the normalization constant $`A`$ (See Fig. 8c). If the stretched exponential conjecture is correct, $`\mathrm{\Pi }(x)`$ should have two straight line asymptotes for $`\mathrm{log}x+\mathrm{}`$ with the slope $`\psi _{\text{tm}}`$ and for $`\mathrm{log}x\mathrm{}`$ with the slope $`\varphi _{\text{tm}}`$. The slopes $`\varphi _{\text{tm}}`$ and $`\psi _{\text{tm}}`$ of the straight line fits depend weakly on the value of $`A`$. Using $`A=0.14`$, we obtain the longest regimes of straight line behavior. For this $`A`$ we obtained $`\varphi _{\text{tm}}3.0`$ and $`\psi _{\text{tm}}3.0`$. With the same assumptions used to derive Eq. (14), we can derive a similar expression for $`\varphi _{\text{tm}}`$
$$\varphi _{\text{tm}}=\frac{1}{d_{\text{tm}}1},$$
(49)
which yields a predicted value of $`\varphi _{\text{tm}}`$ of $`3.0`$ in good agreement with our simulation result.
#### 2 Off-Critical Behavior
Finally, in order to test the dependence of $`P^{}(t_{\text{min}}|r)`$ on $`p`$ we obtain data for large system size $`L`$ ($`L=1000`$) and for several values of $`pp_c`$. As we do for the shortest length, we analyze the behavior of $`t_{\text{min}}`$ in three regimes determined by the relation of the value of the connectedness length, $`\xi `$, to the values of $`r`$ and $`L`$.
* $`\xi >L>r`$. In this regime, the fact that $`pp_c`$ cannot be detected because the connectedness length is larger than the other relevant variables.
* $`L>\xi >r`$. In this case, the upper cutoff of the distribution Eq. (42) is governed by $`f_3`$ and the functional form of the rescaled probability $`\mathrm{\Phi }`$ is given by
$$\mathrm{\Phi }(t_{\text{min}}/r^{d_{\text{tm}}})f_1\left(\frac{t_{\text{min}}}{r^{d_{\text{tm}}}}\right)f_3\left(\frac{t_{\text{min}}}{\xi ^{d_{\text{tm}}}}\right).$$
(50)
For large $`t_{\text{min}}`$, we suggest an exponential decay of $`\mathrm{\Phi }`$
$$\mathrm{\Phi }(t_{\text{min}}/r^{d_{\text{tm}}})\mathrm{exp}\left(c\frac{t_{\text{min}}}{\xi ^{d_{\text{tm}}}}\right).$$
(51)
Semi-logarithmic plots of $`\mathrm{log}\mathrm{\Phi }(t_{\text{min}}/r^{d_{\text{tm}}})`$ versus $`t_{\text{min}}`$ for $`p>p_c`$ and $`p<p_c`$ shown in Fig. 9a and 9b, respectively, can be approximated by straight lines with slopes which approach zero as $`pp_c`$. According to Eq. (51), this slope $`k(p)`$ should be proportional to
$$\xi ^{d_{\text{tm}}}=|pp_c|^{d_{\text{tm}}\nu }|pp_c|^{1.77}.$$
(52)
Figure 9c shows double logarithmic plots of $`|k(p)|`$ versus $`|pp_c|`$ for $`p<p_c`$ and $`p>p_c`$, which can be well approximated by straight lines with slopes 1.81 and 1.77, respectively, in good agreement with the scaling conjecture, Eq. (52). As was the case with the analysis of $`P^{}(\mathrm{}|r)>p_c`$ for $`d=3`$ (see Sec. II.C.2.ii), we cannot determine the parameters that govern the large $`t_{\text{min}}`$ behavior of $`P^{}(t_{\text{min}})`$ because of limitations on the size of the system we can simulate.
* $`L>r>\xi `$. When the connectedness length is smaller than the distance between the wells, the behavior of the system is the same as a homogeneous system . This can be seen in Fig. 10a in which we plot $`P(t_{\text{min}}|r)`$ for various values of $`r`$ at $`p=0.6`$. As $`r`$ increases from below the connectedness length to above the connectedness length, the form of the distribution changes from the power law distribution of Eq. (45) to a distribution with a pronounced peak, a characteristic of homogeneous systems. In Fig. 10b, in order to eliminate the finite-size effect, we select $`L=r+2`$ so that the distribution $`P(t|r)`$ does not have a power-law regime, even for small $`r`$. In this case, as shown in Fig. 10c, the fractal dimension of the minimal traveling time crosses over from $`d_{\text{tm}}=1.33`$ to $`d_{\text{tm}}=2.0`$, characteristic of a homogeneous system . The same considerations that we use to derive the behavior of the mean and variance of the shortest path can be applied to the mean and variance of the minimal time. At the moment of breakthrough, i.e., when the first tracer particle reaches the second well, the part of the system filled with tracer particles consists of $`n_b=(r/\xi )^d`$ independent blobs, each of which having a certain number of bonds $`(t_{\text{min}})_b`$ with an average $`(t_{\text{min}})_b=\xi ^{d_{\text{tm}}}`$ and a variance $`\sigma _b^2=\xi ^{2d_{\text{tm}}}`$. Thus the average minimal time for the entire system scales as
$$t_{\text{min}}=n_b\xi ^{d_{\text{tm}}}=r^d\xi ^{d_{\text{tm}}d},$$
(53)
with a variance
$$\sigma ^2=n_b\xi ^{2d_{\text{tm}}}=r^d\xi ^{2d_{\text{tm}}d}.$$
(54)
The scaling plot (Fig. 10d) of $`t_{\text{min}}`$ versus $`|pp_c|`$ shows good agreement with the theoretical prediction of Eq. (53)
$`{\displaystyle \frac{t_{\text{min}}}{r^d}}`$ $`=`$ $`(pp_c)^{(dd_{\text{tm}})\nu }`$ (55)
$`=`$ $`(pp_c)^{0.89}[d=2].`$ (56)
The graph of $`\sigma `$ versus $`r`$ (see Fig. 10e) shows linear behavior of $`\sigma `$ versus $`r`$ in agreement with Eq. (54). Equation (54) also predicts that the slope of this linear dependence decays as
$$|pp_c|^{[d_{\text{tm}}(d/2)]\nu }=|pp_c|^{0.42}[d=2].$$
(57)
However the measured slope has very small variation with $`|pp_c|`$ which is beyond the accuracy of our data points.
As mentioned above, the minimal traveling time is the sum of the inverse local velocities over the fastest path where the fastest path is statistically identical to the shortest path. While the velocity distribution has been studied extensively , because the velocities along the path are correlated, the relation between the minimum traveling time distribution and the local velocity distribution is an open challenge for further research.
The analysis for three dimensions is completely analogous to that for two dimensions. Our results are shown in Figs. 11 and 12 and the scaling parameters found are included in Table I.
### C Fastest Path
We observe that the path which takes minimal time is not always the shortest path. However analysis of the distributions of $`\mathrm{}_{\text{min}}`$ yields parameters identical to those for the distribution of the shortest paths between points separated by distance $`r`$ studied in detail in Ref. . Thus, statistically, the path which takes the shortest time is one of the paths of shortest length.
In many transport problems, the characteristic time $`t^{}`$ scales with the characteristic length $`\mathrm{}^{}`$ with a power law,
$$t^{}(\mathrm{}^{})^z.$$
(58)
Since $`t^{}`$ scales as $`r^{d_t}`$ and $`\mathrm{}^{}`$ scales as $`r^{d_{\text{min}}}`$,
$$z=d_t/d_{\text{min}}.$$
(59)
Since $`t_{\text{min}}`$ and $`\mathrm{}_{\text{min}}`$ are strongly correlated, the distributions $`P(\mathrm{}_{\text{min}})`$ and $`P(t_{\text{min}})`$ satisfy
$$P(\mathrm{}_{\text{min}})d\mathrm{}_{\text{min}}=P(t_{\text{min}})dt_{\text{min}}$$
(60)
Combining Eqs. (58)-(60) and the equations for the respective distributions, we obtain a scaling relation between exponents,
$$(g_{\mathrm{}_{\text{min}}}1)d_{\mathrm{}_{\text{min}}}=(g_{\text{tm}}1)d_{\text{tm}}.$$
(61)
These scaling relations are well satisfied by the set of scaling exponents given in Table I.
## IV Conclusions
By modeling porous media by bond percolation and using concepts of percolation theory, we study the flow of fluid in porous media in 2 and 3 dimensions between two “wells” separated by Euclidean distance $`r`$. We investigate the distribution function of the shortest path connecting the two sites, and propose a scaling Ansatz that accounts for the dependence of this distribution (i) on $`L`$, the size of the system, and (ii) on $`p`$, the bond occupancy probability. The finite size of the system, $`L`$ corresponds to the size of the oil field and the bond occupancy probability, $`p`$ is related to the good sand to impermeable rock ratio of the porous media. We confirm by extensive simulations that the Ansatz holds for $`d=2,3`$, and we calculate the relevant scaling parameters.
In order to understand the properties of the flow of oil displaced by fluid or gas, we study the dynamics of flow on percolation clusters. We study two dynamical quantities: the minimal traveling time and the length of the path corresponding to the minimal traveling time. Because of the approximate parallel between the shortest path and the minimal traveling time of flow, the study of the shortest path is the first step in understanding the properties of the oil fields. In particular, a scaling Ansatz for these dynamical quantities includes the effect of finite system size and off-critical bond occupation probability. We find that the scaling form for the distribution functions for these dynamical quantities for $`d=2,3`$ is similar to, but not identical to, that for the shortest path. In addition to calculating the relevant distribution functions and scaling relations, we also calculate a number of new exponents (see Table I).
## V acknowledgments
We thank BP Amoco for financial support, and M. Barthélémy, A. Coniglio, J. Koplik, S. Redner, and R. M. Ziff for helpful discussions.
## A Sampling all size clusters versus sampling only infinite clusters
Since $`P^{}(\mathrm{}|r)`$ is the distribution of shortest paths between point in clusters of all sizes (not just the infinite cluster), we can write
$$P^{}(\mathrm{}|r)=_0^{\mathrm{}}P(\mathrm{}|r,M)P(M)𝑑M,$$
(A1)
where $`P(\mathrm{}|r,M)`$ is the distribution of shortest paths in clusters of mass $`M`$.
The size $`L`$ of a cluster of mass $`M`$ is of the order $`M^{1/d_f}`$ and the shortest path between two points in the cluster is larger than $`L^{d_{\text{min}}}`$ drops off exponentially, so $`P(\mathrm{}|r,M)`$ will effectively be zero when $`M`$ is less than $`\mathrm{}^{d_f/d_{\text{min}}}`$. On the other hand, when $`M`$ is greater than this value, $`P(\mathrm{}|r,M)`$ will be the same distribution as the distribution for the case when the two points are in an infinite cluster. This can be taken into account by replacing $`P(\mathrm{}|r,M)`$ with $`P(\mathrm{}|r)`$ and replacing the zero lower limit of the integral by $`\mathrm{}^{d_f/d_{\text{min}}}`$. Since $`P(M)M^\tau `$, where $`\tau =d/d_f`$, we have
$`P^{}(\mathrm{}|r)`$ $`=`$ $`{\displaystyle _{\mathrm{}^{d_f/d_{\text{min}}}}^{\mathrm{}}}P(\mathrm{}|r)M^{d/d_f}𝑑M`$ (A2)
$`=`$ $`P(\mathrm{}|r)\mathrm{}^{(d_fd)/d_{\text{min}}},`$ (A3)
from which follows
$$g_{\mathrm{}}^{}=g_{\mathrm{}}+\frac{dd_f}{d_{\text{min}}}=g_{\mathrm{}}+(dd_f)\stackrel{~}{\nu }.$$
(A4)
This can be generalized for any quantity, $`x`$, with fractal dimension $`d_x`$ and with distributions
$$P(x|r)=\frac{1}{r^{d_x}}\left(\frac{x}{r^{d_x}}\right)^{g_x}$$
(A5)
and
$$P^{}(x|r)=\frac{1}{r^{d_x}}\left(\frac{x}{r^{d_x}}\right)^{g_x^{}},$$
(A6)
where “infinite clusters” and all finite-sized clusters, respectively, are sampled. This generalization results in
$$g_x^{}=g_x+\frac{dd_f}{d_x}.$$
(A7)
Table I. Summary of exponents and coefficients in scaling form
$`P(x|r)(1/r^{d_x})(x/r^{d_x})^{g_x}f_1(x/r^{d_x})f_2(x/L^{d_x})f_3(x/\xi ^{d_x})`$
where $`f_1(y)=\mathrm{exp}(a_xy^{\varphi _x})`$, $`f_2(y)=\mathrm{exp}(b_xy^{\psi _x})`$, $`f_3(y)=\mathrm{exp}(c_xy)`$. Here $`x`$ denotes one of the quantities, $`\mathrm{}`$ or $`t_{\text{min}}`$. The notation N/A means Not Applicable (since no theoretical value exists), while the notation (+/-) indicates above or below $`p_c`$.
(a) $`d=2`$ results
| $`x`$ | $`\mathrm{}`$ | | $`t_{\text{min}}`$ | |
| --- | --- | --- | --- | --- |
| exponent | sim. | theory | sim. | theory |
| $`d_x`$ | $`1.13\pm 0.01`$ | $`N/A`$ | $`1.33\pm 0.05`$ | $`N/A`$ |
| $`g_x^{}`$ | $`2.14\pm 0.02`$ | 2.11 | $`2.0\pm 0.1`$ | $`N/A`$ |
| $`a_x`$ | $`0.5`$ | $`N/A`$ | $`1.1`$ | $`N/A`$ |
| $`\varphi _x`$ | $`7.3\pm 0.5`$ | $`1/(d_x1)=7.69`$ | $`3.0`$ | $`3.0`$ |
| $`b_x`$ | $`3.5`$ | $`N/A`$ | $`5.0`$ | $`N/A`$ |
| $`\psi _x`$ | $`4.0\pm 0.5`$ | $`N/A`$ | $`3.0`$ | $`N/A`$ |
| $`c_x`$ | $`2.4(),3.7(+)`$ | $`N/A`$ | $`1.6(),2.6(+)`$ | $`N/A`$ |
(b) $`d=3`$ results
| $`x`$ | $`\mathrm{}`$ | | $`t_{\text{min}}`$ | |
| --- | --- | --- | --- | --- |
| exponent | sim. | theory | sim. | theory |
| $`d_x`$ | $`1.39\pm 0.05`$ | $`N/A`$ | $`1.45\pm 0.1`$ | $`N/A`$ |
| $`g_x^{}`$ | $`2.3\pm 0.1`$ | $`2.23`$ | $`2.1\pm 0.1`$ | $`N/A`$ |
| $`a_x`$ | $`1.4`$ | $`N/A`$ | $`2.5`$ | $`N/A`$ |
| $`\varphi _x`$ | $`2.1\pm 0.5`$ | $`1/(d_x1)=2.56`$ | $`1.6`$ | $`2.0`$ |
| $`b_x`$ | $`2.0`$ | $`N/A`$ | $`2.3`$ | $`N/A`$ |
| $`\psi _x`$ | $`2.5\pm 0.5`$ | $`N/A`$ | $`2.0`$ | $`N/A`$ |
| $`c_x`$ | $`3.1()`$ | $`N/A`$ | $`2.9()`$ | $`N/A`$ |
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# The Semiclassical Propagator for Spin Coherent States
## I Introduction
Coherent-state path integrals for spin were introduced by Klauder, and by Kuratsuji and Suzuki . Related phase space path integrals were introduced by Jevicki and Papanicolaou, and by Nielsen and Röhrlich. For a review see . These path integrals have attracted attention in connection with geometric quantization, and for providing examples hinting at possible infinite-dimensional extensions of the Duistermat-Heckman theorem on conditions for the exactness of the stationary phase approximation. Perhaps their most significant practical applications, however, have been in computations of spin tunnelling in the semiclassical limit. Here the spin path-integral formalism gives a good qualitative description of the tunneling process, including the simplest and most vivid picture of the topological quenching of spin tunneling that has recently been seen in the magnetic molecule Fe<sub>8</sub>. When we require precise quantitative results, however, the spin coherent-state path integral runs into problems: A straight forward application of instanton methods to compute the tunnel splitting yields answers that are incorrect beyond the leading exponential order. A full derivation of the splitting, including the correct prefactor, has only recently been provided by Belinicher, Providencia and Providencia. These authors showed that the continuum limit of the discrete path integral is rather delicate, and in their computation the simplicity of the instanton method is lost. These difficulties have lead to the spin path integral acquiring a reputation for being unreliable—or, even worse, being meaningful only in its discrete-time form. Many workers in the field have sought alternatives to path integrals such as discrete WKB methods .
This paper is intended to effect a rehabilitation of the continuous-time spin coherent-state path integral. We advertise and explain the origin of a previously discovered, but largely unknown, correction to the naïve form of the semiclassical propagator. This “extra phase” was obtained by Solari as a result of a careful evaluation of the discrete path form of the path integral. It also appears, as a product of a manipulation, apparently carried out for convenience, in a paper by Kochetov. We derive it here by pointing out that the functional determinant resulting from the fluctuation integral about the classical path possesses an anomaly. Regulating the determinant in a manner consistent with the underlying causal structure leads to the extra contribution.
The structure of the paper is as follows: In section two we review spin coherent states built on highest- or lowest-weight spin-$`j`$ states. We focus primarily on their holomorphic properties. In section three we review the properties of the classical action that appears in the path integral for spin, stressing the importance of boundary terms in avoiding the over-determination problem. In section four we compute the gaussian integral over small fluctuations about the classical path, and obtain the extra-phase correction. In section five we verify that, once the extra contribution is taken into account, the semiclassical propagator has the correct short time behaviour. This verification is immediate at first order in $`T`$, but the agreement between our expression and the exact result at $`O(T^2)`$ provides a significant test of the correctness of our result. In section six we check the consistency of the expression for the propagator under the dissection of the path. We find that our semiclassical propagator does not pass this test unless we repartition terms between the exponent and the prefactor. This forces us to regard the large parameter in the semiclassical expansion as being $`j+1/2`$, rather than $`j`$. As a byproduct, this observation resolves the mystery of the divergent normalization factor that appears in most treatments of the semiclassical propagator. Finally, in section seven, we compute the semiclassical propagator for the hamiltonian $`\widehat{H}=\nu J_z^2`$. We confirm that our expression obtains the correct leading and next to leading terms in the large-$`j`$ expansion.
## II Spin Coherent States
We define a family of spin coherent states by
$$|z=\mathrm{exp}(z\widehat{J}_+)|j,j.$$
(1)
These states are not normalized, but have the advantage of being holomorphic in the parameter $`z`$. Consequently, matrix elements such as $`z^{}|\widehat{O}|z`$ will be holomorphic functions of the variable $`z`$, and anti-holomorphic functions of the variable $`z^{}`$.
The inner product of two of these states is
$$z^{}|z=(1+\overline{z}^{}z)^{2j},$$
(2)
and the left eigenstates $`j,m|`$ of $`\widehat{J}^2`$ and $`\widehat{J}_3`$ have coherent-state wavefunctions
$$\psi _m^{(1)}(z)j,m|z=\sqrt{\frac{2j!}{(jm)!(j+m)!}}z^{j+m}.$$
(3)
This means that a general element of the spin-$`j`$ Hilbert space may be represented by a polynomial in $`z`$ of degree $`n2j`$.
As with any family of generalized coherent states derived from a unitary irreducible representation of a compact group, Shur’s lemma provides us with an overcompleteness relation. In the present case this reads
$$\text{1}=\frac{2j+1}{\pi }\frac{d^2z}{(1+\overline{z}z)^{2j+2}}|zz|.$$
(4)
Here $`2j+1`$ appears because it is the dimension of the representation. The symbol $`d^2z`$ is shorthand for $`dxdy`$, and the factor $`1/(1+\overline{z}z)^2`$ combines with this to make the invariant measure on the coset $`SU(2)/U(1)`$. This coset is, of course, the two-sphere, $`S^2`$, equipped with stereographic coordinates. The south pole, corresponding to spin down, is at $`z=0`$, while the north pole, spin up, is at $`z=\mathrm{}`$ — the one-point compactification of the complex plane. The remaining factor in the measure, $`1/(1+\overline{z}z)^{2j}`$, serves to normalize the states.
The wavefunctions $`\psi _m^{(1)}(z)`$ are singular at the north pole, $`z=\mathrm{}`$. Indeed there is no actual state $`|\mathrm{}`$ because the phase of this putative limiting state would depend on the direction from which we approach the point at infinity. We may, however, define a second family of states
$$|z_2=\mathrm{exp}(z\widehat{J}_{})|j,j,$$
(5)
and form the wavefunctions
$$\psi _m^{(2)}(z)=j,m|z_2.$$
(6)
These states and wavefunctions are well defined in the vicinity of the north pole, but singular near the south pole.
To find the relation between $`\psi ^{(2)}(z)`$ and $`\psi ^{(1)}(z)`$ we note that the matrix identity
$$\left[\begin{array}{cc}1& z\\ 0& 1\end{array}\right]\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]=\left[\begin{array}{cc}1& 0\\ z^1& 1\end{array}\right]\left[\begin{array}{cc}z& 0\\ 0& z^1\end{array}\right]\left[\begin{array}{cc}1& z^1\\ 0& 1\end{array}\right],$$
(7)
coupled with the faithfulness of the spin-$`\frac{1}{2}`$ representation of $`SU(2)`$, implies the relation
$$\mathrm{exp}(z\widehat{J}_+)\widehat{w}=\mathrm{exp}(z^1\widehat{J}_{})(z)^{2\widehat{J}_3}\mathrm{exp}(z^1\widehat{J}_+),$$
(8)
where $`\widehat{w}=\mathrm{exp}(i\pi \widehat{J}_2)`$ is the generator of the Weyl group of $`SU(2)`$. We also note that
$$\widehat{w}|j,j=(1)^{2j}|j,j,\widehat{w}|j,j=|j,j.$$
(9)
Thus,
$`\psi _m^{(1)}(z)`$ $`=`$ $`j,m|e^{z\widehat{J}_+}|j,j`$ (10)
$`=`$ $`(1)^{2j}j,m|e^{z\widehat{J}_+}\widehat{w}|j,j`$ (11)
$`=`$ $`(1)^{2j}j,m|e^{z^1\widehat{J}_{}}(z)^{2\widehat{J}_3}e^{z^1\widehat{J}_+}|j,j`$ (12)
$`=`$ $`(1)^{2j}(z)^{2j}j,m|e^{z^1\widehat{J}_{}}|j,j`$ (13)
$`=`$ $`z^{2j}\psi _m^{(2)}(z^1).`$ (14)
The coherent-state wavefunctions $`\psi _m^{(1)}`$ and $`\psi _m^{(2)}`$ may therefore be regarded as composing a single global section, $`\psi _m`$, of a holomorphic line bundle with transition function $`z^{2j}`$ relating its components $`\psi _m^{(1)}(z)`$, and $`\psi _m^{(2)}(\zeta 1/z)`$ in the two coordinate patches. It is the requirement that the transition function and its inverse be holomorphic and single valued in the overlap of the coordinate patches that forces $`2j`$ to be an integer. In the sequel, all coherent states, unless otherwise specified, will drawn from the first family, $`|z`$.
The above construction is an example of the Borel-Weil realization of representations of compact groups as sections of holomorphic bundles. It serves as the paradigm for the more general theory of geometric quantization. Because global analyticity is characteristic of the minimal-spread coherent states built on highest- (or lowest-) weight states, and also serves (via the transition function) to specify the Hilbert space, it is a property that should be maintained order-by-order in any approximation scheme.
For physical interpretations we must normalize the coherent states. This we do by multiplying them by
$$N(\overline{z},z)=(1+\overline{z}z)^j.$$
(15)
For example,
$$N^2z|\widehat{J}_3|z=j\frac{\overline{z}z1}{\overline{z}z+1},\text{and}N^2z|\widehat{J}_+|z=\frac{2j\overline{z}}{\overline{z}z+1}.$$
(16)
If we recall the connection between stereographic and spherical polar coordinates,
$$z=e^{i\varphi }\mathrm{cot}\frac{\theta }{2},$$
(17)
we see that
$$j\frac{\overline{z}z1}{\overline{z}z+1}=j\mathrm{cos}\theta ,\text{and}\frac{2jz}{\overline{z}z+1}=je^{i\varphi }\mathrm{sin}\theta .$$
(18)
We also note that
$`N^2z|\widehat{J}_3^2|z`$ $`=`$ $`j(j{\displaystyle \frac{1}{2}})\left({\displaystyle \frac{\overline{z}z1}{\overline{z}z+1}}\right)^2+{\displaystyle \frac{1}{2}}j`$ (19)
$`=`$ $`j(j{\displaystyle \frac{1}{2}})\mathrm{cos}^2\theta +{\displaystyle \frac{j}{2}}.`$ (20)
Similarly
$`N^2z|\widehat{J}_1^2|z`$ $`=`$ $`j(j{\displaystyle \frac{1}{2}})\mathrm{sin}^2\theta \mathrm{cos}^2\varphi +{\displaystyle \frac{j}{2}},`$ (21)
$`N^2z|\widehat{J}_2^2|z`$ $`=`$ $`j(j{\displaystyle \frac{1}{2}})\mathrm{sin}^2\theta \mathrm{sin}^2\varphi +{\displaystyle \frac{j}{2}}.`$ (22)
Thus $`N^2z|\widehat{𝐉}^2|z=j(j+1)`$, as it should.
The normalized wavefunctions $`N(\overline{z},z)\psi _m^{(1)}(z)`$ have their maximum amplitude on the lines of latitude
$$|z|^2=|z_m|^2=\frac{j+m}{jm}$$
(23)
corresponding to the polar angle $`\theta _m=\mathrm{cos}^1m/j`$. Note that
$$N^2z_m|\widehat{J}_3|z_m=j\frac{|z_m|^21}{|z_m|^2+1}=m.$$
(24)
The variance, in terms of $`m`$, is given by
$$\left(N^2\widehat{J}_3^2N^4\widehat{J}_3^2\right)=\frac{1}{2}j\left(1\left(\frac{\overline{z}z1}{\overline{z}z+1}\right)^2\right)=\frac{1}{2}j(1\mathrm{cos}^2\theta ).$$
(25)
Since $`mj\mathrm{cos}\theta `$, the normalized wavefunctions have zonal spread $`\mathrm{\Delta }\theta 1/\sqrt{j}`$. As $`j`$ becomes large the quantum spin becomes more localized, and more classical.
## III Spin Action
We wish to find a semiclassical approximation for the propagator
$$K(\overline{\zeta }_f,\zeta _i,T)=\zeta _f|e^{i\widehat{H}T}|\zeta _i$$
(26)
in the form
$$K_{\mathrm{scl}}(\overline{\zeta }_f,\zeta _i,T)=K_{\mathrm{reduced}}\mathrm{exp}\left\{S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,T)\right\}.$$
(27)
Here $`S_{\mathrm{cl}}`$ is the action for a classical path going from the point $`z=\zeta _i`$ to the point $`z=\zeta _f`$ in time $`T`$. The action functional is expected to be that appearing in the path integral representation of the exact propagator. The amplitude $`K_{\mathrm{reduced}}`$, the pre-exponential factor, is then given by a gaussian approximation to the integral over deviations from the classical trajectory. Such a semiclassical approximation should be accurate when $`j`$ is large.
If a continuous-time path integral is “derived” by inserting $`N`$ intermediate overcompleteness relations into (26) and taking a formal limit $`N\mathrm{}`$, then we find
$$K(\overline{\zeta }_f,\zeta _i,T)=_{\zeta _i}^{\overline{\zeta }_f}𝑑\mu (\overline{z},z)\mathrm{exp}\{S(\overline{z}(t),z(t))\},$$
(28)
where the path measure $`d\mu `$ is
$$d\mu (\overline{z}(t),z(t))=\underset{N\mathrm{}}{lim}\underset{n=1}{\overset{N}{}}\frac{2j+1}{\pi }\frac{d^2z_n}{(1+\overline{z}_nz_n)^{2j+2}},$$
(29)
and the action $`S(\overline{z}(t),z(t))`$ is given by
$$S(\overline{z}(t),z(t))=j\left\{\mathrm{ln}(1+\overline{\zeta }_fz(T))+\mathrm{ln}(1+\overline{z}(0)\zeta _i)\right\}+_0^T\left\{j\frac{\dot{\overline{z}}z\overline{z}\dot{z}}{1+\overline{z}z}iH(\overline{z},z)\right\}𝑑t.$$
(30)
Here the classical hamiltonian, $`H(\overline{z},z)`$, is related to the quantum $`\widehat{H}`$ by
$$H(\overline{z},z)=z|\widehat{H}|z/z|z.$$
(31)
The paths $`z(t)`$, $`\overline{z}(t)`$ obey the boundary conditions $`z(0)=\zeta _i`$, $`\overline{z}(T)=\overline{\zeta }_f`$, but $`\overline{z}(0)`$, $`z(T)`$, being actually $`\overline{z}(0+ϵ)`$ and $`z(Tϵ)`$, are unconstrained, and are to be integrated over.
When we regard $`S`$ as the phase-space action for a classical system, the explicit boundary terms, which appear naturally in the discretized path integral, serve to ensure that both the first-order Hamilton equations and their boundary conditions are compatible with the action principle. To see this, make a general variation in the trajectory, including variations in the endpoints. We find that
$`\delta S`$ $`=`$ $`{\displaystyle \frac{2jz(T)}{1+\overline{\zeta }_fz(T)}}\delta \overline{\zeta }_f+{\displaystyle \frac{2j\overline{z}(0)}{1+\overline{z}(0)\zeta _i}}\delta \zeta _i`$ (33)
$`+{\displaystyle _0^T}\left\{\delta z(t)\left({\displaystyle \frac{2j\dot{\overline{z}}}{(1+\overline{z}z)^2}}i{\displaystyle \frac{H}{z}}\right)+\delta \overline{z}(t)\left({\displaystyle \frac{2j\dot{z}}{(1+\overline{z}z)^2}}i{\displaystyle \frac{H}{\overline{z}}}\right)\right\}𝑑t.`$
There are no boundary contributions proportional to $`\delta \overline{z}(0)`$ or $`\delta z(T)`$ because of a cancellation of such terms arising from an integration by parts against those arising from the variation of the explicit boundary terms. Equating the variation of the action to zero therefore requires the classical path to obey the Hamilton equations
$$\dot{\overline{z}}=i\frac{(1+\overline{z}z)^2}{2j}\frac{H}{z},\dot{z}=i\frac{(1+\overline{z}z)^2}{2j}\frac{H}{\overline{z}},$$
(34)
together with boundary conditions that fix $`z(0)=\zeta _i`$, and $`\overline{z}(T)=\overline{\zeta }_f`$.
The quantities $`\overline{z}(0)`$ and $`z(T)`$ are not fixed by the boundary conditions, but can be found by solving the equations of motion. If we know the action for the classical path, they can also be read off from the Hamilton-Jacobi equations that follow from (33), viz:
$$\frac{S_{\mathrm{cl}}}{\overline{\zeta }_f}=\frac{2jz(T)}{1+\overline{\zeta }_fz(T)},\frac{S_{\mathrm{cl}}}{\zeta _i}=\frac{2j\overline{z}(0)}{1+\overline{z}(0)\zeta _i}.$$
(35)
In general $`\overline{z}(0)`$ will not be the complex conjugate of $`z(0)\zeta _i`$, nor will $`z(T)`$ be the complex conjugate of $`\overline{z}(T)\overline{\zeta }_f`$. This means that if we write $`z`$ as $`x+iy`$ and $`\overline{z}=xiy`$, then, except in special cases, $`x`$ and $`y`$ are not real numbers.
The Hamilton-Jacobi relations also tell us that
$$\frac{S_{\mathrm{cl}}}{\overline{\zeta }_i}=\frac{S_{\mathrm{cl}}}{\zeta _f}=0,$$
(36)
showing that $`S_{\mathrm{cl}}`$ is a holomorphic function of $`\zeta _i`$, and an anti-holomorphic function of $`\zeta _f`$. These analyticity properties of $`S_{\mathrm{cl}}`$ coincide with those of $`K`$. This is reasonable since $`\mathrm{exp}S_{\mathrm{cl}}`$ is the leading approximation to $`K`$, and we would expect analyticity to preserved term-by-term in the large $`j`$ expansion. Finally
$$\frac{S_{\mathrm{cl}}}{T}=iH(\overline{\zeta }_f,z(T)).$$
(37)
The leading semiclassical approximation is exact when the quantum hamiltonian $`\widehat{H}`$ is an element of the Lie algebra of $`SU(2)`$. For example, if $`\widehat{H}=\omega \widehat{J}_3`$, then
$$H(\overline{z},z)=N^2z|\widehat{H}|z=\omega j\frac{\overline{z}z1}{\overline{z}z+1}$$
(38)
and
$$\frac{H}{z}=\frac{2j\omega \overline{z}}{(1+\overline{z}z)^2},\frac{H}{\overline{z}}=\frac{2j\omega z}{(1+\overline{z}z)^2},$$
(39)
The equations of motion are therefore
$$\dot{\overline{z}}=i\omega \overline{z},\dot{z}=i\omega z.$$
(40)
The solutions obeying the appropriate boundary conditions are
$$z(t)=e^{i\omega t}\zeta _i,\overline{z}(t)=e^{i\omega (tT)}\overline{\zeta }_f,$$
(41)
so
$$z(T)=e^{i\omega T}\zeta _i,\overline{z}(0)=e^{i\omega T}\overline{\zeta }_f.$$
(42)
It will only be in exceptional circumstances that $`z(T)=(\overline{\zeta }_f)^{}`$ or $`\overline{z}(0)=(\zeta _i)^{}`$.
Inserting the solutions (41) into the action we find
$`S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,T)`$ $`=`$ $`j\left\{\mathrm{ln}(1+\overline{\zeta }_f\zeta _ie^{i\omega T})+\mathrm{ln}(1+\overline{\zeta }_f\zeta _ie^{i\omega T})\right\}+{\displaystyle _0^T}\left\{ij\omega {\displaystyle \frac{2\overline{z}z}{1+\overline{z}z}}ij\omega {\displaystyle \frac{\overline{z}z1}{zz+1}}\right\}𝑑t`$ (43)
$`=`$ $`2j\mathrm{ln}(1+\overline{\zeta }_f\zeta _ie^{i\omega T})+ij\omega T.`$ (44)
This is to be compared with the exact propagator
$$K=\zeta _f|e^{i\widehat{H}T}|\zeta _i=e^{i\omega jT}(1+e^{i\omega T}\overline{\zeta }_f\zeta _i)^{2j}=\mathrm{exp}S_{\mathrm{cl}}.$$
(45)
When the hamiltonian is a more general element of the enveloping algebra (i.e. a polynomial in the generators) there will be corrections to this simple result.
## IV Fluctuation Determinant
The prefactor in the semiclassical propagator comes from integration over gaussian fluctuations about the classical trajectory. To evaluate these, we consider the second variation of the classical action, holding $`z(0)=\zeta _i`$ and $`\overline{z}(T)=\overline{\zeta }_f`$ fixed. We will write
$$S=S_{\mathrm{cl}}+\delta S+\frac{1}{2}\delta ^2S+\mathrm{},$$
(46)
where
$$\delta ^2S=i_0^T\frac{2j}{(1+\overline{z}z)^2}\left(\begin{array}{cc}\delta \overline{z}& \delta z\end{array}\right)\left[\begin{array}{cc}i_t+A& B\\ \overline{B}& i_t+A\end{array}\right]\left(\begin{array}{c}\delta z\\ \delta \overline{z}\end{array}\right)𝑑t.$$
(47)
Here
$`A`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{}{\overline{z}}}{\displaystyle \frac{(1+\overline{z}z)^2}{2j}}{\displaystyle \frac{H}{z}}+{\displaystyle \frac{}{z}}{\displaystyle \frac{(1+\overline{z}z)^2}{2j}}{\displaystyle \frac{H}{\overline{z}}}\right),`$ (48)
$`B`$ $`=`$ $`{\displaystyle \frac{}{\overline{z}}}{\displaystyle \frac{(1+\overline{z}z)^2}{2j}}{\displaystyle \frac{H}{\overline{z}}},`$ (49)
$`\overline{B}`$ $`=`$ $`{\displaystyle \frac{}{z}}{\displaystyle \frac{(1+\overline{z}z)^2}{2j}}{\displaystyle \frac{H}{z}}.`$ (50)
When $`z(t)`$, $`\overline{z}(t)`$ are the classical path, then $`\delta S=0`$.
On making a change of variables
$`\delta z`$ $`=`$ $`(1+\overline{z}z)\eta ,`$ (51)
$`\delta \overline{z}`$ $`=`$ $`(1+\overline{z}z)\overline{\eta },`$ (52)
we see that we have to compute the quadratic path integral
$$K_{\mathrm{reduced}}d[\eta ]d[\overline{\eta }]\mathrm{exp}\left\{2ij_0^T\frac{1}{2}\left(\begin{array}{cc}\overline{\eta }& \eta \end{array}\right)\left[\begin{array}{cc}i_t+A& B\\ \overline{B}& i_t+A\end{array}\right]\left(\begin{array}{c}\eta \\ \overline{\eta }\end{array}\right)𝑑t\right\}.$$
(53)
This path integral is proportional to $`\mathrm{Det}^{\frac{1}{2}}𝒟`$, where the matrix differential operator
$$𝒟=\left[\begin{array}{cc}i_t+A& B\\ \overline{B}& i_t+A\end{array}\right]=i\sigma _3_t+M$$
(54)
is subject to the boundary conditions $`\eta (0)=0`$ and $`\overline{\eta }(T)=0`$. (We will use the symbol “$`\mathrm{Det}`$” for functional determinants and “$`\mathrm{det}`$” to denote the determinant of a finite matrix. Similarly “$`\mathrm{Tr}`$” and “$`\mathrm{tr}`$”.)
There are several subtleties involved in calculating $`\mathrm{Det}𝒟`$. The most obvious is that the boundary conditions imposed on $`𝒟`$ are not in the class that make it self adjoint. Although $`𝒟`$ and $`𝒟^{}`$ are formally the same differential operator, self-adjointness requires, in addition, that their domains of definition coincide. It is not hard to see that the only boundary condition on $`𝒟`$ that leads to an identical boundary condition for $`𝒟^{}`$ is $`\eta (0)=e^{i\theta _0}\overline{\eta }(0)`$ and $`\eta (T)=e^{i\theta _T}\overline{\eta }(T)`$ for some real angles $`\theta _0`$, $`\theta _T`$. Indeed if $`B=\overline{B}=0`$, then $`𝒟`$ with our boundary conditions has no eigenfunctions — never mind a complete set. The determinant cannot be expressed as an infinite product of eigenvalues, therefore. Diagonalizability is not, however, a fundamental requirement for defining a determinant. There exists a well-defined Green function $`G=𝒟^1`$, and we should be able to obtain the determinant by varying the parameters and using the identity $`\delta \mathrm{ln}\mathrm{Det}𝒟=\mathrm{Tr}\{𝒟^1\delta 𝒟\}`$, which holds even if $`𝒟`$ is not diagonalizable.
A potential pitfall in this approach is that the variation $`\delta \mathrm{ln}\mathrm{Det}𝒟`$ is given by
$$\delta \mathrm{ln}\mathrm{Det}𝒟=\mathrm{Tr}\{𝒟^1\delta 𝒟\}=_0^T\mathrm{tr}\{G(t,t)\delta M\}𝑑t,$$
(55)
but the Green function $`G(t,t^{})`$ is discontinuous at $`t=t^{}`$. We might have a different expression for the variation depending on whether we choose to evaluate $`G(t,t)`$ as $`G(t,t+ϵ)`$ or as $`G(t,tϵ)`$. The jump in $`G`$ is, however, proportional to $`\sigma _3`$, and $`\mathrm{tr}\{\sigma _3\delta M\}0`$, so we have reason to hope that there is no actual ambiguity.
If we agree to interpret $`G(t,t)`$ as $`\frac{1}{2}(G(t,t+ϵ)+G(t,tϵ))`$, then the formal calculation is straight forward, and we merely summarize the results:
We begin by defining the matrix
$$\mathrm{\Phi }(t)=\left(\begin{array}{cc}\eta _T(t)& \eta _0(t)\\ \overline{\eta }_T(t)& \overline{\eta }_0(t)\end{array}\right).$$
(56)
Here the column vector $`(\eta _0(t),\overline{\eta }_0(t))^T`$ is a solution of $`𝒟\mathrm{\Psi }=0`$ obeying the boundary condition $`\eta _0(0)=0`$, $`\overline{\eta }_0(0)=1`$, and $`(\eta _T(t),\overline{\eta }_T(t))^T`$ is a solution with $`\eta _T(T)=1`$, $`\overline{\eta }_T(T)=0`$. The determinant of $`\mathrm{\Phi }(t)`$ is an analogue of the wronskian and is independent of $`t`$. We find that $`\mathrm{Det}𝒟=C\mathrm{det}\mathrm{\Phi }`$, where $`C`$ is some constant independent of $`H`$.
Since $`\mathrm{det}\mathrm{\Phi }`$ is time independent, we may conveniently evaluate it at $`t=T`$, where
$$C^1\mathrm{Det}𝒟=\left|\begin{array}{cc}1& \eta _0(T)\\ 0& \overline{\eta }_0(T)\end{array}\right|=\overline{\eta }_0(T),$$
(57)
or at $`t=0`$, where
$$C^1\mathrm{Det}𝒟=\left|\begin{array}{cc}\eta _T(0)& 0\\ \overline{\eta }_T(0)& 1\end{array}\right|=\eta _T(0).$$
(58)
By relaxing the conditions that $`\eta (T)=\overline{\eta }(0)=1`$, we may interpret these results in terms of the variation of the endpoints of the classical trajectory as we vary the initial points. That is
$$C^1\mathrm{Det}𝒟=\left(\frac{\overline{\eta }(0)}{\overline{\eta }(T)}\right)^1=\left(\frac{\eta (T)}{\eta (0)}\right)^1,$$
(59)
or, in terms of the original variables,
$$C^1\mathrm{Det}𝒟=\frac{1+\overline{z}(0)\zeta _i}{1+\overline{\zeta }_fz(T)}\left(\frac{\overline{z}(0)}{\overline{\zeta }_f}\right)^1=\frac{1+\overline{\zeta }_fz(T)}{1+\overline{z}(0)\zeta _i}\left(\frac{z(T)}{\zeta _i}\right)^1.$$
(60)
The equivalence of these two expressions for the determinant is not immediately obvious, but from the Hamilton-Jacobi relations
$$\frac{S_{\mathrm{cl}}}{\overline{\zeta }_f}=\frac{2jz(T)}{1+\overline{\zeta }_fz(T)},\frac{S_{\mathrm{cl}}}{\zeta _i}=\frac{2j\overline{z}(0)}{1+\overline{z}(0)\zeta _i},$$
(61)
and the equality of mixed partials of $`S_{\mathrm{cl}}`$, we obtain
$$\frac{^2S_{\mathrm{cl}}}{\zeta _i\overline{\zeta }_f}=\frac{2j}{(1+\overline{\zeta }_fz(T))^2}\frac{z(T)}{\zeta _i}=\frac{2j}{(1+\overline{z}(0)\zeta _i)^2}\frac{\overline{z}(0)}{\overline{\zeta }_f}.$$
(62)
Both expressions in (60) thus reduce to
$$C\mathrm{Det}^1𝒟=\frac{(1+\overline{\zeta }_fz(T))(1+\overline{z}(0)\zeta _i)}{2j}\frac{^2S}{\zeta _i\overline{\zeta }_f}.$$
(63)
Our calculation of the fluctuation determinant suggests, therefore, that
$$K_{\mathrm{scl}}(\overline{\zeta }_f,\zeta _i,T)\stackrel{\mathrm{?}}{=}\left(\frac{(1+\overline{\zeta }_fz(T))(1+\overline{z}(0)\zeta _i)}{2j}\frac{^2S_{\mathrm{cl}}}{\zeta _i\overline{\zeta }_f}\right)^{\frac{1}{2}}\mathrm{exp}S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,T).$$
(64)
(The proportionality constant is fixed by the requirement that this expression reduces to $`\zeta _f|\zeta _i`$ when $`T=0`$).
As indicated by the “?” over the equals sign, there are problems with this expression, and it is not quite correct.
The first problem is that, despite the optimism expressed above, there is a degree of indeterminacy in the calculation of the functional determinant. To see this, make the substitution
$`\eta (t)`$ $``$ $`e^{i\theta (t)}\eta (t),`$ (65)
$`\overline{\eta }(t)`$ $``$ $`e^{i\theta (t)}\overline{\eta }(t).`$ (66)
in the path integral (53). The measure is unchanged, but we replace $`𝒟`$ with $`\stackrel{~}{𝒟}`$, where $`\stackrel{~}{𝒟}`$ is the matrix operator $`𝒟`$ with
$`A\stackrel{~}{A}`$ $`=`$ $`A+_t\theta ,`$ (67)
$`B\stackrel{~}{B}`$ $`=`$ $`e^{2i\theta (t)}B,`$ (68)
$`\overline{B}\stackrel{~}{\overline{B}}`$ $`=`$ $`e^{2i\theta (t)}\overline{B}.`$ (69)
The value of the path integral must be unaltered by this change of integration variables, but the solution to
$$\left[\begin{array}{cc}i_t+\stackrel{~}{A}& \stackrel{~}{B}\\ \stackrel{~}{\overline{B}}& i_t+\stackrel{~}{A}\end{array}\right]\left(\begin{array}{c}\eta (t)\\ \overline{\eta }(t)\end{array}\right)=0$$
(70)
with $`\eta (0)=0`$, $`\overline{\eta }(0)=1`$ is now $`(e^{i(\theta (t)\theta (0))}\eta _0(t),e^{i(\theta (t)\theta (0))}\overline{\eta }_0(t))^T`$. The determinant, as we have calculated it, is therefore not invariant, but ends up multiplied by $`e^{i(\theta (T)\theta (0))}`$. Our expression for the functional determinant has an “anomaly” therefore.
The anomaly arises because the argument we made about the harmlessness of the discontinuity in $`G`$ depends on our defining $`G(t,t)`$ as $`G(t,t\pm ϵ)`$ with the same choice of sign in front of the $`ϵ`$ in both entries in the trace. If we examine the discrete version of path integral we see that, on the contrary, one of the entries should be evaluated with a plus, and one with a minus. Our calculation of the determinant assumed that we could interpret $`G(t,t)`$ as $`\frac{1}{2}(G(t,t+ϵ)+G(t,tϵ))`$, so our formula for the determinant is only correct if both terms in $`\mathrm{tr}\{\sigma _3\delta M\}`$ are separately zero. This will only be the case for operators $`𝒟`$ with $`A0`$. Fortunately the discrete path integral does permit the change of variables described above, and we may use this freedom to force the diagonal entries, $`\stackrel{~}{A}`$, to zero before computing the determinant. The correctly regulated functional determinant therefore differs from its naïve value by a multiplicative factor.
Including the correction to the fluctuation determinant, the semiclassical propagator becomes
$$K_{\mathrm{scl}}(\overline{\zeta }_f,\zeta _i,T)=\left(\frac{(1+\overline{\zeta }_fz(T))(1+\overline{z}(0)\zeta _i)}{2j}\frac{^2S_{\mathrm{cl}}}{\zeta _i\overline{\zeta }_f}\right)^{\frac{1}{2}}\mathrm{exp}\left\{S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,T)+\frac{i}{2}_0^TA(t)𝑑t\right\}.$$
(71)
where
$$A(\overline{z},z)=\frac{1}{2}\left(\frac{}{\overline{z}}\frac{(1+\overline{z}z)^2}{2j}\frac{H}{z}+\frac{}{z}\frac{(1+\overline{z}z)^2}{2j}\frac{H}{\overline{z}}\right),$$
(72)
is the coefficient appearing in (50)
The manoeuvre of setting $`\stackrel{~}{A}`$ to zero before evaluating the fluctuation determinant appears (although without explanation as to why it was necessary) in the previously cited paper by Kochetov that provided part of the motivation for our present work. Kochetov therefore gets the corrected expression (71). It seems, however, that the “extra phase” (it is a phase only in the simplest cases), $`\frac{i}{2}_0^TA(t)𝑑t`$, was first obtained by Solari from a careful evaluation of the discrete determinant. Solari also pointed out the necessity of a similar correction in the harmonic oscillator coherent-state path integral, which has a flat phase space. Kochetov’s discovery of the correction seems to have been independent of this earlier work.
Because of the extra phase, (71) gives the correct, indeed exact, semiclassical propagator for the case $`\widehat{H}=\omega \widehat{J}_z`$, and also for any hamiltonian consisting of (possibly time dependent) elements of the Lie algebra of $`SU(2)`$ .
## V Short Time Accuracy
The Solari-Kochetov phase also solves a second problem with (64). In contrast to the configuration space propagator, which diverges as $`T^{\frac{1}{2}}`$, the coherent-state propagator $`K(\zeta _f,\zeta _i,T)`$ is analytic in $`T`$ near $`T=0`$. This is because of the finite spread of the coherent-state wavefunctions. To first order in $`T`$ we have
$`K(\overline{\zeta }_f,\zeta _i,T)\zeta _f|e^{i\widehat{H}T}|\zeta _i`$ $``$ $`\zeta _f|\zeta _iiT\zeta _f|\widehat{H}|\zeta _i`$ (73)
$`=`$ $`\zeta _f|\zeta _i\left(1iTH(\overline{\zeta }_f,\zeta _i)\right).`$ (74)
(In the last equality we have exploited analyticity to observe that the off-diagonal $`\zeta _f|\widehat{H}|\zeta _i`$, is obtained from the diagonal $`\zeta |\widehat{H}|\zeta `$ by the the simple replacement $`\zeta \zeta _i`$, $`\overline{\zeta }\overline{\zeta }_f`$ .)
Now, from the Hamilton-Jacobi equation
$$\frac{S_{\mathrm{cl}}}{T}=iH(\overline{\zeta }_f,z(T)),$$
(75)
we have
$$S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,T)=S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,0)iTH(\overline{\zeta }_f,\zeta _i)+O(T^2),$$
(76)
while
$$S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,0)=2j\mathrm{ln}(1+\overline{\zeta }_f\zeta _i)=\mathrm{ln}\zeta _f|\zeta _i.$$
(77)
Thus, in order to get agreement between (71) and (74), the fluctuation determinant must make no $`O(T)`$ contribution to the propagator. A short calculation shows, however, that
$$\frac{(1+\overline{\zeta }_fz(T))(1+\overline{z}(0)\zeta _i)}{2j}\frac{^2S_{\mathrm{cl}}}{\zeta _i\overline{\zeta }_f}=1iTA(\overline{\zeta }_f,\zeta _i)+O(T^2).$$
(78)
Fortunately this contribution is exactly cancelled by the $`O(T)`$ contribution from the Solari-Kochetov extra phase.
We now ask how well does the semiclassical propagator do at next order in the short-time expansion. In order to provide a systematic grading for the terms, we will regard the hamiltonian $`\widehat{H}`$ as being $`O(j)`$. The entire action is then homogeneous of degree one in $`j`$. With this assumption, and by analogy with the usual semiclassical expansion in powers of $`\mathrm{}`$, we expect that
$$K(\overline{\zeta }_f,\zeta _i,T)=K_{\mathrm{reduced}}\mathrm{exp}\{S_{\mathrm{cl}}\}\left[1+O\left(\frac{1}{j}\right)\right],$$
(79)
where $`S_{\mathrm{cl}}`$ is $`O(j)`$, while the prefactor, $`K_{\mathrm{reduced}}`$, is $`O(j^0)`$.
At short time the exact coherent-state propagator is certainly of this form. To demonstrate this, expand
$$\zeta _f|e^{i\widehat{H}T}|\zeta _i=\zeta _f|\zeta _iiT\zeta _f|\widehat{H}|\zeta _i\frac{T^2}{2}\zeta _f|\widehat{H}^2|\zeta _i+\mathrm{}.$$
(80)
Now $`\zeta _f|\widehat{H}|\zeta _i=\zeta _f|\zeta _iH(\overline{\zeta }_f,\zeta _i)`$, but some work is needed to evaluate $`\zeta _f|\widehat{H}^2|\zeta _i`$.
Inserting an overcompleteness integral, we have
$`\zeta _f|\widehat{H}^2|\zeta _i`$ $`=`$ $`{\displaystyle \frac{2j+1}{\pi }}{\displaystyle \frac{d^2z}{(1+\overline{z}z)^{2j+2}}\zeta _f|\widehat{H}|zz|\widehat{H}|\zeta _i}`$ (81)
$`=`$ $`{\displaystyle \frac{2j+1}{\pi }}{\displaystyle \frac{d^2z}{(1+\overline{z}z)^{2j+2}}(1+\overline{\zeta }_fz)^{2j}(1+\overline{z}\zeta _i)^{2j}H(\overline{\zeta }_f,z)H(\overline{z},\zeta _i)}.`$ (82)
We now perform a steepest descent expansion in the integral over the intermediate states, and obtain the first three terms in its asymptotic expansion in powers of $`1/j`$. This computation is greatly simplified by using two shortcuts: First we need calculate only the diagonal matrix element $`\zeta |\widehat{H}^2|\zeta `$. Given this, we may appeal to analyticity and obtain the general matrix element by setting $`\overline{\zeta }\overline{\zeta }_f`$ and $`\zeta \zeta _i`$. Next we rotate the sphere so as to centre the coordinate system on the point $`\zeta `$. Thus $`\zeta 0`$, and the coordinate system is locally geodetic. In these coordinates the saddle point of the $`z`$ integral is at $`\zeta =z=0`$, and far fewer terms have to taken into consideration.
To return to the original coordinates, we need to be able to recognize some $`SO(3)SU(2)`$ invariant combinations of derivatives and $`(1+\overline{z}z)^2`$ factors.
One easily establishes that, under the Möbius mapping
$$zz^{}=\frac{az+b}{cz+d},\text{where}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]SU(2),$$
(83)
we have
$$\frac{d^2z}{(1+\overline{z}z)^2}=\frac{d^2z^{}}{(1+\overline{z}^{}z^{})^2}$$
(84)
together with
$`(1+\overline{z}z)^2{\displaystyle \frac{f(\overline{z},z)}{z}}{\displaystyle \frac{g(\overline{z},z)}{\overline{z}}}`$ $`=`$ $`(1+\overline{z}^{}z^{})^2{\displaystyle \frac{f(\overline{z}^{},z^{})}{z^{}}}{\displaystyle \frac{g(\overline{z}^{},z^{})}{\overline{z}^{}}},`$ (85)
$`(1+\overline{z}z)^2{\displaystyle \frac{^2f(\overline{z},z)}{z\overline{z}}}`$ $`=`$ $`(1+\overline{z}^{}z^{})^2{\displaystyle \frac{^2f(\overline{z},^{}z^{})}{z^{}\overline{z}^{}}},`$ (86)
and that the combination
$$Z=\left(\frac{}{z}(1+\overline{z}z)^2\frac{}{z}f\right)\left(\frac{}{\overline{z}}(1+\overline{z}z)^2\frac{}{\overline{z}}g\right)$$
(87)
is similarly invariant. Thus, when we see the term $`_{zz}^2f_{\overline{z}\overline{z}}^2g`$ appearing in the expansion about the stationary point $`z=0`$, we realize that in the integral for the general matrix element (where the saddle point is at $`z=\zeta _i`$, $`\overline{z}=\overline{\zeta }_f`$) we should replace it by (87).
Proceeding in this manner we find
$`\zeta _f|\widehat{H}^2|\zeta _i`$ $`=`$ $`\zeta _f|\zeta _i\{H^2(\overline{\zeta }_f,\zeta _i)+{\displaystyle \frac{(1+\overline{\zeta }_f\zeta _i)^2}{2j}}{\displaystyle \frac{H}{\zeta _i}}{\displaystyle \frac{H}{\overline{\zeta }_f}}`$ (89)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{(2j)^2}}\left({\displaystyle \frac{}{\overline{\zeta }_f}}(1+\overline{\zeta }_f\zeta _i)^2{\displaystyle \frac{H}{\overline{\zeta }_f}}\right)\left({\displaystyle \frac{}{\zeta _i}}(1+\overline{\zeta }_f\zeta _i)^2{\displaystyle \frac{H}{\zeta _i}}\right)+O\left({\displaystyle \frac{1}{j}}\right)\}.`$
The three terms in braces in this expression are of $`O(j^2)`$, $`O(j)`$, and of $`O(j^0)`$ respectively.
We may now re-exponentiate (89) as
$`\zeta _f|e^{i\widehat{H}T}|\zeta _i`$ $`=`$ $`\mathrm{exp}\left\{\mathrm{ln}\zeta _f|\zeta _iiTH(\overline{\zeta }_f,\zeta _i){\displaystyle \frac{1}{2}}T^2{\displaystyle \frac{(1+\overline{\zeta }_f\zeta _i)^2}{2j}}{\displaystyle \frac{H}{\zeta _i}}{\displaystyle \frac{H}{\overline{\zeta }_f}}+\mathrm{}\right\}`$ (91)
$`\times \left[1{\displaystyle \frac{T^2}{4}}{\displaystyle \frac{1}{(2j)^2}}\left({\displaystyle \frac{}{\zeta _i}}(1+\overline{\zeta }_f\zeta _i)^2{\displaystyle \frac{}{\zeta _i}}H\right)\left({\displaystyle \frac{}{\overline{\zeta }_f}}(1+\overline{\zeta }_f\zeta _i)^2{\displaystyle \frac{}{\overline{\zeta }_f}}H\right)+\mathrm{}\right]`$
Again using the Hamilton Jacobi equation,
$$\frac{S_{\mathrm{cl}}}{T}=iH(\overline{\zeta }_f,z(T)),$$
(92)
and the equation of motion for $`z(t)`$, we may generate the Taylor series for $`S_{\mathrm{cl}}(T)`$. We immediately verify the term in the exponential is the classical action to $`O(T^2)`$:
$$S_{\mathrm{cl}}=\mathrm{ln}\zeta _f|\zeta _iiTH(\overline{\zeta }_f,\zeta _i)\frac{1}{2}T^2\frac{(1+\overline{\zeta }_f\zeta _i)^2}{2j}\frac{H}{\zeta _i}\frac{H}{\overline{\zeta }_f}+O(T^3).$$
(93)
The expression in the square brackets in (91) must be the prefactor, and is manifestly $`O(j^0)`$. It is a little tedious to verify that our formula for the pre-exponential factor, including the Solari-Kochetov correction, reduces to exactly this, but it is so. To collapse the terms, it helps to use the identity
$`(1+\overline{z}z)^2{\displaystyle \frac{^2}{\overline{z}z}}\left((1+\overline{z}z)^2{\displaystyle \frac{H}{\overline{z}}}{\displaystyle \frac{H}{z}}\right)=`$ (94)
$`2(1+\overline{z}z)^2{\displaystyle \frac{H}{\overline{z}}}{\displaystyle \frac{H}{z}}+\left((1+\overline{z}z)^2{\displaystyle \frac{^2H}{\overline{z}z}}\right)^2+(1+\overline{z}z)^2{\displaystyle \frac{H}{\overline{z}}}{\displaystyle \frac{}{z}}\left((1+\overline{z}z)^2{\displaystyle \frac{^2H}{\overline{z}z}}\right)`$ (95)
$`+(1+\overline{z}z)^2{\displaystyle \frac{H}{z}}{\displaystyle \frac{}{\overline{z}}}\left((1+\overline{z}z)^2{\displaystyle \frac{^2H}{\overline{z}z}}\right)+\left({\displaystyle \frac{}{z}}(1+\overline{z}z)^2{\displaystyle \frac{}{z}}H\right)\left({\displaystyle \frac{}{\overline{z}}}(1+\overline{z}z)^2{\displaystyle \frac{}{\overline{z}}}H\right),`$ (96)
which is most easily established by noting that all terms are invariant, and, at $`z=0`$, both sides reduce to
$$\left(\frac{^2}{\overline{z}z}+2\right)\frac{H}{\overline{z}}\frac{H}{z}.$$
(97)
The semiclassical expression, therefore, has errors of at most $`O(j^1)`$ at short time. Our expectation is, of course, that it has this degree of accuracy uniformly in $`T`$.
## VI Consistency
A further test of the correctness of (71) is to verify its consistency under dissection of the classical trajectory. The exact propagator must satisfy the sewing condition
$$K(\overline{\zeta }_f,\zeta _i,t_1+t_2)=\frac{2j+1}{\pi }\frac{d^2\xi }{(1+\overline{\xi }\xi )^{2j+2}}K(\overline{\zeta }_f,\xi ,t_2)K(\overline{\xi },\zeta _i,t_1),$$
(98)
which follows from the definition of $`K`$ and the overcompleteness condition (4). The semiclassical approximation to $`K`$ should obey a similar condition, but with the exact integration over the intermediate states replaced by a suitable stationary phase approximation.
Since $`K_{\mathrm{scl}}\mathrm{exp}S_{\mathrm{cl}}`$, we begin with the relationship between the action for the total path from $`\zeta _i`$ to $`\zeta _f`$, and the actions for the two segments from $`\zeta _i`$ to the intermediate point $`\xi `$, and from $`\xi `$ to $`\zeta _f`$. To eliminate the redundant intermediate-point boundary terms we must define
$$S(\overline{\zeta }_f,\zeta _i,t_1+t_2)=S(\overline{\zeta }_f,\xi ,t_2)+S(\overline{\xi },\zeta _i,t_1)2j\mathrm{ln}(1+\overline{\xi }\xi ).$$
(99)
We will write this compactly as
$$S_{\mathrm{tot}}=S_2+S_12j\mathrm{ln}(1+\overline{\xi }\xi ).$$
(100)
In writing (99) we have tacitly assumed that our chosen starting $`\xi `$ of the second path segment coincides with the dynamically determined endpoint $`z(t_1)`$ of the first path segment, and that the dynamically determined starting $`\overline{z}(t_1)`$ of the second path segment coincides with our chosen $`\overline{\xi }`$ endpoint of the first path segment. This will not generally be the case — but it will be when $`\overline{\xi }`$, $`\xi `$ obey the stationary-phase equations
$$\frac{S_{\mathrm{tot}}}{\xi }=\frac{S_{\mathrm{tot}}}{\overline{\xi }}=0.$$
(101)
Taking into account the analyticity properties of $`S_1`$ and $`S_2`$, these are
$`0`$ $`=`$ $`{\displaystyle \frac{S_2(\overline{\zeta }_f,\xi )}{\xi }}{\displaystyle \frac{2j\overline{\xi }}{1+\overline{\xi }\xi }},`$ (102)
$`0`$ $`=`$ $`{\displaystyle \frac{S_1(\overline{\xi },\zeta _i)}{\overline{\xi }}}{\displaystyle \frac{2j\xi }{1+\overline{\xi }\xi }}.`$ (103)
Comparing (103) with the Hamilton-Jacobi equations confirms that $`\xi _c=z(t_1)`$ and $`\overline{\xi }_c=\overline{z}(t_1)`$, where $`\xi _c`$, $`\overline{\xi }_c`$ is the stationary phase point.
To evaluate the integral over small deviations from the classical stationary phase point, we set $`\xi =\xi _c+\eta `$, $`\overline{\xi }=\overline{\xi }_c+\overline{\eta }`$. We expand
$$S_{\mathrm{tot}}=S_{\mathrm{tot}}|_{\overline{\xi }_c,\xi _c}\frac{1}{2}\frac{2j}{(1+\overline{\xi }_c\xi _c)^2}\left(\begin{array}{cc}\overline{\eta },& \eta \end{array}\right)\left[\begin{array}{cc}1& \alpha \\ \beta & 1\end{array}\right]\left(\begin{array}{c}\eta \\ \overline{\eta }\end{array}\right),$$
(104)
where
$$\alpha =\frac{(1+\overline{\xi }_c\xi _c)^2}{2j}\frac{^2S_1}{\overline{\xi }_c^2}+\xi _c^2=\frac{1}{2j}(1+\overline{\xi }_c\xi _c)\frac{}{\overline{\xi }_c}(1+\overline{\xi }_c\xi _c)\frac{S_1}{\overline{\xi }_c},$$
(105)
and
$$\beta =\frac{(1+\overline{\xi }_c\xi _c)^2}{2j}\frac{^2S_2}{\xi _c^2}+\overline{\xi }_c^2=\frac{1}{2j}(1+\overline{\xi }_c\xi _c)\frac{}{\xi _c}(1+\overline{\xi }_c\xi _c)\frac{S_2}{\xi _c}.$$
(106)
(The second equality in these equations uses the stationary phase equations.)
We now put together two semiclassical propagators and perform the gaussian integral over the deviation from the stationary phase point. Using the semiclassical Solari-Kochetov form (71) for the propagators on the right-hand side of Eq. (98), we get (with $`T=t_1+t_2`$),
$`K_{\mathrm{comb}}`$ $`=`$ $`{\displaystyle \frac{2j+1}{\pi }}{\displaystyle \frac{d^2\eta }{(1+\overline{\xi }_c\xi _c)^2}\mathrm{exp}\left\{S_1+S_22j\mathrm{ln}(1+\overline{\xi }_c\xi _c)+\frac{i}{2}_0^TA𝑑t\frac{1}{2}\delta ^2S\right\}}`$ (108)
$`\times \left({\displaystyle \frac{(1+\overline{\zeta }_fz(T))(1+\overline{\xi }_c\xi _c)}{2j}}{\displaystyle \frac{^2S_2}{\overline{\zeta }_f\xi _c}}{\displaystyle \frac{(1+\overline{\xi }_c\xi _c)(1+\overline{z}(0)\zeta _i)}{2j}}{\displaystyle \frac{^2S_1}{\overline{\xi }_c\zeta _i}}\right)^{\frac{1}{2}}.`$
Notice that, as with consistency test of the ordinary Feynman path integral, the measure and the prefactors, including the Solari-Kochetov “extra-phase” term, are all being treated as constants. The integration involves only the variation of the classical action
$$\delta ^2S=\frac{2j}{(1+\overline{\xi }_c\xi _c)^2}\left(\begin{array}{cc}\overline{\eta },& \eta \end{array}\right)\left[\begin{array}{cc}1& \alpha \\ \beta & 1\end{array}\right]\left(\begin{array}{c}\eta \\ \overline{\eta }\end{array}\right),$$
(109)
and yields, along with other factors, the inverse square-root of the determinant
$$D=\left|\begin{array}{cc}1& \alpha \\ \beta & 1\end{array}\right|.$$
(110)
We now use the result, established in the appendix, that
$$\frac{^2S_{\mathrm{tot}}}{\overline{\zeta }_f\zeta _i}=\frac{(1+\overline{\xi }_c\xi _c)^2}{2j}\frac{^2S_2}{\overline{\zeta }_f\xi _c}\frac{^2S_1}{\overline{\xi }_c\zeta _i}\left|\begin{array}{cc}1& \alpha \\ \beta & 1\end{array}\right|^1,$$
(111)
to obtain
$$K_{\mathrm{comb}}=\left(\frac{2j+1}{2j}\right)\left(\frac{(1+\overline{\zeta }_fz(T))(1+\overline{z}(0)\zeta _i)}{2j}\frac{^2S_{\mathrm{tot}}}{\overline{\zeta }_f\zeta _i}\right)^{\frac{1}{2}}\mathrm{exp}\left\{S_{\mathrm{tot}}(\overline{\zeta }_f,\zeta _i,T)+\frac{i}{2}_0^TA𝑑t\right\}.$$
(112)
The semiclassical approximation therefore reproduces itself except for a niggling factor of $`(2j+1)/2j`$, which is due to a conflict between the normalization of the measure and the $`2j`$ appearing in the exponent.
Although this discrepant factor approaches unity in the large-$`j`$ limit, it is nonetheless disturbing. Each of the infinitely many gaussian integrations that constitute the semiclassical approximation to the path integral ought to be indistinguishable from our single gaussian integration over the intermediate point $`\xi `$. We should, therefore, be able to dissect the path into arbitrarily many parts without affecting the final answer. This is not currently so, and, in particular, the limit of large $`j`$ does not commute with the limit of a large number of intermediate points.
The origin of the discrepancy is not hard to find. In the large-$`j`$ limit the effective radius of our spherical phase space becomes large, and, near $`z=0`$, the spin-$`j`$ reproducing-kernel relation
$$\frac{2j+1}{\pi }\frac{d^2z}{(1+\overline{z}z)^2}(1+\overline{z}z)^{2j}\zeta _2|zz|\zeta _1=\zeta _2|\zeta _1,$$
(113)
or more explicitly,
$$\frac{2j+1}{\pi }\frac{d^2z}{(1+\overline{z}z)^2}(1+\overline{z}z)^{2j}(1+\overline{\zeta }_2z)^{2j}(1+\overline{z}\zeta _1)^{2j}=(1+\overline{\zeta }_2\zeta _1)^{2j},$$
(114)
should contract to a suitably scaled version of its flat-phase-space analogue
$$\frac{d^2z}{\pi }e^{\overline{z}z}e^{\overline{\zeta }_2z}e^{\overline{z}\zeta _1}=e^{\overline{\zeta }_2\zeta _1}.$$
(115)
Because it is a gaussian integral, the leading stationary phase “approximation” to (115) is exact.
If we make the obvious large $`j`$ estimates
$$(1+\overline{z}z)^{2j}e^{2j\overline{z}z},(1+\overline{\zeta }_2z)^{2j}e^{2j\overline{\zeta }_2z},(1+\overline{z}\zeta _1)^{2j}e^{2j\overline{z}\zeta _1},$$
(116)
while regarding the sphere measure $`(1+\overline{z}z)^2`$ as a prefactor, we do not get exactly
$$\frac{2j}{\pi }d^2ze^{2j\overline{z}z}e^{2j\overline{\zeta }_2z}e^{2j\overline{z}\zeta _1}=e^{2j\overline{\zeta }_2\zeta _1},$$
(117)
but instead $`(2j+1)/2j`$ times this.
If we keep terms higher order in $`1/2j`$, both those coming from the measure and those from going beyond the quadratic approximation to the exponent, they will of course correct the error. What we really need, however, is a partitioning of the integral on the LHS of (114) such that the leading steepest descent approximation will agree with the RHS. This will happen if regard the expansion parameter as $`2j+1`$ and not $`2j`$. To see this, break up
$$I=\frac{2j+1}{\pi }\frac{d^2z}{(1+\overline{z}z)^2}(1+\overline{z}z)^{2j}(1+\overline{\zeta }_2z)^{2j}(1+\overline{z}\zeta _1)^{2j}$$
(118)
as
$$I=\frac{2j+1}{\pi }\frac{d^2z}{(1+\overline{z}z)^2}g^1(\overline{z},z)e^{(2j+1)\mathrm{ln}g(\overline{z},z)}$$
(119)
with
$$g(\overline{z},z)=(1+\overline{z}z)^1(1+\overline{\zeta }_2z)(1+\overline{z}\zeta _1).$$
(120)
The critical point of the function in the exponential is at $`\overline{z}=\overline{\zeta }_2`$, $`z=\zeta _1`$, and
$$g(\overline{\zeta }_2,\zeta _1)=(1+\overline{\zeta }_2\zeta _1),$$
(121)
$$\frac{^2\mathrm{ln}g}{z\overline{z}}|_{\overline{z}=\overline{\zeta }_2,z=\zeta _1}=\frac{1}{(1+\overline{\zeta }_2\zeta _1)^2}.$$
(122)
Thus
$`I`$ $``$ $`{\displaystyle \frac{2j+1}{\pi }}{\displaystyle \frac{1}{(1+\overline{\zeta }_2\zeta _1)^3}}(1+\overline{\zeta }_2\zeta _1)^{2j+1}{\displaystyle d^2ze^{\frac{2j+1}{(1+\overline{\zeta }_2\zeta _1)^2}(\overline{z}\overline{\zeta }_2)(z\zeta _1)}}`$ (123)
$`=`$ $`{\displaystyle \frac{2j+1}{\pi }}(1+\overline{\zeta }_2\zeta _1)^{2j2}{\displaystyle \frac{\pi }{2j+1}}(1+\overline{\zeta }_2\zeta _1)^2`$ (124)
$`=`$ $`(1+\overline{\zeta }_2\zeta _1)^{2j}.`$ (125)
The leading term of the asymptotic expansion of $`I`$ in powers of $`1/(2j+1)`$ is therefore exact.
This observation suggests rewriting the semiclassical approximation to $`K`$ as
$$K_{\mathrm{scl}}(\overline{\zeta }_f,\zeta _i,T)=\frac{1}{\sqrt{2j+1}}\left(\frac{^2\stackrel{~}{S}_{\mathrm{cl}}}{\overline{\zeta }_f\zeta _i}\right)^{\frac{1}{2}}\mathrm{exp}\left\{\stackrel{~}{S}_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,T)+\frac{i}{2}_0^TQ𝑑t\right\},$$
(126)
where $`\stackrel{~}{S}_{\mathrm{cl}}=(2j+1)S_{\mathrm{cl}}/(2j)`$, and
$$Q=\frac{1}{j}\left(\frac{(1+\overline{z}z)^2}{2}\frac{^2H}{\overline{z}z}+H(\overline{z},z)\right)$$
(127)
is the term required to make (126) numerically equal to (71).
With this repartitioning of terms between the exponent and the prefactor we have exactly the same classical equations of motion, but now
$`K_{\mathrm{comb}}`$ $`=`$ $`{\displaystyle \frac{2\stackrel{~}{j}}{\pi }}{\displaystyle \frac{d^2\eta }{(1+\overline{\xi }_c\xi _c)^2}(1+\overline{\xi }_c\xi _c)\mathrm{exp}\left\{\stackrel{~}{S}_1+\stackrel{~}{S}_2(2\stackrel{~}{j})\mathrm{ln}(1+\overline{\xi }_c\xi _c)\frac{1}{2}\delta ^2\stackrel{~}{S}\right\}}`$ (129)
$`\times {\displaystyle \frac{1}{(2\stackrel{~}{j})}}\left({\displaystyle \frac{^2\stackrel{~}{S}_2}{\overline{\zeta }_f\xi _c}}{\displaystyle \frac{^2\stackrel{~}{S}_1}{\overline{\xi }_c\zeta _i}}\right)^{\frac{1}{2}}\mathrm{exp}\left\{{\displaystyle \frac{i}{2}}{\displaystyle _0^T}Q𝑑t\right\},`$
where
$$\stackrel{~}{j}=j+\frac{1}{2},$$
(130)
and
$$\delta ^2\stackrel{~}{S}=\frac{2\stackrel{~}{j}}{(1+\overline{\xi }_c\xi _c)^2}\left(\begin{array}{cc}\overline{\eta },& \eta \end{array}\right)\left[\begin{array}{cc}1& \stackrel{~}{\alpha }\\ \stackrel{~}{\beta }& 1\end{array}\right]\left(\begin{array}{c}\eta \\ \overline{\eta }\end{array}\right).$$
(131)
The quantities $`\stackrel{~}{\alpha }`$ and $`\stackrel{~}{\beta }`$ are obtained from Eqs. (105) and (106) by putting tildes on $`S_1`$, $`S_2`$, and $`j`$. Note, though, that $`\stackrel{~}{\alpha }=\alpha `$, and $`\stackrel{~}{\beta }=\beta `$. Note also, that we have inserted a factor of $`(1+\overline{\xi }_c\xi _c)`$ in the integral to compensate for the extra factor of $`(1+\overline{\xi }\xi )`$ that was taken from the measure into the exponential to complete $`\stackrel{~}{S}_{\mathrm{tot}}`$. Thus part of both the measure and the prefactor are varied in determining the stationary phase, and get integrated over, while part is regarded as a constant.
The integration in Eq. (129) can be done at once by noting that all equations in Appendix A are unchanged if we put tildes on the actions, $`j`$, $`\alpha `$, and $`\beta `$ everywhere. In particular, the identity (A16) holds with tildes. We thus obtain
$$K_{\mathrm{comb}}=\frac{1}{\sqrt{2j+1}}\left(\frac{^2\stackrel{~}{S}_{\mathrm{tot}}}{\overline{\zeta }_f\zeta _i}\right)^{\frac{1}{2}}\mathrm{exp}\left\{\stackrel{~}{S}_{\mathrm{tot}}(\overline{\zeta }_f,\zeta _i,T)+\frac{i}{2}_0^TQ𝑑t\right\},$$
(132)
all unwanted factors of $`2j+1`$ and $`(1+\overline{\xi }_c\xi _c)`$, having cancelled. Thus, with this form of stationary-phase integration, the propagator reproduces itself exactly.
What this means is that the semiclassical approximation must be tacitly using (125) in making each of the many integrations that go into the gaussian approximation to the path integral. Once we realize this, we see that there is no need for the mysterious divergent normalization factor, $`𝒩=lim_N\mathrm{}(1+1/2j)^N`$, that plagues most treatments of the semiclassical spin propagator.
The appearance of $`j+1/2`$ as the large parameter in the fluctuation integral has been remarked on before by Ercolessi et al. and by Funahashi et al.. The former worry that it is inconsistent to include fluctuations of the measure in the gaussian integral without also considering their effect in the saddle point equations. In our case all terms that are being integrated over do appear also in the equations determining the saddle point.
Note that the correction $`Q`$ vanishes for Larmor precession where $`\widehat{H}=\omega \widehat{J}_3`$. In this case, as we have seen earlier,
$$S_{\mathrm{cl}}=2j\mathrm{ln}(1+\overline{\zeta }_f\zeta _ie^{i\omega T})+ij\omega T.$$
(133)
$`\stackrel{~}{S}`$ is obtained from this by the substitution $`jj+\frac{1}{2}`$, so
$$\frac{^2\stackrel{~}{S}}{\overline{\zeta }_f\zeta _i}=e^{i\omega T}\frac{2j+1}{(1+\overline{\zeta }_f\zeta _ie^{i\omega T})^2}.$$
(134)
Thus
$`{\displaystyle \frac{1}{\sqrt{2j+1}}}\left({\displaystyle \frac{^2\stackrel{~}{S}}{\overline{\zeta }_f\zeta _i}}\right)^{\frac{1}{2}}e^{\stackrel{~}{S}(\overline{\zeta }_f,\zeta _i,T)}`$ $`=`$ $`e^{i\omega T/2}(1+\overline{\zeta }_f\zeta _ie^{i\omega T})^1(1+\overline{\zeta }_f\zeta _ie^{i\omega T})^{2j+1}e^{i\omega (j+\frac{1}{2})T}`$ (135)
$`=`$ $`e^{i\omega T}(1+\overline{\zeta }_f\zeta _ie^{i\omega T})^{2j},`$ (136)
which is the exact answer.
## VII An Example: $`\widehat{H}=\nu \widehat{J}_3^2`$
As an application of the semiclassical formalism consider $`\widehat{H}=\nu \widehat{J}_3^2`$. This hamiltonian is time reversal invariant, and we might worry that a hidden shift $`jj+1/2`$ would compromise the Kramers degeneracy expected when $`j`$ is half integral.
The classical hamiltonian corresponding to $`\widehat{H}=\nu \widehat{J}_3^2`$ is
$$H(\overline{z},z)=\frac{z|\nu \widehat{J}_3^2|z}{z|z}=\nu \left(j(j\frac{1}{2})\left(\frac{\overline{z}z1}{\overline{z}z+1}\right)^2+\frac{1}{2}j\right).$$
(137)
This should be compared with the “naïve” classical hamiltonian
$$H_{\mathrm{naive}}=\nu j^2\left(\frac{\overline{z}z1}{\overline{z}z+1}\right)^2,$$
(138)
which is what we would get if we simply expressed the classical direction-dependent energy $`\nu j^2\mathrm{cos}^2\theta `$ in terms of the stereographic coordinates on $`S^2`$.
The hamiltonian (137) leads to the classical equations of motion
$$\dot{\overline{z}}=i\omega (\overline{z},z)\overline{z},\dot{z}=i\omega (\overline{z},z)z,$$
(139)
where, with $`\mu =\nu j(j1/2)`$,
$$\omega (\overline{z},z)=\left(\frac{2\mu }{j}\right)\left(\frac{\overline{z}z1}{\overline{z}z+1}\right).$$
(140)
Since these equations imply the time independence of the product $`\overline{z}z`$, $`\omega `$ is itself time independent and the solutions may be written down directly as
$$z(t)=e^{i\omega t}\zeta _i,\overline{z}(t)=e^{i\omega (tT)}\overline{\zeta }_f.$$
(141)
Here $`\omega `$ is to be determined by the self-consistency condition
$$\omega =\left(\frac{2\mu }{j}\right)\left(\frac{e^{i\omega T}\overline{\zeta }_f\zeta _i1}{e^{i\omega T}\overline{\zeta }_f\zeta _i+1}\right).$$
(142)
As we will see below, this equation has an infinite family of solutions. Here, we wish to consider how various quantities scale with $`j`$. By demanding that Eqs. (139) continue to be meaningful as $`j\mathrm{}`$, we see that we must have $`\mu =O(j)`$, $`\omega =O(1)`$, and $`\nu =O(1/j)`$.
The classical action for the solution (141) is
$`S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,T)`$ $`=`$ $`2j\mathrm{ln}(1+e^{i\omega T}\overline{\zeta }_f\zeta _i)`$ (144)
$`+{\displaystyle _0^T}\left\{j\left({\displaystyle \frac{2i\omega e^{i\omega T}\overline{\zeta }_f\zeta _i}{1+e^{i\omega T}\overline{\zeta }_f\zeta _i}}\right)i\mu \omega ^2\left({\displaystyle \frac{j}{2\mu }}\right)^2{\displaystyle \frac{i}{2}}j\nu \right\}𝑑t`$
$`=`$ $`2j\mathrm{ln}(1+e^{i\omega T}\overline{\zeta }_f\zeta _i)+iT\{j\omega +{\displaystyle \frac{j^2}{4\mu }}\omega ^2{\displaystyle \frac{1}{2}}j\nu \}.`$ (145)
The apparently cosmetic rewrite in the last line leads to a useful way of looking at the problem. Define
$$S_\omega (\overline{\zeta }_f,\zeta _i,T)=2j\mathrm{ln}(1+e^{i\omega T}\overline{\zeta }_f\zeta _i)+iT\{j\omega +\frac{j^2}{4\mu }\omega ^2\},$$
(146)
where we regard $`\omega `$ as an independent variable. The equation
$$\frac{S_\omega (\overline{\zeta }_f,\zeta _i,T)}{\omega }=iTj\left\{\frac{e^{i\omega T}\overline{\zeta }_f\zeta _i1}{e^{i\omega T}\overline{\zeta }_f\zeta _i+1}+\left(\frac{j}{2\mu }\right)\omega \right\}$$
(147)
then shows that the consistency condition on $`\omega `$ is equivalent to $`S_\omega /\omega =0`$. We can also use $`S_\omega (\overline{\zeta }_f,\zeta _i,T)`$ to express the second variation of $`S_{\mathrm{cl}}`$ required for the prefactor $`A`$. By differentiating the Jacobi equation (35) we have
$$\frac{^2S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,T)}{\overline{\zeta }_f\zeta _i}=\frac{2j}{(1+\overline{\zeta }_fz(T))^2}\frac{z(T)}{\zeta _i},$$
(148)
and from this we find, with Eq. (141), that
$$\frac{^2S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,T)}{\overline{\zeta }_f\zeta _i}=\frac{2j}{(1+\overline{\zeta }_fz(T))^2}\left\{e^{i\omega T}+e^{i\omega T}\zeta _i\left(iT\frac{\omega }{\zeta _i}\right)\right\}.$$
(149)
We now differentiate the condition $`S_\omega /\omega =0`$ with respect to $`\zeta _i`$. This yields
$$\frac{^2S_\omega }{\zeta _i\omega }+\frac{^2S_\omega }{\omega ^2}\frac{\omega }{\zeta _i}=0.$$
(150)
Using this result to eliminate $`(\omega /\zeta _i)`$ in Eq. (149), we find, after a little algebra, that
$$\frac{^2S_{\mathrm{cl}}(\overline{\zeta }_f,\zeta _i,T)}{\overline{\zeta }_f\zeta _i}=\frac{2je^{i\omega T}}{(1+\overline{\zeta }_fz(T))^2}\frac{iTj^2}{2\mu }\left(\frac{^2S_\omega }{\omega ^2}\right)^1.$$
(151)
Substituting Eqs. (141), (145), and (151) into the basic semiclassical form (71) for the propagator, we obtain
$$K_{\mathrm{scl}}=\underset{\omega }{}\left(\frac{iTj^2}{2\mu }\right)^{1/2}\left(\frac{^2S_\omega }{\omega ^2}\right)^{1/2}\mathrm{exp}\left\{S_\omega \frac{iT}{2}(\omega +j\nu )+\frac{i}{2}_0^TA𝑑t\right\}.$$
(152)
The sum over $`\omega `$ is to be performed over all solutions to Eq. (142).
The utility of $`S_\omega (\overline{\zeta }_f,\zeta _i,T)`$ is not hard to understand. We are trying to evaluate
$$\zeta _f|e^{i\nu \widehat{J}_3^2T}|\zeta _i=\underset{m=j}{\overset{m=j}{}}(\overline{\zeta }_f\zeta _i)^{j+m}\frac{2j!}{(j+m)!(jm)!}e^{i\nu m^2T},$$
(153)
while we already know that
$`\zeta _f|e^{i\omega \widehat{J}_3T}|\zeta _i`$ $`=`$ $`{\displaystyle \underset{m=j}{\overset{m=j}{}}}(\overline{\zeta }_f\zeta _i)^{j+m}{\displaystyle \frac{2j!}{(j+m)!(jm)!}}e^{i\omega mT}`$ (154)
$`=`$ $`(1+e^{i\omega T}\overline{\zeta }_f\zeta _i)^{2j}e^{i\omega jT}`$ (155)
$`=`$ $`\mathrm{exp}S_{\omega 0}(\overline{\zeta }_f,\zeta _i,T),`$ (156)
where
$$S_{\omega 0}(\overline{\zeta }_f,\zeta _i,\omega )=2j\mathrm{ln}(1+e^{i\omega T}\overline{\zeta }_f\zeta _i)+iTj\omega .$$
(157)
From the identity
$$e^{i\nu m^2T}=e^{i\frac{\pi }{4}}\sqrt{\frac{T}{4\pi \nu }}_{\mathrm{}}^{\mathrm{}}𝑑\omega e^{i\omega mT}e^{i\omega ^2T/4\nu }$$
(158)
we have the exact relation
$`\zeta _f|e^{i\nu \widehat{J}_3^2T}|\zeta _i`$ $`=`$ $`e^{i\frac{\pi }{4}}\sqrt{{\displaystyle \frac{T}{4\pi \nu }}}{\displaystyle 𝑑\omega \zeta _f|e^{i\omega \widehat{J}_3}|\zeta _ie^{i\frac{\omega ^2T}{4\nu }}}`$ (159)
$`=`$ $`e^{i\frac{\pi }{4}}\sqrt{{\displaystyle \frac{T}{4\pi \nu }}}{\displaystyle 𝑑\omega \mathrm{exp}\left\{S_{\omega 0}(\overline{\zeta }_f,\zeta _i,\omega )+i\frac{\omega ^2T}{4\nu }\right\}}`$ (160)
$`=`$ $`e^{i\frac{\pi }{4}}\sqrt{{\displaystyle \frac{T}{4\pi \nu }}}{\displaystyle 𝑑\omega \mathrm{exp}\left\{2j\mathrm{ln}(1+e^{i\omega T}\overline{\zeta }_f\zeta _i)+iT\{j\omega +\frac{\omega ^2}{4\nu }\}\right\}}.`$ (161)
Given the form of the classical action (145), that $`\mu j^2\nu `$, and the occurrence of $`(^2S_\omega /\omega ^2)^{1/2}`$ in the prefactor, it is clear that the semiclassical approximation is attempting a stationary phase approximation to this integral over $`\omega `$. That this approximation is indeed indicated can be seen by evaluating $`(^2S_\omega /\omega ^2)`$. From Eqs. (147) and (142), we find
$$\frac{^2S_\omega }{\omega ^2}=\frac{iTj^2}{2\mu }\frac{1}{2}jT^2\left(1\frac{j^2\omega ^2}{4\mu ^2}\right),$$
(162)
which scales as $`j`$ as $`j\mathrm{}`$.
We now write the exponent in Eq. (161) as $`S_\omega iTj\omega ^2/8\mu `$. Since the second term is $`O(j^0)`$ as $`j\mathrm{}`$, we may regard it as part of the pre-exponential factor in carrying out the stationary phase integral. In this way, we obtain
$$K_{\mathrm{exact}}\underset{\omega }{}\left(\frac{iT}{2\nu }\right)^{1/2}\left(\frac{^2S_\omega }{\omega ^2}\right)^{1/2}\mathrm{exp}\left\{S_\omega iT\frac{j\omega ^2}{8\mu }\right\}.$$
(163)
The pre-exponential factors in the preceding equation agree with those in Eq. (152) to terms of order unity. To see whether the exponents agree, we must discuss the effect of the Solari-Kochetov phase. We find that
$`A`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{}{\overline{z}}}{\displaystyle \frac{(1+\overline{z}z)^2}{2j}}{\displaystyle \frac{H}{z}}+{\displaystyle \frac{}{z}}{\displaystyle \frac{(1+\overline{z}z)^2}{2j}}{\displaystyle \frac{H}{\overline{z}}}\right)`$ (164)
$`=`$ $`\omega +{\displaystyle \frac{4\mu }{j}}{\displaystyle \frac{\overline{z}z}{(1+\overline{z}z)^2}}`$ (165)
$`=`$ $`\left(\omega +{\displaystyle \frac{\mu }{j}}\right){\displaystyle \frac{j\omega ^2}{4\mu }}.`$ (166)
The term in parentheses serves to cancel \[up to $`O(1)`$\] the second term in the exponent in Eq. (152), and the $`j\omega ^2/4\mu `$ term serves to correct $`S_\omega `$ as needed in Eq. (163). Thus our semiclassical formula is indeed accurate up to $`O(1)`$ as $`j\mathrm{}`$, and we may be confident that spectral properties (Kramers degeneracy in particular) derived from it by constructing, say, the Green’s function or density of states, will be faithfully given.
Having demonstrated the formal equivalence of $`K_{\mathrm{scl}}`$ and $`K_{\mathrm{exact}}`$, we turn to the actual nature of the solution. Let us first rewrite the self-consistency condition (142) as
$$2i\stackrel{~}{\omega }\tau +\mathrm{ln}\left(\frac{1+\stackrel{~}{\omega }}{1\stackrel{~}{\omega }}\right)=\mathrm{ln}\alpha ,$$
(167)
where $`\stackrel{~}{\omega }=j\omega /2\mu `$, $`\tau =\mu T/j`$, and $`\alpha =\overline{\zeta }_f\zeta _i`$. In the limit $`\tau \mathrm{}`$, the left hand side of (167) must remain finite, suggesting that $`\stackrel{~}{\omega }1/\tau `$. A development in powers of $`1/\tau `$ shows that we may write
$$\stackrel{~}{\omega }\frac{i}{2}\frac{\mathrm{ln}\alpha }{\tau i}\frac{1}{24}\frac{(\mathrm{ln}\alpha )^3}{\tau ^4}+O(\tau ^5).$$
(168)
Since no restriction has been placed on which branch of $`\mathrm{ln}\alpha `$ is to be taken, this solution is infinitely multivalued, as asserted above. To leading order in $`1/\tau `$, different solutions differ by additive amounts $`n\pi /\tau `$, where $`n`$ is an integer.
On the other hand, at $`\tau =0`$, Eq. (167) has a unique solution, $`\stackrel{~}{\omega }=(\alpha 1)/(\alpha +1)`$. The apparent contradiction with the earlier argument for an infinite number of solutions is resolved by noting that if, as $`\tau 0`$, we allow $`\stackrel{~}{\omega }`$ to diverge as $`1/\tau `$, the left hand side of (167) again remains finite. Another development in powers of $`\tau `$ reveals that
$$\stackrel{~}{\omega }\frac{i}{2}\frac{\mathrm{ln}(\alpha )}{\tau }\frac{2}{\mathrm{ln}(\alpha )}\frac{8i}{[\mathrm{ln}(\alpha )]^3}\tau +\mathrm{},$$
(169)
which is also multivalued on account of the infinitely many branches of $`\mathrm{ln}(\alpha )`$.
We can gain further insight into the nature of the propagator and the values of $`\omega `$ at the relevant stationary-phase points by working with initial and final states on the equator of the sphere: $`\zeta _i=e^{i\varphi _i}`$, $`\overline{\zeta }_f=e^{i\varphi _f}`$. When $`j`$ is large, the problem should be essentially equivalent to a massive particle constrained to move on a ring of circumference $`2\pi `$. If we write the hamiltonian for the latter as $`L^2/2M`$, where $`L`$ is the orbital angular momentum, and $`M`$ the mass, we expect the results for the two problems to be similar with $`M=2\nu `$.
We start by considering the propagator for Larmor precession. Employing the leading large-$`j`$ estimate
$$\frac{2j!}{(j+m)!(jm)!}\frac{2^{2j}}{\sqrt{\pi j}}e^{m^2/j},$$
(170)
and using the shorthand $`\mathrm{\Delta }\varphi =\varphi _f\varphi _i`$, we may write
$`\zeta _f|e^{i\omega \widehat{J}_3}|\zeta _i`$ $`=`$ $`{\displaystyle \underset{m=j}{\overset{m=j}{}}}(\overline{\zeta }_f\zeta _i)^{j+m}{\displaystyle \frac{2j!}{(j+m)!(jm)!}}e^{i\omega mT}`$ (171)
$``$ $`e^{ij\mathrm{\Delta }\varphi }{\displaystyle \frac{2^{2j}}{\sqrt{\pi j}}}{\displaystyle \underset{m}{}}e^{im(\mathrm{\Delta }\varphi +\omega T)}e^{m^2/j}.`$ (172)
If $`Tj^{1/2}/\omega `$, the summand will have widely varying phases over the range of $`m`$ values that contributes to the sum, $`|m|\sqrt{j}`$. By extending the sum over $`m`$ to infinity and using the Poisson summation formula (taking care that $`m`$ takes half-integer values when $`j`$ is half integral), we find
$$\zeta _f|e^{i\omega \widehat{J}_3}|\zeta _ie^{ij\mathrm{\Delta }\varphi }2^{2j}\underset{n}{}e^{\frac{j}{4}(\mathrm{\Delta }\varphi +\omega T2\pi n)^2}\times (1)^n,$$
(173)
where the $`(1)^n`$ factor is present only when $`j`$ is half integral. This form is better suited to studying the large $`j`$ limit (for fixed $`T`$). In that case, (173), regarded as a function of $`\omega `$, is sharply peaked at $`\omega =\overline{\omega }_n=(2\pi n\mathrm{\Delta }\varphi )/T`$. These are the angular frequencies that allow uniform precession between $`\varphi _i`$ and $`\varphi _f`$ in time $`T`$. We now recall that Eq. (173) is nothing but $`\mathrm{exp}(S_{\omega 0})`$. If we substitute this form into Eq. (161), and take into account the factor $`\mathrm{exp}\{i\omega ^2T/4\nu \}`$ in determing the saddle-point frequencies, we find that they become complex
$`\omega _n`$ $`=`$ $`\overline{\omega }_n\left(1{\displaystyle \frac{i}{\nu Tj}}\right)^1`$ (174)
$``$ $`\overline{\omega }_n+{\displaystyle \frac{i\overline{\omega }_n}{\nu Tj}}.`$ (175)
Not surprisingly, this is just what we found in Eq. (168). The result reflects the fact that, to move at the required speed, the hamiltonian trajectories must move off the equator. There is then no real trajectory between the classical endpoints, and we must exploit the freedom to have trajectories where $`\overline{\zeta }_fz(T)^{}`$. When $`j`$ is large, however, Hamilton’s equations provide large velocities close to the equator, and the imaginary parts of $`\omega `$ are correspondingly small. Performing the integration in Eq. (161), we find
$$e^{i\varphi _f}|e^{i\nu TJ_z^2}|e^{i\varphi _i}2^{2j}e^{ij\mathrm{\Delta }\varphi }\frac{1}{(1+ij\nu T)^{\frac{1}{2}}}\underset{n}{}e^{\frac{j}{4}(\mathrm{\Delta }\varphi 2\pi n)^2/(1+ij\nu T)}\times (1)^n,$$
(176)
where, again, the last factor is only present when $`j`$ is half-integral. This form should be compared with that for the massive particle
$$\varphi _f|e^{iL^2T/2M}|\varphi _i=\frac{1}{(2\pi iMT)^{1/2}}\underset{n}{}\mathrm{exp}\left(in\mathrm{\Phi }+i\frac{M(\mathrm{\Delta }\varphi 2n\pi )^2}{2T}\right).$$
(177)
We have incorporated an Aharonov-Bohm phase $`\mathrm{\Phi }`$ into the result. This phase should be $`\pi `$ when we compare with half-integer spins, and the resulting pairwise degeneracy of the energy levels is the particle-on-a-ring analogue of Kramers degeneracy.
The similarity between Eqs. (176) and (177) is evident. Notice how $`j`$ sets the time scale for the crossover between the large-$`T`$ regime, where the spin behaves essentially as a particle of mass $`2\nu `$ on the ring, and the short-time regime where the finite range of the coherent-state wavefunctions cuts off the $`1/\sqrt{T}`$ divergence.
Note that we have ignored the difference between $`\mu /j^2`$ and $`\nu `$ in the above comparison, since as discussed while showing the equivalence of $`K_{\mathrm{scl}}`$ and $`K_{\mathrm{exact}}`$, the error incurred is of order $`1/j^2`$ relative to the leading term in the action. The semiclassical approximation therefore correctly obtains the first two terms in the large-$`j`$ expansion.
## VIII Discussion
In the previous sections we have used the continuous-time path integral to motivate a semiclassical approximation to the coherent-state propagator for spin $`j`$. Although our derivation of the semiclassical propagator is purely formal, and the resulting expression must initially have only the status of a conjecture, we have demonstrated its correctness by verifying its short-time accuracy to $`O(T^2)`$, and checking its consistency under dissection of the path. From these two properties we may conclude that our expression is accurate to $`O(j^0)`$ uniformly in time.
In our derivation it was necessary to take into account an “anomaly” in the evaluation of the functional determinant of the Jacobi operator. This is the only place where we had to appeal to details of the discrete version of the path integral. Regulating the determinant in a manner consistent with the discrete path integral results in a correction to the naïve expression for the prefactor. This correction had been noted before, by Solari and by Kochetov, but its importance does not seem to have been widely appreciated.
We have also discussed an example where an infinite number of classical trajectories contribute to the propagator. Here we again saw how the Solari-Kochetov factor is essential in obtaining the correct result.
A calculation of the Solari-Kochetov correction to the tunnel splitting between classically degenerate spin states will be reported in a separate publication.
###### Acknowledgements.
Work at Urbana and Evanston was supported by the National Science Foundation under grants DMR-98-17941 and DMR-9616749, respectively. MS would also like to thank the TCM group at the Cavendish Laboratory, Cambridge, England, for hospitality, and the EPSRC for financial support under grant number GR/N00364.
## A Composition of path-density factors
In the this appendix we derive (111). We begin by restating the stationary phase conditions (103):
$`0`$ $`=`$ $`{\displaystyle \frac{S_2(\overline{\zeta }_f,\xi )}{\xi }}{\displaystyle \frac{2j\overline{\xi }}{1+\overline{\xi }\xi }},`$ (A1)
$`0`$ $`=`$ $`{\displaystyle \frac{S_1(\overline{\xi },\zeta _i)}{\overline{\xi }}}{\displaystyle \frac{2j\xi }{1+\overline{\xi }\xi }}.`$ (A2)
Consider how the first of these evolves as we vary $`\overline{\zeta }_f`$. We find that
$`0`$ $`=`$ $`{\displaystyle \frac{}{\overline{\zeta }_f}}\left({\displaystyle \frac{S_2}{\xi _c}}{\displaystyle \frac{2j\overline{\xi }_c}{1+\overline{\xi }_c\xi _c}}\right)`$ (A3)
$`=`$ $`{\displaystyle \frac{^2S_2}{\overline{\zeta }_f\xi _c}}+{\displaystyle \frac{^2S_2}{\xi _c^2}}{\displaystyle \frac{\xi _c}{\overline{\zeta }_f}}+{\displaystyle \frac{2j\overline{\xi }_c^2}{(1+\overline{\xi }_c\xi _c)^2}}{\displaystyle \frac{\xi _c}{\overline{\zeta }_f}}{\displaystyle \frac{2j}{(1+\overline{\xi }_c\xi _c)^2}}{\displaystyle \frac{\overline{\xi }_c}{\overline{\zeta }_f}}`$ (A4)
$`=`$ $`{\displaystyle \frac{^2S_2}{\overline{\zeta }_f\xi _c}}+{\displaystyle \frac{\xi _c}{\overline{\zeta }_f}}\left({\displaystyle \frac{^2S_2}{\xi _c^2}}+{\displaystyle \frac{2j\overline{\xi }_c^2}{(1+\overline{\xi }_c\xi _c)^2}}\right){\displaystyle \frac{2j}{(1+\overline{\xi }_c\xi _c)^2}}{\displaystyle \frac{\overline{\xi }_c}{\overline{\zeta }_f}}.`$ (A5)
In the last line, we recognize the expression in parentheses to be $`2j\beta /((1+\overline{\xi }_c\xi _c)^2`$, where $`\beta `$ is the coefficient appearing in (104). By differentiating each of the two stationary phase conditions with respect to $`\overline{\zeta }_f`$ and $`\zeta _i`$, we get a total of four such equations. These may be summarized as
$$\left(\begin{array}{cc}1& \alpha \\ \beta & 1\end{array}\right)\left(\begin{array}{cc}\frac{\xi _c}{\zeta _i}& \frac{\xi _c}{\overline{\zeta }_f}\\ \frac{\overline{\xi }_c}{\zeta _i}& \frac{\overline{\xi }_c}{\overline{\zeta }_f}\end{array}\right)=\frac{(1+\overline{\xi }_c\xi _c)^2}{2j}\left(\begin{array}{cc}\frac{^2S_1}{\overline{\xi }_c\zeta _i}& 0\\ 0& \frac{^2S_2}{\overline{\zeta }_f\xi _c}\end{array}\right).$$
(A6)
Taking determinants, we obtain
$$\left|\begin{array}{cc}1& \alpha \\ \beta & 1\end{array}\right|\left|\begin{array}{cc}\frac{\xi _c}{\zeta _i}& \frac{\xi _c}{\overline{\zeta }_f}\\ \frac{\overline{\xi }_c}{\zeta _i}& \frac{\overline{\xi }_c}{\overline{\zeta }_f}\end{array}\right|=\frac{(1+\overline{\xi }_c\xi _c)^4}{(2j)^2}\frac{^2S_1}{\overline{\xi }_c\zeta _i}\frac{^2S_2}{\overline{\zeta }_f\xi _c},$$
(A7)
We now recall that the gaussian integration in Eq. (108) leads to the inverse-square root of the precisely the first determinant in Eq. (A7). This equation expresses this determinant in terms of the second derivatives of $`S_1`$ and $`S_2`$, and the jacobian $`(\xi _c,\overline{\xi }_c)/(\zeta _i,\overline{\zeta }_f)`$. The derivatives of $`S_1`$ and $`S_2`$ will cancel with the prefactors in Eq. (108), leaving only the jacobian. We therefore turn to its evaluation, and show that it can be written in terms of the second derivatives of $`S_{\mathrm{tot}}`$ with respect to $`\overline{\zeta }_f`$ and $`\zeta _i`$. We express $`S_{\mathrm{tot}}`$ as
$$S_{\mathrm{tot}}=S_2+S_12j\mathrm{ln}(1+\overline{\xi }\xi ),$$
(A8)
and take note of the fact that both $`\xi _c`$ and $`\overline{\xi }_c`$ vary as we vary $`\overline{\zeta }_f`$ and $`\zeta _i`$. We have
$`{\displaystyle \frac{^2S_{\mathrm{tot}}}{\overline{\zeta }_f\zeta _i}}`$ $`=`$ $`{\displaystyle \frac{}{\overline{\zeta }_f}}\left({\displaystyle \frac{S_2}{\xi _c}}{\displaystyle \frac{\xi _c}{\zeta _i}}{\displaystyle \frac{2j\overline{\xi }_c}{1+\overline{\xi }\xi }}{\displaystyle \frac{\xi _c}{\zeta _i}}+{\displaystyle \frac{S_1}{\overline{\xi }_c}}{\displaystyle \frac{\overline{\xi }_c}{\zeta _i}}+{\displaystyle \frac{S_1}{\zeta _i}}{\displaystyle \frac{2j\xi _c}{1+\overline{\xi }\xi }}{\displaystyle \frac{\overline{\xi }_c}{\zeta _i}}\right)`$ (A9)
$`=`$ $`{\displaystyle \frac{}{\overline{\zeta }_f}}\left({\displaystyle \frac{\xi _c}{\zeta _i}}\left\{{\displaystyle \frac{S_2}{\xi _c}}{\displaystyle \frac{2j\overline{\xi }_c}{1+\overline{\xi }\xi }}\right\}+\left({\displaystyle \frac{S_1}{\zeta _i}}\right)_{\overline{\xi }_c}+\left\{{\displaystyle \frac{S_1}{\overline{\xi }_c}}{\displaystyle \frac{2j\xi _c}{1+\overline{\xi }\xi }}\right\}{\displaystyle \frac{\overline{\xi }_c}{\zeta _i}}\right).`$ (A10)
The expressions in braces in the last line are the stationary phase conditions, so they are zero, as are their derivatives. Thus:
$$\frac{^2S_{\mathrm{tot}}}{\overline{\zeta }_f\zeta _i}=\frac{}{\overline{\zeta }_f}\left(\frac{S_1}{\zeta _i}\right)_{\overline{\xi }_c}=\frac{^2S_1}{\overline{\xi }_c\zeta _i}\frac{\overline{\xi }_c}{\overline{\zeta }_f}.$$
(A11)
Taking note of the fact that the derivative of $`S_1`$ with respect to $`\zeta _i`$ is at fixed $`\overline{\xi }_c`$, while we have useful expressions for the derivative including the variation of $`\overline{\xi }_c`$, we interchange the order of differentiation, and write
$`{\displaystyle \frac{^2S_{\mathrm{tot}}}{\overline{\zeta }_f\zeta _i}}`$ $`=`$ $`{\displaystyle \frac{\overline{\xi }_c}{\overline{\zeta }_f}}\left({\displaystyle \frac{}{\zeta _i}}\left({\displaystyle \frac{S_1}{\overline{\xi }_c}}\right){\displaystyle \frac{^2S_1}{\overline{\xi }_c^2}}{\displaystyle \frac{\overline{\xi }_c}{\zeta _i}}\right)`$ (A12)
$`=`$ $`{\displaystyle \frac{\overline{\xi }_c}{\overline{\zeta }_f}}{\displaystyle \frac{}{\zeta _i}}\left({\displaystyle \frac{2j\xi _c}{1+\overline{\xi }_c\xi _c}}\right){\displaystyle \frac{\overline{\xi }_c}{\overline{\zeta }_f}}{\displaystyle \frac{^2S_1}{\overline{\xi }_c^2}}{\displaystyle \frac{\overline{\xi }_c}{\zeta _i}}`$ (A13)
$`=`$ $`{\displaystyle \frac{\overline{\xi }_c}{\overline{\zeta }_f}}{\displaystyle \frac{}{\zeta _i}}\left({\displaystyle \frac{2j\xi _c}{1+\overline{\xi }_c\xi _c}}\right){\displaystyle \frac{\overline{\xi }_c}{\zeta _i}}{\displaystyle \frac{}{\overline{\zeta }_f}}\left({\displaystyle \frac{2j\xi _c}{1+\overline{\xi }_c\xi _c}}\right)`$ (A14)
$`=`$ $`{\displaystyle \frac{2j}{(1+\overline{\xi }_c\xi _c)^2}}\left({\displaystyle \frac{\overline{\xi }_c}{\overline{\zeta }_f}}{\displaystyle \frac{\xi _c}{\zeta _i}}{\displaystyle \frac{\overline{\xi }_c}{\zeta _i}}{\displaystyle \frac{\xi _c}{\overline{\zeta }_f}}\right).`$ (A15)
In going from the second line to the third, we used one of the equations from (A6).
Putting this together with (A7) yields
$$\frac{^2S_{\mathrm{tot}}}{\overline{\zeta }_f\zeta _i}=\frac{(1+\overline{\xi }_c\xi _c)^2}{2j}\frac{^2S_2}{\overline{\zeta }_f\xi _c}\frac{^2S_1}{\overline{\xi }_c\zeta _i}\left|\begin{array}{cc}1& \alpha \\ \beta & 1\end{array}\right|^1$$
(A16)
which is identical to Eq. (111)
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# Quintessence and Gravitational Waves
## I Introduction
Recent astrophysical and cosmological observations such as the luminosity distance-redshift relation for the supernovae type Ia , the recent observations of the cosmic microwave background temperature anisotropies , gravitional lensing and velocity fields tend to indicate that a large fraction of the matter of the universe today is composed of matter with negative pressure (see e.g. for a comparison of the different observations). Recent analyses seem to indicate that the energy density $`\rho `$ and the pressure $`P`$ of this fluid satisfies
$$1P/\rho 0.6,$$
(1)
which is compatible with a cosmological constant $`\mathrm{\Lambda }`$ for which $`P/\rho =1`$ (see also for arguments in favour of $`P/\rho <1`$). A typical value of $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$ for its energy density in units of the critical density of the universe corresponds to an energy scale of order $`5\times 10^{47}\text{GeV}^4`$ which is very far from what is expected from high energy physics; this is the well known cosmological constant problem . To circumvent this problem different solutions have been proposed starting from the idea of a dynamical cosmological constant to lead to the class of models known as quintessence , where a spatially homogeneous scalar field $`\varphi `$ is rolling down a potential decreasing when $`\varphi `$ tends to infinity. An example of such a potential which has been widely studied is the inverse power law potential. It can be obtained from some high energy physics models, e.g. where supersymmetry is broken through fermion condensates . Recently, it has been argued that supergravity has to be taken into account since today one expects the scalar field to be of order of the Planck mass $`M_{\mathrm{Pl}}`$ and corrections to the potential appear at this energy. This leads to a better agreement with observations .
An important point about this family of models is the existence of scaling solutions (refered to as tracking solutions), i.e. such that $`\varphi `$ evolves as the scale factor of the universe at a given power. These solutions are attractors of the dynamical system describing the evolution of the scale factor and of the scalar field. This implies that the present time behaviour of the field is almost independent of its initial conditions . This property allows to address (i) the coincidence problem, i.e. the fact that $`\varphi `$ starts to dominate today and (ii) the fine tuning problem, i.e. the fact that one does not have to fine tune the initial condition of the field $`\varphi `$.
One of us extended these models to include a non-minimal coupling $`\xi \overline{R}f(\varphi )`$ between the scalar field and the scalar curvature $`\overline{R}`$ . Such a coupling term appears e.g. when quantising fields in curved spacetime and in multi-dimensional theories . It was shown that when $`f(\varphi )=\varphi ^2/2`$ tracking solutions still exist and this result was generalised to any coupling function $`f`$ and potential $`V`$ satisfying $`V(\varphi )f^n(\varphi )`$. However, such a coupling is constrained by the variations of the constants of nature which fix bounds on $`\xi `$ . A way to circumvent this problem is to consider quintessence models in the framework of scalar-tensor theories where a double attractor mechanism can occur, i.e. of the scalar-tensor theory towards general relativity and of the scalar field $`\varphi `$ towards its tracking solution.
Among all the possible observations of cosmology, gravitational waves give an insight on epochs where there was a variation of the background dynamics since every such variation affects the shape of the stochastic graviton background spectrum . We can then view our universe as containing a sea of stochastic gravitational waves from primordial origin, as predicted by most models of structure formations such as inflation (see also for a review) and topological defects scenarios . Their spectrum extends typically from $`10^{18}`$ Hz (for wavelengths of order of the size of the Hubble radius today) to $`10^{10}`$ Hz (the smallest mode that has been inflated out of the Hubble radius) and they could be detected or constrained by coming experiments such as LIGO , VIRGO (at $`10^2`$ Hz) and LISA (at $`10^4`$ Hz). Gravitational waves, which are perturbations in the metric of the universe have also an effect on the cosmic microwave (CMB) temperature anisotropy and polarisation allowing to extract information on their amplitude from the measure of the CMB anisotropies. For instance, bounds on the energy density spectrum of these cosmological gravitational waves in units of the critical density, $`\mathrm{\Omega }_{\mathrm{GW}}`$, have been obtained from the CMB
$$\frac{\mathrm{d}\mathrm{\Omega }_{\mathrm{GW}}}{\mathrm{d}\mathrm{ln}\omega }|_{10^{18}\mathrm{Hz}}10^{10}.$$
Gravitational waves are also a very good probe of the conditions in the early universe since they decouple early in its history and can help e.g. testing the initial conditions of $`\varphi `$. An example was put forward by Giovannini who showed that in a class of quintessential inflation models there was an era dominated by the scalar field $`\varphi `$ before the radiation dominated era which implies that a large part of the gravitational wave energy of order $`\mathrm{\Omega }_{\mathrm{GW}}10^6`$ (about eight orders of magnitude higher than for standard inflation) was in the GHz region. This may happen in any scenario where the inflation ends with a kinetic phase or when the dominant energy condition is violated . On the other hand, the CMB temperature fluctuations give information on the history of the gravitational waves in between the last scattering surface and today through the integrated Sachs-Wolfe effect, whereas the polarisation of the CMB radiation gives mainly information on the gravitational waves at decoupling. These three observables (energy spectrum, CMB temperature and polarisation anisotropies) are thus complementary and we aim to present here a global study of the cosmological properties of the gravitational waves.
The goals of this article are (i) to study in more details the cosmology with a non-minimal quintessence field and (ii) to study gravitational waves in this class of models. In §II we set up the general framework and describe the two potentials we shall consider. In §III we introduce and define the observable quantities associated with the gravitational waves: their energy density spectrum and their imprint on the CMB radiation anisotropies and its polarisation. In §IV, we point out the general mechanism of damping by the anisotropic stress of the radiation. In §V we discuss the parameters of the problem and investigate the tuning of the potential parameters; we also describe the evolution of the background spacetime and show that a non-minimally coupled quintessence field is a candidate for a ($`\omega <1`$)-matter. In §VII we describe the main properties of the gravitational waves. We finish in §VIII by presenting numerical results and we underline the complementarity of the different observational quantities.
This work gives a detailed study of the observational effects of gravitational waves in the framework of quintessence, including some recent developments, and allowing for non-minimal coupling. This extends the work on quintessential inflation by including the effects on the CMB. It also extends the studies on $`\mathrm{\Lambda }`$CDM to quintessence and is, as far as we know, a more complete study of the effect of gravitational waves on the CMB polarisation. We hope to show that a joint study of the gravitational wave detection experiments of the CMB experiments and of the polarisation experiments can lead to a better determination of their properties.
## II General framework
### A Background spacetime
We consider a universe described by a Friedmann-Lemaître model with Euclidean spatial sections so that the metric takes the form
$$\mathrm{d}s^2=a^2(\eta )\left[\mathrm{d}\eta ^2+\delta _{ij}\mathrm{d}x^i\mathrm{d}x^j\right]\overline{g}_{\mu \nu }\mathrm{d}x^\mu \mathrm{d}x^\nu ,$$
(2)
where $`a`$ is the scale factor and $`\eta `$ the conformal time. Greek indices run from 0 to 3 and latin indices from 1 to 3.
We assume that the matter content of the universe can be described by a mixure of matter and radiation (mainly baryons, CDM, photons and three families of massless, non-degenerate neutrinos) and a scalar field $`\varphi `$ non-minimally coupled to gravity evolving in a potential $`V(\varphi )`$ that will be described later. The action for this system is
$$S=\mathrm{d}^4x\sqrt{\overline{g}}\left[\frac{\overline{R}}{2\kappa }\xi \overline{R}f(\varphi )\frac{1}{2}_\mu \varphi ^\mu \varphi V(\varphi )+_{\mathrm{matter}}\right],$$
(3)
with $`\kappa 8\pi G`$, $`G`$ being the Newton constant, and where $`_{\mathrm{matter}}`$ is the Lagrangian of the ordinary matter which is uncoupled to the scalar field and $`f(\varphi )`$ is an arbitrary function of the scalar field that will be specified later. The action (3) can be rewritten under the interesting form
$$S=\mathrm{d}^4x\sqrt{\overline{g}}\left[\frac{\overline{R}}{2\kappa _{\mathrm{eff}}[\varphi ]}\frac{1}{2}_\mu \varphi ^\mu \varphi V(\varphi )+_{\mathrm{matter}}\right],$$
(4)
with
$$\kappa _{\mathrm{eff}}[\varphi ]\frac{\kappa }{12\xi \kappa f(\varphi )}.$$
(5)
The stress-energy tensor of the scalar field is obtained by varying its Lagrangian \[$`\xi \overline{R}f(\varphi )\frac{1}{2}_\mu \varphi ^\mu \varphi V(\varphi )`$\] to get
$`T_{\mu \nu }^{(\varphi )}=\overline{}_\mu \varphi \overline{}_\nu \varphi {\displaystyle \frac{1}{2}}\overline{g}_{\mu \nu }\overline{}_\lambda \varphi \overline{}^\lambda \varphi V(\varphi )\overline{g}_{\mu \nu }+2\xi \left[\overline{g}_{\mu \nu }\overline{}_\lambda \varphi \overline{}^\lambda \varphi \overline{}_\mu \varphi \overline{}_\nu \varphi \varphi \overline{}_\mu \overline{}_\nu \varphi +\varphi \mathrm{}\varphi \overline{g}_{\mu \nu }+\overline{G}_{\mu \nu }f(\varphi )\right]`$ (6)
where $`\overline{G}_{\mu \nu }`$ is the Einstein tensor of the metric $`\overline{g}_{\mu \nu }`$, $`\overline{}`$ its covariant derivative and $`\mathrm{}\overline{}_\mu \overline{}^\mu `$.
The equations governing the evolution of the background spacetime are then obtained by varying (3) with respect to $`\overline{g}_{\mu \nu }`$, $`\varphi `$ and the ordinary matter fields to get respectively the Friedmann equations, the Klein-Gordon equation and the fluid conservation equation
$`^2`$ $`=`$ $`{\displaystyle \frac{\kappa a^2}{3}}(\rho +\rho _\varphi ),`$ (7)
$`\dot{}^2`$ $`=`$ $`{\displaystyle \frac{\kappa a^2}{2}}(\rho +P+\rho _\varphi +P_\varphi ),`$ (8)
$`\ddot{\varphi }`$ $`+`$ $`2\dot{\varphi }+a^2{\displaystyle \frac{\mathrm{d}V}{\mathrm{d}\varphi }}+6\xi (2^2+\dot{})=0,`$ (9)
$`\dot{\rho }`$ $`=`$ $`3(\rho +P).`$ (10)
A dot denotes a derivative with respect to the conformal time and $`\dot{a}/a`$ is the comoving Hubble parameter. The matter fluid energy density $`\rho `$ and pressure $`P`$ are assumed to satisfy the equation of state $`P=\omega \rho `$. The factor $`\omega `$ varies from $`1/3`$ deep in the radiation era to $`0`$ in the matter era. The scalar field energy density $`\rho _\varphi `$ and pressure $`P_\varphi `$ are obtained from its stress-energy tensor (6) and are then explicitely given by
$`\rho _\varphi `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\dot{\varphi }^2}{a^2}}+V(\varphi )+{\displaystyle \frac{2\xi }{a^2}}\left[3^2f(\varphi )+3\dot{f}(\varphi )\right],`$ (11)
$`P_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\dot{\varphi }^2}{a^2}}V(\varphi ){\displaystyle \frac{2\xi }{a^2}}\left[(2\dot{}+3^2)f(\varphi )+\dot{f}(\varphi )+\ddot{f}(\varphi )\right].`$ (12)
We stress that the conservation equation derived from (6) reduces to the Klein-Gordon equation (9). For each matter component, $`X`$ say, we introduce the density parameter $`\mathrm{\Omega }_X`$ defined as
$$\mathrm{\Omega }_X\frac{\kappa a^2\rho _X}{3^2}.$$
(13)
To completely specify the model, we have to fix the potential $`V(\varphi )`$. Following our previous work and as discussed in the introduction we choose it to behave as
$$V(\varphi )=\mathrm{\Lambda }^4\left(\frac{\mathrm{\Lambda }}{\varphi }\right)^\alpha ,\alpha >0.$$
(14)
where $`\mathrm{\Lambda }`$ is an energy scale. As shown in , such a potential leads to the existence of tracking solutions whatever the value of $`\xi `$ and for which the scalar field behaves as a barotropic fluid of equation of state (as long as the background fluid dominates)
$$P_\varphi =\omega _\varphi \rho _\varphi \text{with}\omega _\varphi =1+\frac{\alpha (1+\omega )}{\alpha +2}.$$
(15)
We also consider another class of potentials arising when one takes supergravity into account and given by
$$\stackrel{~}{V}(\varphi )=\mathrm{\Lambda }^4\left(\frac{\mathrm{\Lambda }}{\varphi }\right)^\alpha \mathrm{exp}(\kappa \varphi ^2/2),\alpha >0.$$
(16)
The effect of the exponential term is important only at late time so that the scaling properties of the tracking solution are not affected during the matter and radiation era. However when the field starts to dominate it leads to its stabilisation which has an effect on the effective equation of state of the scalar fluid.
### B Gravitational waves
In this article, we want to focus on the properties of the gravitational waves which are tensorial perturbations. At linear order, the metric is expanded as
$$g_{\mu \nu }=\overline{g}_{\mu \nu }+f_{\mu \nu }$$
(17)
where $`f_{\mu \nu }`$ is a transverse traceless (TT) perturbation, i.e. satisfying
$$f_{00}=f_{0i}=0,f_{\mu \nu }\overline{g}^{\mu \nu }=0,\overline{}_\mu f^{\mu \nu }=0.$$
(18)
It is also useful to define the perturbation $`h_{\mu \nu }`$ by
$$f_{\mu \nu }a^2h_{\mu \nu }$$
(19)
which, from (18), satisfies
$$h_{00}=h_{0i}=0,h_{kl}\delta ^{kl}=0,_kh^{kl}=0.$$
(20)
The equation of evolution of $`h_{kl}`$ is obtained by considering the TT part of the perturbed Einstein equation (see e.g. ) which leads to
$$\ddot{h}_{kl}+2\left[\kappa _{\mathrm{eff}}\xi \dot{f}(\varphi )\right]\dot{h}_{kl}\mathrm{\Delta }h_{kl}=2\kappa Pa^2\overline{\pi }_{kl}$$
(21)
where $`\mathrm{\Delta }_i^i`$ is the Laplacian and where the anisotropic stress tensor of the matter, $`\overline{\pi }_{kl}`$, is defined as the tensor component of the matter stress-energy tensor
$$\delta T_j^iP\overline{\pi }_j^i,\overline{\pi }_k^k=_i\overline{\pi }_k^i=0.$$
(22)
The anisotropic stress of the matter fluid is dominated by the contribution of the neutrinos and of the photons and its form can be obtained by describing these relativistic fluids by a Boltzmann equation (see §III B below).
## III Observational quantities
The goal of this section is to define the observable quantities related to the gravitational waves. We start by reviewing the computation of the energy density $`\rho _{\mathrm{GW}}`$ of a stochastic background of gravitational waves and finish by describing their effect on the CMB radiation, namely we present the computation of the coefficients $`C_{\mathrm{}}`$ of the development of the angular correlation function of the CMB temperature anisotropy and polarisation.
### A Gravitational waves energy density
The definition of the gravitational waves stress-energy tensor $`t_{\mu \nu }`$ can be found in e.g. in the case of a Minkowski background spacetime, in e.g. in the case of a Friedmann-Lemaître spacetime and a general discussion can be found in e.g. .
To define the gravitational waves stress-energy tensor, we have to expand the Einstein-Hilbert action (4) to second order in the perturbation $`f_{\mu \nu }`$, which implies to develop the curvature scalar $`R`$ at second order in the perturbations (see ) to get (up to divergence terms and forgetting the contribution arising from $`_{\mathrm{matter}}`$)
$$\delta ^{(2)}S=\frac{1}{4\kappa _{\mathrm{eff}}[\varphi ]}\overline{}_\mu f_{\alpha \beta }\overline{}^\mu f^{\alpha \beta }\sqrt{\overline{g}}\mathrm{d}^4x.$$
(23)
This expression is valid whatever the background metric as long as $`f_{\mu \nu }`$ is a transverse traceless perturbation. Note that contrarily to the “standard” situation, $`\kappa `$ now depends on $`\varphi `$ because of the non-minimal coupling with the scalar field. Using the fact that $`\overline{}_\mu f_{\alpha \beta }=a^2_\mu h_{\alpha \beta }`$ we can rewrite the previous expression as
$$\delta ^{(2)}S=\frac{1}{4\kappa _{\mathrm{eff}}[\varphi ]}_\mu h_{kl}^\mu h^{kl}\sqrt{\overline{g}}\mathrm{d}^4x,$$
(24)
which assumes a Friedmann-Lemaître background and the decomposition (20). Now, we decompose $`h_{kl}`$ on its two polarisations as
$$h_{kl}=\underset{\lambda =+,\times }{}h^{(\lambda )}(\eta ,𝐱)ϵ_{kl}^{(\lambda )}(𝐱),$$
(25)
where $`ϵ_{kl}^{(\lambda )}(𝐱)`$ is the polarisation tensor defined as
$$ϵ_{kl}^{(\lambda )}(𝐱)\left(e_k^1e_l^1e_k^2e_l^2\right)\delta _\times ^\lambda +\left(e_k^1e_l^2+e_l^1e_k^2\right)\delta _+^\lambda $$
(26)
for a wave propagating along the direction $`𝐞_3`$ and where $`(𝐞_1,𝐞_2,𝐞_3)`$ is a local orthonormal basis. Since this basis and the polarisation tensor satisfy
$$e_i^ae_b^i=\delta _b^a,ϵ_{kl}^{(\lambda )}ϵ_{(\lambda ^{})}^{kl}=2\delta _\lambda ^{}^\lambda ,$$
(27)
we can rewrite the action (24) of the graviton as
$$\delta ^{(2)}S=\underset{\lambda }{}\frac{1}{2\kappa _{\mathrm{eff}}[\varphi ]}_\mu h^{(\lambda )}^\mu h^{(\lambda )}\sqrt{\overline{g}}\mathrm{d}^4x,$$
(28)
which is the action for two massless scalar fields $`h_\lambda `$ evolving in the background spacetime, as first noticed by Grishchuk . By varying this action with respect to the background metric, we then deduce the stress-energy tensor of the gravitational waves
$$t_{\mu \nu }=\frac{1}{2\kappa _{\mathrm{eff}}[\varphi ]}\underset{\lambda }{}\left(_\mu h^{(\lambda )}_\nu h^{(\lambda )}\overline{g}_{\mu \nu }_\alpha h^{(\lambda )}_\alpha h^{(\lambda )}\right).$$
(29)
If we decompose $`h^{(\lambda )}`$ in Fourier modes as
$$h^{(\lambda )}(𝐱,\eta )=\frac{\mathrm{d}^3𝐤}{(2\pi )^3}\widehat{h}^{(\lambda )}(𝐤,\eta )\text{e}^{i𝐤.𝐱},$$
(30)
we can relate $`\widehat{h}^{(\lambda )}(𝐤,\eta )`$ to its initial value $`\widehat{h}^{(\lambda )}(𝐤,\eta _{\mathrm{in}})`$ , i.e. its value deep in the radiation era (e.g. at the end of the inflationary phase) through the transfer function $`T(k,\eta )`$ by solving (21) to get
$$\widehat{h}^{(\lambda )}(𝐤,\eta )=T(k,\eta )\widehat{h}^{(\lambda )}(𝐤,\eta _{\mathrm{in}}).$$
(31)
Defining the initial power spectrum of the tensor modes as
$$\widehat{h}^{(\lambda )}(𝐤,\eta _{\mathrm{in}})\widehat{h}_{(\lambda ^{})}^{}(𝐤^{},\eta _{\mathrm{in}})k^3P_h(k)\delta (𝐤𝐤^{})\delta _\lambda ^{}^\lambda ,$$
(32)
($`\delta `$ is the Dirac distribution), we can express the space average of $`t_0^0(𝐱,\eta )`$ as
$$t_0^0(𝐱,\eta )=\frac{1}{2\kappa _{\mathrm{eff}}[\varphi ]a^2}\underset{\lambda }{}_ih^{(\lambda )}_jh^{(\lambda )}\delta ^{ij}=\frac{1}{\kappa _{\mathrm{eff}}[\varphi ]}\frac{k^2}{2\pi a^2}P_h(k)T^2(k,\eta )\mathrm{d}\mathrm{ln}(k),$$
(33)
where we used an ergodic hypothesis to replace the space average by an ensemble average. Now, since $`t_0^0(𝐱,\eta )`$ oscillates, we define the energy density of the gravitational waves by taking the average of (33) over $`n`$ periods. It follows that
$$\rho _{\mathrm{GW}}(\eta )=\frac{1}{\kappa _{\mathrm{eff}}[\varphi ]}\frac{k^2}{2\pi a^2}P_h(k)\overline{T}^2(k,\eta )\mathrm{d}\mathrm{ln}(k),$$
(34)
where $`\overline{T}(k,\eta )`$ is the root mean square of $`T(k,\eta )`$ over $`n`$ periods which is well defined as long as the amplitude of the wave varies slowly with respect to its period.
The energy density $`\rho _{\mathrm{GW}}`$ and energy density parameter $`\mathrm{\Omega }_{\mathrm{GW}}`$ by frequency band are then obtained (after averaging on several periods of the wave) by
$`{\displaystyle \frac{\mathrm{d}\rho _{\mathrm{GW}}(k,\eta )}{\mathrm{d}\mathrm{ln}(k)}}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2\kappa _{\mathrm{eff}}[\varphi ]}}\left({\displaystyle \frac{k}{a}}\right)^2P_h(k)\overline{T}^2(k,\eta ),`$ (35)
$`{\displaystyle \frac{\mathrm{d}\mathrm{\Omega }_{\mathrm{GW}}(k,\eta )}{\mathrm{d}\mathrm{ln}(k)}}`$ $`=`$ $`{\displaystyle \frac{1}{6\pi ^2}}\left({\displaystyle \frac{\kappa }{\kappa _{\mathrm{eff}}[\varphi ]}}\right)\left({\displaystyle \frac{k}{}}\right)^2P_h(k)\overline{T}^2(k,\eta ).`$ (36)
Let us stress some important points. Since we have to average on several periods, these expressions are valid only in a “shortwave limit” (see for discussion) where (i) the amplitude of the perturbation is small and (ii) the wavelength of the wave is small compared to the typical radius of the background spacetime. In our case this can be rephrased as $`k/>1`$ which implies that the expressions (35-36) are valid only for modes which are “subhorizon” today, i.e. whose wavelength is smaller than the Hubble radius today. For such modes the ergodic hypothesis is well justified. In fact, because of the averaging procedure of the transfer function, we have to restrict to modes such that $`k/_060`$ if we want to average on about ten periods. Again, we emphasize that there is an explicit dependence of $`\mathrm{\Omega }_{\mathrm{GW}}`$ and $`\rho _{\mathrm{GW}}`$ on the scalar field $`\varphi `$ because of the non-minimal coupling and our expressions reduce to the standard ones when $`\xi =0`$. We have described the gravitational waves by two stochastic variables $`\widehat{h}^{(\lambda )}`$ which can be understood as being the classical limit of a complete quantum description of the gravitational waves (see e.g. for details).
Before turning to the effects of the gravitational waves on the CMB, let us make a comment that will lead us to introduce some new notations. In the previous analysis we decomposed the metric perturbation $`h_{ij}`$ on the basis $`\stackrel{~}{Q}_{ij}^\lambda (𝐱,𝐤)ϵ_{ij}^\lambda \text{exp}(i𝐤.𝐱)`$ of TT eigenfunctions of the Laplacian, i.e. such that
$$\mathrm{\Delta }\stackrel{~}{Q}_{ij}^\lambda =k^2\stackrel{~}{Q}_{ij}^\lambda \text{with}^i\stackrel{~}{Q}_{ij}^\lambda =\delta ^{ij}\stackrel{~}{Q}_{ij}^\lambda =0.$$
(37)
Such a decomposition is indeed not unique and we could have chosen any other such basis. In the CMB literature, one often prefers to use the basis $`Q_{ij}^{(\pm 2)}(𝐱,𝐤)`$ defined by
$$Q_{ij}^{(\pm 2)}\sqrt{\frac{3}{8}}(e_1\pm ie_2)_i(e_1\pm ie_2)_j\text{e}^{i𝐤.𝐱},$$
(38)
the vectors $`e_1`$ and $`e_2`$ being defined above (26). If we decompose $`h_{ij}`$ on the latter basis as
$$h_{ij}=\underset{m=\pm 2}{}\frac{\mathrm{d}^3𝐤}{(2\pi )^3}2H^{(m)}Q_{ij}^{(m)}(𝐱,𝐤),$$
(39)
then the two decompositions are related by
$$\widehat{h}^{(\times )}=\sqrt{\frac{3}{2}}[H^{(+2)}+H^{(2)}],\widehat{h}^{(+)}=\sqrt{\frac{3}{2}}i[H^{(+2)}H^{(2)}].$$
(40)
The two polarisations $`H^{(\pm 2)}`$ are then solution of (21) which reads
$$\ddot{H}^{(m)}+2\left[\kappa _{\mathrm{eff}}\xi \dot{f}(\varphi )\right]\dot{H}^{(m)}+k^2H^{(m)}=\kappa Pa^2\pi ^{(m)}$$
(41)
where $`\pi ^{(m)}`$ is the coefficient of the development of $`\delta T_{ij}`$ as in (39) so that the transfer functions for $`H^{(m)}`$ and $`\widehat{h}^{(\lambda )}`$ are the same. If we now define the power spectrum of $`H^{(m)}(𝐤,\eta _{\mathrm{in}})`$ as
$$H^{(m_1)}(𝐤,\eta _{\mathrm{in}})H^{(m_2)}(𝐤^{},\eta _{\mathrm{in}})=(2\pi )^3k^3P_T(k)\delta (𝐤𝐤^{})\delta _{m_1,m_2}$$
(42)
one can easily check that
$$P_T(k)=\frac{1}{3}P_h(k)$$
(43)
and that if the two polarisations $`+`$ and $`\times `$ are independent then $`H^{(+2)}`$ and $`H^{(2)}`$ are also independent. With these notations, the energy density spectra are given as
$`{\displaystyle \frac{\mathrm{d}\rho _{\mathrm{GW}}}{\mathrm{d}\mathrm{ln}(k)}}`$ $`=`$ $`{\displaystyle \frac{3}{2\pi ^2\kappa }}\left({\displaystyle \frac{k}{a}}\right)^2P_T(k)\overline{T}^2(k,\eta ),`$ (44)
$`{\displaystyle \frac{\mathrm{d}\mathrm{\Omega }_{\mathrm{GW}}}{\mathrm{d}\mathrm{ln}(k)}}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}\left({\displaystyle \frac{\kappa }{\kappa _{\mathrm{eff}}}}\right)\left({\displaystyle \frac{k}{}}\right)^2P_T(k)\overline{T}^2(k,\eta ).`$ (45)
Indeed, this does not change the result but we found interesting to make the link between the notations used in the gravitational waves literature and in the CMB literature , specially because we want to present both in a unified framework and language. From now on, we use the second decomposition and its interest will be enlightened by the study of the CMB anisotropies.
### B CMB temperature and polarisation anisotropies
Gravitational waves, being metric perturbations, have an effect on the temperature and polarisation of the CMB photons. Any metric perturbation induces a fluctuation on the CMB temperature $`\mathrm{\Theta }`$ through the Sachs-Wolfe effect and any anisotropic distribution of photons scattered by electrons will become polarised and vice-versa. Since Thomson scattering generates linear polarisation, we only need to consider the Stokes parameters $`Q`$ and $`U`$ and more conveniently their two combinations $`Q\pm iU`$ which are invariant under rotation.
Following Hu and White , we decompose the tensorial part of the temperature anisotropies according to
$`\mathrm{\Theta }(\eta ,𝐱,\widehat{n})`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^3𝐤}{(2\pi )^3}\underset{\mathrm{}}{}\underset{m=\pm 2}{}\mathrm{\Theta }_{\mathrm{}}^{(m)}(k,\eta )G_{\mathrm{}}^m(𝐤,𝐱,\widehat{n})},`$ (46)
$`(Q\pm iU)(\eta ,𝐱,\widehat{n})`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^3𝐤}{(2\pi )^3}\underset{\mathrm{}}{}\underset{m=\pm 2}{}(E_{\mathrm{}}^{(m)}\pm iB_{\mathrm{}}^{(m)})(k,\eta )_{\pm 2}G_{\mathrm{}}^m(𝐤,𝐱,\widehat{n})}.`$ (47)
The coefficients $`E_{\mathrm{}}^{(m)}`$ and $`B_{\mathrm{}}^{(m)}`$ transform as $`E_{\mathrm{}}()^{\mathrm{}}E_{\mathrm{}}`$ and $`B_{\mathrm{}}()^{\mathrm{}}B_{\mathrm{}}`$ under parity and are called the “electric” and “magnetic” part of the polarisation. The functions $`G_{\mathrm{}}^m`$, $`{}_{\pm 2}{}^{}G_{\mathrm{}}^{m}`$ form three independent sets of orthonormal functions and depend both on the position $`𝐱`$ and on the direction of propagation of the photons $`\widehat{n}`$ and are defined as
$`G_{\mathrm{}}^m(𝐤,𝐱,\widehat{n})`$ $``$ $`()^{\mathrm{}}\sqrt{{\displaystyle \frac{4\pi }{2\mathrm{}+1}}}Y_{\mathrm{}}^m(\widehat{n})\mathrm{exp}(i𝐤.𝐱),`$ (48)
$`{}_{\pm 2}{}^{}G_{\mathrm{}}^{m}(𝐤,𝐱,\widehat{n})`$ $``$ $`()^{\mathrm{}}\sqrt{{\displaystyle \frac{4\pi }{2\mathrm{}+1}}}_{\pm 2}Y_{\mathrm{}}^m(\widehat{n})\mathrm{exp}(i𝐤.𝐱),`$ (49)
where the functions $`Y_{\mathrm{}}^m(\widehat{n})`$ are the standard spherical harmonics and the functions $`{}_{\pm 2}{}^{}Y_{\mathrm{}}^{m}(\widehat{n})`$ are the spin-weighted spherical harmonics . Note that the decomposition on the basis $`Q_{ij}^{(\pm 2)}`$ in the previous section is enlightened by the fact that $`Q_{ij}^{(m)}n^in^j=G_2^{(m)}`$ .
The angular correlation function of these temperature/polarisation anisotropies are observed on a 2-sphere around us and can be decomposed in Legendre polynomials $`P_{\mathrm{}}`$ as
$$\frac{\delta U}{T}(\widehat{\gamma }_1)\frac{\delta V}{T}(\widehat{\gamma }_2)=\frac{1}{4\pi }\underset{\mathrm{}}{}(2\mathrm{}+1)C_{\mathrm{}}^{UV}P_{\mathrm{}}(\widehat{\gamma }_1.\widehat{\gamma }_2),$$
(50)
where $`U`$, $`V`$ stand for $`\mathrm{\Theta }`$, $`E`$, or $`B`$. Now the brackets stand for an average on the sky, i.e. on all pairs $`(\widehat{\gamma }_1,\widehat{\gamma }_2)`$ such that $`\widehat{\gamma }_1.\widehat{\gamma }_2=\mathrm{cos}\theta _{12}`$. Using the orthonormality properties of the eigenfunctions $`G`$, equations (46-47) can be inverted to extract the angular power spectra $`C_{\mathrm{}}^{UV}`$ of the temperature and polarisation anisotropies as
$$T_0^2(2\mathrm{}+1)^2C_{\mathrm{}}^{UV}(\eta _0)=\frac{2}{\pi }\frac{\mathrm{d}k}{k}\underset{m=\pm 2}{}k^3U_{\mathrm{}}^{(m)}(\eta _0,k)V_{\mathrm{}}^{(m)}(\eta _0,k).$$
(51)
The equations of evolution of the temperature and polarisation multipoles $`\mathrm{\Theta }_{\mathrm{}}^{(m)}`$, $`E_{\mathrm{}}^{(m)}`$ and $`B_{\mathrm{}}^{(m)}`$ can be obtained by decomposing the Boltzmann equation satisfied by the photon (or neutrino) distribution function on the eigenfunctions $`G`$ to get
$`\dot{\mathrm{\Theta }}_{\mathrm{}}^{(m)}`$ $`=`$ $`k\left[{\displaystyle \frac{{}_{0}{}^{}\kappa _{\mathrm{}}^{m}}{2\mathrm{}1}}\mathrm{\Theta }_\mathrm{}1^{(m)}{\displaystyle \frac{{}_{0}{}^{}\kappa _{\mathrm{}+1}^{m}}{2\mathrm{}+3}}\mathrm{\Theta }_{\mathrm{}+1}^{(m)}\right]\dot{\tau }\mathrm{\Theta }_{\mathrm{}}^{(m)}+\delta _{\mathrm{},2}S^{(m)},`$ (52)
$`\dot{E}_{\mathrm{}}^{(m)}`$ $`=`$ $`k\left[{\displaystyle \frac{{}_{2}{}^{}\kappa _{\mathrm{}}^{m}}{2\mathrm{}1}}E_\mathrm{}1^{(m)}{\displaystyle \frac{2m}{\mathrm{}(\mathrm{}+1)}}B_{\mathrm{}}^{(m)}{\displaystyle \frac{{}_{2}{}^{}\kappa _{\mathrm{}+1}^{m}}{2\mathrm{}+3}}E_{\mathrm{}+1}^{(m)}\right]\dot{\tau }\left(E_{\mathrm{}}^{(m)}+\delta _{\mathrm{},2}\sqrt{6}P^{(m)}\right),`$ (53)
$`\dot{B}_{\mathrm{}}^{(m)}`$ $`=`$ $`k\left[{\displaystyle \frac{{}_{2}{}^{}\kappa _{\mathrm{}}^{m}}{2\mathrm{}1}}B_\mathrm{}1^{(m)}+{\displaystyle \frac{2m}{\mathrm{}(\mathrm{}+1)}}E_{\mathrm{}}^{(m)}{\displaystyle \frac{{}_{2}{}^{}\kappa _{\mathrm{}+1}^{m}}{2\mathrm{}+3}}B_{\mathrm{}+1}^{(m)}\right]\dot{\tau }B_{\mathrm{}}^{(m)},`$ (54)
where
$`S^{(m)}`$ $``$ $`\dot{\tau }P^{(m)}\dot{H}^{(m)},`$ (55)
$`P^{(m)}`$ $``$ $`{\displaystyle \frac{1}{10}}\left(\mathrm{\Theta }_2^{(m)}\sqrt{6}E_2^{(m)}\right),`$ (56)
$`{}_{s}{}^{}\kappa _{\mathrm{}}^{m}`$ $``$ $`\mathrm{}\sqrt{\left(1{\displaystyle \frac{m^2}{\mathrm{}^2}}\right)\left(1{\displaystyle \frac{s^2}{\mathrm{}^2}}\right)}.`$ (57)
The differential optical depth $`\dot{\tau }`$ vanishes for neutrinos (except in the very early universe, but the observable wavelengths in the CMB are not affected by this) and is proportional to the free electron density in the case of photons. It has to be calculated by solving the relevant kinetic recombination equations for hydrogen and helium .
The quantity $`P^{(m)}`$ represents the coupling between temperature and polarisation. Due to our choice of decomposition, only the electric part of the polarisation is affected by temperature anisotropies. However, electric and magnetic part of polarisation couple themselves as photons propagate. The Clebsch-Gordan coefficients $`{}_{s}{}^{}\kappa _{\mathrm{}}^{m}`$ arise from product properties of spherical harmonics. They are obtained in the same way as for the scalar modes .
$`{\displaystyle \frac{\mathrm{\Theta }_{\mathrm{}}^{(m)}(\eta _0,k)}{2\mathrm{}+1}}`$ $`=`$ $`{\displaystyle _0^{\eta _0}}d\eta e^\tau S^{(m)}j_{\mathrm{}}^{(m)}(k(\eta _0\eta )),`$ (58)
$`{\displaystyle \frac{E_{\mathrm{}}^{(m)}(\eta _0,k)}{2\mathrm{}+1}}`$ $`=`$ $`\sqrt{6}{\displaystyle _0^{\eta _0}}d\eta \dot{\tau }e^\tau P^{(m)}ϵ_{\mathrm{}}^{(m)}(k(\eta _0\eta )),`$ (59)
$`{\displaystyle \frac{B_{\mathrm{}}^{(m)}(\eta _0,k)}{2\mathrm{}+1}}`$ $`=`$ $`\sqrt{6}{\displaystyle _0^{\eta _0}}d\eta \dot{\tau }e^\tau P^{(m)}\beta _{\mathrm{}}^{(m)}(k(\eta _0\eta )).`$ (60)
The functions $`j_{\mathrm{}}^{(m)}`$, $`ϵ_{\mathrm{}}^{(m)}`$ and $`\beta _{\mathrm{}}^{(m)}`$ are defined in terms of the spherical Bessel functions, $`j_{\mathrm{}}(x)`$ as
$`j_{\mathrm{}}^{(\pm 2)}(x)`$ $``$ $`\sqrt{{\displaystyle \frac{3}{8}}{\displaystyle \frac{(\mathrm{}+2)!}{(\mathrm{}2)!}}}{\displaystyle \frac{j_{\mathrm{}}(x)}{x^2}},`$ (61)
$`ϵ_{\mathrm{}}^{(\pm 2)}(x)`$ $``$ $`{\displaystyle \frac{1}{4}}\left[j_{\mathrm{}}(x)+j_{\mathrm{}}^{\prime \prime }(x)+2{\displaystyle \frac{j_{\mathrm{}}(x)}{x^2}}+4{\displaystyle \frac{j_{\mathrm{}}^{}(x)}{x}}\right],`$ (62)
$`\beta _{\mathrm{}}^{(\pm 2)}(x)`$ $``$ $`\pm {\displaystyle \frac{1}{2}}\left[j_{\mathrm{}}^{}(x)+2{\displaystyle \frac{j_{\mathrm{}}(x)}{x}}\right]`$ (63)
where a prime denotes a derivative with respect to the argument $`x`$.
Beside the small contribution due to the polarisation in Eq. (58), the temperature fluctuation in the direction $`\widehat{\gamma }`$ reduce to the well known result by Sachs and Wolfe
$$\frac{\delta T}{T}(\widehat{\gamma })=\frac{1}{2}_{\eta _{\mathrm{LSS}}}^{\eta _0}\dot{h}_{ij}\gamma ^i\gamma ^jd\eta ,$$
(64)
where the subscript $`\mathrm{LSS}`$ stands for last scattering surface. The “visibility function” $`\dot{\tau }e^\tau `$ appearing in equations (59-60) takes a non zero value only at the time of decoupling so that, contrarily to temperature anisotropies which are constantly generated by gravitational interactions with photons, polarisation is generated only at the last scattering surface.
## IV Damping of the gravitational waves
We first study the effect of the damping of the gravitational waves due to the anisotropic stress of the photons. Such a damping of the amplitude of gravitational waves in various viscous cosmic media has been already discussed in ; we give a description of this damping in the formalism we use here in order to quantify precisely its effect on CMB anisotropies. On subhorizon scales larger than the diffusion length $`\lambda _\mathrm{D}\dot{\tau }^1`$ of the photons, i.e. such that
$$_{\mathrm{eff}}\tau _{}^1k\dot{\tau },$$
(65)
the set of equations (52-57) implies that
$$\mathrm{\Theta }_2^{(\pm 2)}=\frac{4}{3}\frac{\dot{H}}{\dot{\tau }}\text{and}E_2^{(\pm 2)}=\frac{\sqrt{6}}{4}\mathrm{\Theta }_2^{(\pm 2)}.$$
(66)
Since $`\pi ^{(\pm 2)}`$ is proportional to $`\mathrm{\Theta }_2^{(\pm 2)}`$, we can insert the former expressions in the gravitational waves evolution equation (41) to get the back reaction of the anisotropic stress
$$\ddot{H}+2_{\mathrm{eff}}\dot{H}+k^2H=\frac{32}{15}\kappa Pa^2\frac{\dot{H}}{\dot{\tau }}.$$
(67)
Setting
$$H\overline{H}h,$$
(68)
with $`\overline{H}`$ solution of the homogeneous equation of evolution (i.e. with $`\pi ^{(m)}=0`$), we get
$$\overline{H}\ddot{h}+\left(2\dot{\overline{H}}+2_{\mathrm{eff}}\overline{H}+\frac{32}{15}\kappa Pa^2\frac{\dot{\overline{H}}}{\dot{\tau }}\right)\dot{h}+\frac{32}{15}\kappa Pa^2\frac{\dot{\overline{H}}}{\dot{\tau }}h=0.$$
(69)
Now, (i) since $`\dot{\overline{H}}k\overline{H}`$, we deduce that $`\dot{\overline{H}}_{\mathrm{eff}}\overline{H}`$, (ii) since $`\kappa Pa^2_{\mathrm{eff}}`$, it follows that $`(32/15)\kappa Pa^2\dot{\overline{H}}/\dot{\tau }_{\mathrm{eff}}^2/\dot{\tau }k\overline{H}\dot{\overline{H}}`$ and (iii) since $`\ddot{h}\dot{h}/\tau _\mathrm{D}`$, $`\overline{H}\ddot{h}\dot{\overline{H}}\dot{h}`$ and in conclusion in the limit (65) the equation of evolution of the gravitational waves in presence of the anisotropic stress (69) reduces to
$$\dot{h}=\frac{16}{15}\kappa Pa^2\dot{\tau }h.$$
(70)
We deduce that a mode $`k`$ is damped from the time it enters the Hubble radius, i.e. $`\eta k^1`$ since it happens during the radiation era, to roughly the time when the anisotropic stress becomes negligible, i.e. approximately at the time of last scattering $`\eta _{\mathrm{LSS}}`$. It follows that
$$h(k,\eta _{\mathrm{LSS}})\text{exp}\left(\frac{16}{15}_{1/k}^{\eta _{\mathrm{LSS}}}\frac{\kappa Pa^2}{\dot{\tau }}𝑑\eta \right)h(k,1/k).$$
(71)
This damping of the gravitational waves by the anisotropic stress of the photon fluid is analogous to the damping of the scalar modes (density fluctuations) known as the Silk damping a description of which, in the formalism used here, can be found in . Note however that the origin of the damping is different.
This effect is small but, apart from , was not much emphasized in the literature before. Assuming that the universe is completely ionized until the last scattering surface, the integral of Eq. (71) is of order
$$\frac{1}{3}\left(1\frac{Y_{\mathrm{He}}}{2}\right)\frac{\mathrm{\Omega }_\gamma ^0}{\mathrm{\Omega }_b^0}\frac{m_\mathrm{p}\kappa }{\sigma _{\mathrm{Th}}}a_0(\eta _{\mathrm{LSS}}1/k)10^3(11/k\eta _{\mathrm{LSS}}),$$
(72)
where $`m_\mathrm{p}`$ is the proton mass, $`\sigma _{\mathrm{Th}}`$ is Thomson scattering cross section. The real damping factor is greater than the estimate (72) because the universe becomes neutral at the last scattering surface (so that the term $`\dot{\tau }`$ is smaller). In Fig. 1 \[left\], we plot this damping factor for the modes that entered into the Hubble radius long before the last scattering surface (i.e. such that $`k\eta _{\mathrm{LSS}}^1`$). As a consequence, the comparison between the damped case to the undamped case, shozn on Fig. 1 \[right\] does not show significant differences. The amplitude of the high-$`\mathrm{}`$ tail of the CMB anisotropy spectrum is lowered by roughly $`10\%`$ when one includes this effect. The same occurs of course for polarisation. We emphasize that this result does not depend on any particular model, and is not included in the most recent (3.2) version of CMBFAST.
## V Specification of the model
### A Model parameters
At this stage, the model we are discussing depends on five parameters:
1. $`f(\varphi )`$ which is an arbitrary function of the scalar field $`\varphi `$,
2. $`\xi `$ the coupling of the scalar field with the background spacetime,
3. $`\alpha `$ the slope of the potential (14) or (16),
4. $`\mathrm{\Omega }_\varphi ^0`$ the energy density in the scalar field today,
5. $`P_T(k)`$ the spectrum of the gravitational waves.
Indeed there exist some constraints on these functions and parameters and we make the following assumptions and choices:
1. We assume that $`f(\varphi )=\varphi ^2/2`$; this is the only choice for which the coupling constant $`\xi `$ is dimensionless. Moreover such a choice can be seen as the lowest term in an expansion of $`f`$ in powers of $`\varphi `$. As shown in there exists tracking solutions for the field $`\varphi `$ evolving in the potential (14) with such a coupling.
2. If the scalar field $`\varphi `$ is coupled to the spacetime metric, this coupling must be weak enough so that it does not generate a significant time variation of the constants of nature . Taking into account the bound on the variation $`|\dot{G}/G|`$ of the Newton constant and on the variation $`|\dot{\alpha }/\alpha |`$ of the fine structure constant , it was shown that
$$10^2\xi 10^210^1.$$
(73)
This bound is however sensitive to the shape of the potential. On the other hand the experimental constraints (from the Shapiro effect and the light deflection in the Solar system) on the post-Newtonian parameters imply
$$\left|\xi \right|\frac{3.9\times 10^2}{\sqrt{\alpha (\alpha +2)}}.$$
(74)
However, in this class of models one does not try to have a theory converging towards general relativity at late time and the coupling $`\xi `$ is constant which is the main reasons of the severe bounds on its value. This can be improved by generalising this kind of models by considering them in the framework of scalar-tensor theories .
3. In most models $`\alpha `$ is not constrained theoretically. If the matter content of the universe today is dominated by the matter-radiation fluid then the fact that the observations favour $`1<\omega _\varphi <0.6`$ gives a bound on $`\alpha `$, which is indeeed not the case anymore if the scalar field starts to dominate. In Fig. 2 \[left\], we compare this analytic estimate and the numerical determination of the energy scale $`\mathrm{\Lambda }`$ as a function of the slope $`\alpha `$. We see that if $`\alpha >4`$ then $`\mathrm{\Lambda }`$ is at least larger than 1 TeV (when $`\mathrm{\Omega }_\varphi ^0=0.7`$).
4. The density parameter $`\mathrm{\Omega }_\varphi ^0`$ is not severely constrained theoretically, but observations seem to indicate $`\mathrm{\Omega }_\varphi ^00.7`$. One has to check that if the scalar field was dominating the matter content of the universe at some early stage then it has to be subdominant at the time of nucleosynthesis (see e.g. ). The choice of $`\mathrm{\Omega }_\varphi ^0`$ fixes the value of the energy scale $`\mathrm{\Lambda }`$ in (14) or (16); this is the coincidence problem. On Fig. 2 \[right\], we depict the variation of the energy scale $`\mathrm{\Lambda }`$ with $`\mathrm{\Omega }_\varphi ^0`$ and $`\alpha `$. It is not very sensitive to $`\mathrm{\Omega }_\varphi ^0`$ as long as $`0.1<\mathrm{\Omega }_\varphi ^0<0.9`$. In fact, when the quintessence field starts to dominate the matter content and if we have reached the attractor then $`\mathrm{d}^2V/\mathrm{d}\varphi ^2H^2`$ (see ), and $`H^2V/M_{\mathrm{Pl}}^2`$ so that we can estimate that the variation of $`\mathrm{\Lambda }`$ with $`\alpha `$ follows
$$\mathrm{\Lambda }=\left(\rho _{\mathrm{crit}}M_{\mathrm{Pl}}^\alpha \right)^{\frac{1}{4+\alpha }}.$$
(75)
We conclude that
$$\frac{\delta \mathrm{\Lambda }}{\mathrm{\Lambda }}\frac{1}{4+\alpha }\frac{\delta \mathrm{\Omega }_\varphi ^0}{\mathrm{\Omega }_\varphi ^0}$$
and thus a prescision of 10% on $`\mathrm{\Omega }_\varphi ^0`$ requires to tune $`\mathrm{\Lambda }`$ at a 1% level if e.g. $`\alpha =6`$, which is a less drastic tuning than the usual cosmological constant fine tuning problem.
5. $`P_T`$ has to be determined by a specific model, such as e.g. inflation, and we parametrise it as
$$P_T(k)A_Tk^{n_T}$$
(76)
where $`A_T`$ is a constant dans $`n_T`$ is the tensor mode spectral index. $`A_T`$ is obtained by normalising the CMB temperature anisotropies to the COBE data at $`\mathrm{}=10`$ for which
$$T_0\sqrt{\frac{\mathrm{}(\mathrm{}+1)}{2\pi }C_{\mathrm{}}^{\mathrm{\Theta }\mathrm{\Theta }}}30\mu \text{K}.$$
(77)
Since some measurements tend to show that there is a peak at the degree scale , we conclude that a significant part of the anisotropies may be generated by the scalar modes. In the “standard” slow-roll inflation picture, this is compatible with an almost scale-invariant spectrum with a low tensor contribution, in which case the COBE results would put only an upper limit on the amplitude of the gravitational waves spectrum. Nevertheless, we point out that it is also possible that most of the large scale anisotropies can be generated by gravitational waves. This assumes a strong deviation from scale invariance ($`n_S=1.69`$ and $`n_T=0.0`$), but is in good agreement with observational data .
### B Initial conditions and behaviour of the background spacetime
Concerning the initial conditions for the scalar field $`\varphi `$, we will consider the two extreme cases:
* (IC1) where we assume that the scalar field is at equipartition with the matter (i.e. mainly with the radiation) deep in the radiation era,
* (IC2) where it dominates the matter content of the universe at a very early stage.
Situation (IC1) implies that at the end of reheating,
$$\rho _\varphi 10^4\rho _\gamma ,$$
(78)
where the factor $`10^4`$ is roughly the inverse of the number of degrees of freedom at that time. Since the quintessence field is already subdominant at this epoch, one does not need to care about its effect on nucleosynthesis since it remains subdominant until recently. In the second situation (IC2), the field starts by dominating and inflation ends by a kinetic phase rapidly than $`\rho _\gamma `$ and will thus become subdominant. One has to check that this happens before nucleosynthesis . A realisation of such initial conditions can be obtained in quintessential inflation .
In Fig. 3, we depict the evolution of the energy density of the quintessence field, matter and radiation for the initial conditions (IC1) \[left\] and (IC2) \[right\]. We see that for a very large range of initial conditions (roughly for $`10^{47}\text{GeV}^4\rho _\varphi 10^{113}\text{GeV}^4`$ at a redshift of $`z10^{30}`$) we end up with a quintessence field which starts to dominate today. This explains briefly how the fine tuning problem is solved . We can also check that with these values the scalar field does not dominate the matter content of the universe at nucleosynthesis, i.e. at a redshift of order $`z10^{10}`$.
An interesting point concerns the evolution of the scalar field equation of state in the case (IC2) when $`\xi =0`$. The field rolls down very fastly so that we are first in a regime where
$$\rho P\frac{1}{2}\frac{\dot{\varphi }^2}{a^2}$$
(79)
from which we conclude that its equation of state is $`\omega _\varphi 1`$ (see Fig. 4). But, due to the exponential behaviour contribution of the potential, the field is stopped when $`\varphi M_{\mathrm{Pl}}`$ and then rolls back to smaller values (see Fig. 5) so that the field undergoes a series of damped oscillations (because of the friction term coming from the expansion in the Klein-Gordon equation). This implies that there exist times such that $`\dot{\varphi }0`$ and thus small period around them where the equation of state varies rapidly to $`\omega _\varphi 1`$ (see Fig. 4). This sudden change in the equation of state of $`\varphi `$ happen while it is dominating the matter content of the universe (see Fig. 3) so that it implies variations in the evolution of the scale factor of the universe which, in principle, should let a signature in the gravitational waves energy spectrum. Indeed, this does not happen in standard quintessence and is a specific feature of the SUGRA-quintessence.
When $`\xi 0`$, there are no significant modifications to the background dynamics as long as the field has not reached the Planck mass \[because $`2\xi \kappa f(\varphi )`$ is small compared to unity, see Eq. 5\]. Then, the main difference appears at late times when the field starts to dominate and comes from the fact that the bound $`1<\omega _\varphi <1`$ no longer applies, and one can get lower values of $`\omega _\varphi `$. Equivalently, the equation of state parameter $`\omega (P_{\mathrm{fluid}}+P_\varphi )/(\rho _{\mathrm{fluid}}+\rho _\varphi )`$ for the whole background fluids can reach values smaller than $`1`$ (see Fig. 6 where we plot the variation of $`\omega _\varphi `$ as a function of redshift). As pointed out by Caldwell , such a matter fits the current observational data. Different candidates such as a decaying dark matter component and a kinetic quintessence field were proposed. Here, we show that any non-minimally coupled scalar field may be a good candidate for a component of matter with $`\omega <1`$. The constraints (73) on $`\xi `$ implies that for our class of models
$$3\omega _\varphi <0$$
if the scalar field dominates. We emphasize that $`\omega _\varphi `$ is not uniquely defined according to the way one splits $`T_{(\varphi )}^{\mu \nu }`$ in (6). In Fig. 6, we used the Friedmann equations (7-8) to extract $`\omega `$ from
$$\frac{\dot{}}{^2}1=\frac{3}{2}(1+\omega )\mathrm{\Omega }$$
and then $`\omega _\varphi `$ from
$$\omega \mathrm{\Omega }=\underset{i}{}\omega _i\mathrm{\Omega }_i$$
where $`i`$ runs on all the matter species. This corresponds to the value of $`\omega `$ as it may be reconstructed from observational data such as e.g. the supernovae type Ia.
## VI Qualitative discussion
### A Gravitational waves spectrum
Equation (41) describes the evolution of a damped oscillator. Injecting the ansatz
$$H^{(m)}A^{(m)}\mathrm{exp}(ik\eta )$$
(80)
in (41) and performing a WKB approximation leads to the equation
$`\dot{A}^{(m)}+_{\mathrm{eff}}A^{(m)}`$ $`=`$ $`0`$ (81)
for the evolution of the amplitude $`A^{(m)}`$ where $`_{\mathrm{eff}}\dot{\stackrel{~}{a}}/\stackrel{~}{a}`$, with
$$\stackrel{~}{a}a\sqrt{12\kappa \xi f(\varphi )}.$$
(82)
This WKB approximation holds only for “sub-horizon” modes. Before a mode has a wavelength smaller than the Hubble radius, its amplitude evolves according to
$$\ddot{A}^{(m)}+2_{\mathrm{eff}}\dot{A}^{(m)}=0,$$
(83)
the solutions of which are a constant mode and a decaying mode. Neglecting the decaying mode, we see that the wave is “frozen” as long as its wavelength is larger than the Hubble radius, and that it undergoes damped oscillations once its wavelength is shorter than the Hubble radius. The damping of a mode of wavenumber $`k`$ between the time it enters the Hubble radius and today is then proportional to
$$\frac{\stackrel{~}{a}_k}{\stackrel{~}{a}_0},$$
(84)
where $`\stackrel{~}{a}_k`$ is the scale factor evaluated at the time $`\eta _k`$ when the mode $`k`$ enters the Hubble radius (i.e. when $`=k`$) and $`\stackrel{~}{a}_0`$ is scale factor today. Injecting this behaviour in (36), we obtain that the energy density spectrum of gravitational waves scales as
$$\frac{\mathrm{d}\mathrm{\Omega }_{\mathrm{GW}}}{\mathrm{d}\mathrm{ln}(k)}k^2\stackrel{~}{a}_k^2P_T(k).$$
(85)
First let us assume that $`\xi =0`$. For wavelengths corresponding to modes that have entered the Hubble radius in the matter dominated era (for which $`a\eta ^2`$ and thus $`\eta _kk^1`$), one can easily sort out that
$$\stackrel{~}{a}_kk^2$$
(86)
and the gravitational waves spectrum behaves as
$$\frac{\mathrm{d}\mathrm{\Omega }_{\mathrm{GW}}}{\mathrm{d}\mathrm{ln}(k)}k^2P_T(k).$$
(87)
Equivalently, for wavelengths corresponding to modes entering the Hubble radius in the radiation dominated era (for which $`a\eta `$) one can show that the gravitational waves energy spectrum behaves as
$$\frac{\mathrm{d}\mathrm{\Omega }_{\mathrm{GW}}}{\mathrm{d}\mathrm{ln}(k)}k^0P_T(k).$$
(88)
To finish, if it happens that there exist wavelengths corresponding to modes that have entered the Hubble radius while the scalar field was dominating (for which $`a\sqrt{\eta }`$ since $`\rho _\varphi 1/a^6`$) one obtains that
$$\frac{\mathrm{d}\mathrm{\Omega }_{\mathrm{GW}}}{\mathrm{d}\mathrm{ln}(k)}k^1P_T(k).$$
(89)
In conclusion, we have found three behaviours for the gravitational waves spectrum according to the wavelength. In Fig. 9, we give an example of such a spectrum in a case where one has a scalar field dominating at early stage \[initial condition (IC2)\]. These results hold also when $`\xi 0`$ but the slopes of the spectrum are slightly changed since the time behaviour of $`a`$ has to be replaced by the one of $`\stackrel{~}{a}`$.
### B CMB anisotropies
For scales smaller than the Hubble radius at decoupling, one can follow the same lines to predict the tensor part of the CMB temperature anisotropies. The main difference is that the expression for $`\mathrm{}(\mathrm{}+1)C_{\mathrm{}}`$ does not involve any factor $`k^2`$ as in Eqns. (45), the reason being that Eq. (58) can be integrated by parts to drop the time derivative of $`H^{(m)}`$, which shows that anisotropies are mostly generated on the last scattering surface with an amplitude of $`|H^{(m)}|^2`$. Therefore, the spectrum behaves as
$$\mathrm{}(\mathrm{}+1)C_{\mathrm{}}\mathrm{}^{n_T4},\mathrm{}^{n_T2},\mathrm{}^{n_T1},$$
(90)
for modes which have entered the Hubble radius is the matter dominated, radiation dominated and kinetic scalar field dominated eras respectively. With standard cosmological parameters, the radiation to matter transition occurs soon before the decoupling, and the scalar field dominates only at very early times. As a consequence, one sees almost only the regime $`\mathrm{}(\mathrm{}+1)C_{\mathrm{}}\mathrm{}^{n_T2}`$. For modes which enter into the Hubble radius after the last scattering surface, one can show that the produced spectrum scales as
$$\mathrm{}(\mathrm{}+1)C_{\mathrm{}}\mathrm{}^{n_T}.$$
(91)
Note that this expression is indeed an approximation and that it is not easy to calculate an accurate analytical solution . These results are illustrated on Fig. 7. As already stressed, the result of Eq. 91 applies at large angular scales which have not entered into the Hubble radius at recombination. For standard cosmologies, this occurs for multipoles smaller than $`\mathrm{}100`$ (in addition, there are also some corrections to this rough estimate which occur at the very smallest multipoles and slightly boost the spectrum, as can also be seen on Fig. 7). Then, at higher multipoles the result of Eq. 90 is valid. The matter dominated regime before recombination is rather short, and occurs only between $`\mathrm{}100`$ and $`\mathrm{}200`$ (less than one oscillation in the spectrum). For $`\mathrm{}200`$, one sees the regime $`\mathrm{}(\mathrm{}+1)C_{\mathrm{}}\mathrm{}^{n_T2}`$ (see also Fig. 1 of ).
### C Results of the $`\mathrm{\Lambda }`$CDM model
Before turning to a more general numerical study of the class of models we consider in this article, we recall in Fig. 7 the general results for the temperature and polarisation angular power spectra and the gravitational waves density spectrum for a $`\mathrm{\Lambda }`$CDM model. This spectrum has two branches: a soft branch at lower frequencies (corresponding to the matter dominated era) and a hard branch at higher frequencies (corresponding to modes that entered the horizon in the radiation era). Following , we set the cut-off on this spectrum to the last mode that has been inflated out of the Hubble radius.
## VII Numerical results
### A Field $`\varphi `$ initially at equipartition
Since the scalar field only starts to dominate at very recent time, we expect no effect on the gravitational waves energy spectrum (since at earlier time the universe is always radiation dominated). However, the change in today’s universe equation of state yield a specific signature in the angular diameter-distance relation. Hence, one expects to see the quintessence field behaviour in the positions of the peaks in the CMB anisotropy spectra.
The temperature anisotropies plots of Fig. 8 are therefore identical at high multipoles except for their overall position which are different. At low redshift, the scalar field dominates and the dynamics of the expansion depends explicitely of the value of the coupling $`\xi `$, which cause some slight differences in the CMB anisotropies at the very first multipoles ($`\mathrm{}5`$). We have also seen that the polarisation is generated by gravity and therefore different gravitational constants lead to different normalisation between the polarisation and the temperature spectra. Since we normalise the “bare” Einstein constant $`\kappa `$ so that the effective Einstein constant corresponds to what we measure (in e.g. a Cavendish experiment), models with a different $`\xi `$ have different $`\kappa `$. At decoupling, the scalar field does not dominate and therefore $`\kappa ^{\mathrm{LSS}}=\kappa _{\mathrm{eff}}^{\mathrm{LSS}}`$. This induces different amplitudes for the polarisation anisotropy spectra. For the lowest values of $`\xi `$ there is a factor $`24`$ in amplitude as compared with the $`\xi =0`$ case, which roughly corresponds to the square of the variation of $`\kappa _{\mathrm{eff}}`$ (and, hence $`G`$) between the last scattering surface and now. Note that the effect of $`\xi `$ depends on its sign. This is the reason why the constraint derived by Chiba are stronger for negative values of $`\xi `$. The same can be seen in Fig. 6.
We conclude that the temperature anisotropies and polarisation give mainly information on the spectral index $`n_T`$, the energy density of the scalar field today $`\mathrm{\Omega }_\varphi ^0`$ and its coupling $`\xi `$.
### B Field $`\varphi `$ dominates at early stage
We now turn to the more unusual case where the scalar field dominates at the end of inflation and where the universe undergoes a kinetic phase before the radiation era as in e.g. quintessential inflation . The wavelengths corresponding to the observable CMB multipoles ($`\mathrm{}2000`$) are much larger than the Hubble radius at nucleosynthesis, epoch at which we have to be radiation dominated. As a consequence, we expect no signature from this early phase on the CMB anisotropies and polarisation.
As first pointed out in , if the scalar field dominates at early stage, there is an excess of gravitational waves at high frequency \[see equation (89)\]. On Fig. 9, we present such a spectrum and we will discuss the implication of this excess later.
An interesting effect concerns the difference between the spectra obtained from an inverse power law potential and a SUGRA-like potential. As shown on Fig. 9 \[right\], the amplitude of the spectrum at high frequency in roughly 30% higher for inverse power law potentials. The relative decrease in amplitude at these frequencies for SUGRA-like potentials depends on the dynamics of the scalar field in the bounce (see Figs. 4 and 5) during which the equation of state varies from $`+1`$ to $`1`$ and to $`+1`$ again. Thus, during this time, the modes that had just entered into the Hubble radius (and thus which had just started to undergo damped oscillations) went out of it (during the $`\omega <0`$ epoch) and their amplitude was frozen before re-entering the Hubble radius again. Hence, the modes of larger wavelengths are less damped which explains this decrease in amplitude. Now, if the slope of the potential is less steep, the bounce lasts longer (note that we always reach $`\omega =1`$ at the point where $`\dot{\varphi }=0`$) and thus the damping is stronger. This signature, even if not detectable by coming experiment is nevertheless a clear feature of supergravity.
To finish, let us discuss the total energy density of gravitational waves $`\mathrm{\Omega }_{\mathrm{GW}}^0`$ today. As pointed out in , it has also to be negligible at nucleosynthesis; this constraint is more drastic than the only requirement that $`\mathrm{\Omega }_\varphi ^0`$ be negligible at that time. Let us emphasize that the constraint on $`\mathrm{\Omega }_\varphi ^0`$ cannot be avoided (since it involves background dynamics) whereas the one on $`\mathrm{\Omega }_{\mathrm{GW}}^0`$ depends on $`A_T`$ and $`n_T`$ and thus leads to a combined constraint on the initial conditions of the scalar field and on the initial power spectrum of the gravitational waves. Besides $`A_T`$ and $`n_T`$, $`\mathrm{\Omega }_{\mathrm{GW}}^0`$ mainly depends on the initial values of $`\rho _\varphi `$ and $`\rho _{\mathrm{rad}}`$ which can be parametrised by the reheating temperature $`T_R`$ (related roughly to $`\rho _{\mathrm{rad}}`$ at that time) and the redshift $`z_{}`$ of equality between the kinetic scalar field era and the radiation era (related roughly to $`\rho _\varphi /\rho _{\mathrm{rad}}`$ at the end of reheating). $`\mathrm{\Omega }_{\mathrm{GW}}^0`$ can be estimated by the surface of the spectrum below the part with a positive slope (i.e. the high frequency part; see Fig. 9) and thus of order
$$\mathrm{\Omega }_{\mathrm{GW}}^0\frac{k_R}{k_{}}\frac{\mathrm{d}\mathrm{\Omega }_{\mathrm{GW}}}{\mathrm{d}\mathrm{ln}k}|_k_{}$$
(92)
where $`k_R`$ and $`k_{}`$ are respectively the modes entering the Hubble radius at the reheating and at $`z_{}`$. Thus the “bump” at short wavelength cannot be too high. Moreover, the energy density at the end of reheating cannot be higher that Planck scale, so that it fixes a limit on the shortest mode in which gravitational waves are produced. On Fig. 10, we first plot \[left\] the variation of the gravitational wave spectrum with the parameters $`(T_R,z_{})`$ and we then give \[right\] the “safe” zone of parameters for nucleosynthesis \[for $`n_T=0`$\] and defined by $`\mathrm{\Omega }_{\mathrm{GW}}^010^6`$. Let us briefly explain how this bounds are obtained.
1. We first rephrase in terms of $`T_R`$ the fact that the field is dominating at the end of the inflation phase, i.e.
$$z_{}<z_Rz_{}<\alpha _1T_R,$$
(93)
where $`\alpha _1`$ is some numerical coefficient. This corresponds to the solid line on Fig. 10.
2. We then impose that the scalar field is subdominant at nucleosynthesis, i.e. that
$$z_{}>10^{10}.$$
(94)
This corresponds to the horizontal dash-dot line on Fig. 10.
3. At the end of the inflation phase, we want the energy density to be smaller that the Planck energy density. If the scalar field is dominating if gives
$$\rho _\varphi ^0\frac{(1+z_R)^6}{(1+z_{})^2}(1+z_{\mathrm{eq}})<\rho _{\mathrm{Pl}}z_{}>\alpha _2T_R^3$$
(95)
where $`\alpha _2`$ is another numerical coefficient. This corresponds to the dot line on Fig. 10. Note that since we are in a field dominated era $`H^2a^2`$ and thus on this “Planck limit” we have $`k_Rz_R^2/z_{}`$ and $`k_{}z_{}`$ (and thus $`k_R1/z_R`$ and $`k_{}z_R^3`$) from wich we conclude that the maximum of the power spectrum is roughly located on a curve $`(1/z_R,1/z_R^4)`$ \[see Fig. 10\].
4. To finish, we want that the gravitational waves energy density does not alter nucleosynthesis, i.e. that
$$\mathrm{\Omega }_{\mathrm{GW}}^010^6T_R<\alpha _3z_{}$$
(96)
where $`\alpha _3`$ is a third numerical coefficient. This corresponds to the dot-dash line on Fig. 10.
For all the points $`(T_R,z_{})`$ above the dotted and dot-dashed lines of Fig. 10 \[right\], there is no excess of gravitational waves. The solid line separates the two sets of initial conditions zwe have considered. We must emphasize that this result was obtained for $`n_T=0`$ and that the spectrum can be tilted, which modifies the bounds on the parameter set $`(T_R,z_{})`$ (more precisely, taking smaller $`A_T`$ or negative $`n_T`$ lowers the diagonal dot-dashed line). Such constraints may be important for instance while considering models where a scalar field dominates at baryogenesis .
In the case of a “blue” initial power spectrum (i.e. with $`n_T>0`$ in our notations, or $`\beta >2`$ in the notations of ), as advocated for example in , the flat branch of Fig. 9 (corresponding to the “semi-hard” branch of ) is already tilted, giving as stronger constraint on our model. For instance, if $`n_T=0.2`$, the amplitude at $`\nu =10\mathrm{GHz}`$ is boosted by a factor $`3\times 10^5`$. As a consequence, the quantity of gravitational waves at high frequency cannot be boosted as much as in the case of a scale invariant spectrum, and the allowed range of parameters for our model (see Fig. 10) is narrowed.
## VIII Conclusion
In this article, we have studied some properties of quintessence models with a non-minimally coupled scalar field among which the spectrum of gravitational waves.
We have shown that such a quintessence field can behave as a fluid with $`\omega <1`$ and our models lead to $`3\omega 0`$ when the field dominates. We related the energy scale $`\mathrm{\Lambda }`$ of the potential to its slope $`\alpha `$ and to the scalar field energy density today $`\mathrm{\Omega }_\varphi ^0`$. In particular, we showed that $`\mathrm{\Lambda }`$ is almost independent of $`\mathrm{\Omega }_\varphi ^0`$. The coincidence problem, i.e. the fact that $`\mathrm{\Omega }_\varphi ^01`$ implies a tuning of $`\mathrm{\Lambda }`$ (roughly the precision on $`\mathrm{\Lambda }`$ has to one order of magnitude higher than the one on $`\mathrm{\Omega }_\varphi ^0`$) which is however less severe than the fine tuning needed for a cosmological constant. This being fixed, the tracking mechanism allows to span a very wide range of initial conditions for the scalar field and there is no fine tuning in that respect.
We then showed that the combined study of the gravitational waves energy spectrum and of their imprint on the CMB radiation temperature and polarisation enables to extract many complementary informations on the models:
* the CMB mainly gives results on $`\xi `$, $`\mathrm{\Omega }_\varphi ^0`$ and $`n_T`$,
* the energy spectrum gives results on the initial conditions of the scalar field.
As pointed out in , there is an excess of gravitational waves today if inflation ends by a kinetic phase. In that case, one has to check that both $`\mathrm{\Omega }_\varphi `$ and $`\mathrm{\Omega }_{\mathrm{GW}}`$ are negligible at the time of nucleosynthesis and we relate the amount of gravitational waves today to the reheating temperature and the time of equality between the kinetic scalar era and the radiation era.
We also pointed out that gravitational waves are damped by the anisotropic stress of radiation, which implies that the CMB anisotropy and polarisation spectra are lowered roughly by 10% for high multipoles. It was also shown that the amplitude of the gravitational waves spectrum for inverse power law potentials is $`30`$% higher than for SUGRA-like potentials at high frequency. Indeed this is probably not detectable by coming experiments but it could ultimately lead to a signature of supergravity.
###### Acknowledgements.
It is a pleasure to thank Pierre Binétruy, Nathalie Deruelle, Thibault Damour, David Langlois, Patrick Peter and Filipo Vernizzi for fruitful discussions.
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# 1 Introduction
## 1 Introduction
Previous measurements performed at GSI in Darmstadt by the EPOS and the Orange groups have shown the appearance of narrow (FWHM $`3080`$ keV), and unexpected, $`e^+`$–lines at energies in the range 250–400 keV in the positron spectra, obtained from heavy-ion collisions near the Coulomb barrier . The lines were superimposed on continuous spectral distributions from quasi-atomic positron emission and from nuclear positron background, which has been determined from measured low resolution $`\gamma `$-spectra and theoretical IPC conversion coefficients . No viable explanation could be found to account for these results. Initially, the measurements were performed to investigate the decay of the neutral vacuum of QED predicted by theory , a phenomenon which can lead to the emission of monoenergetic positrons in supercritical collisions . But it was clear that the width of the reported e<sup>+</sup>–lines was too narrow as to be attributed to this effect and, particularly, the occurrence of similar e<sup>+</sup>–lines in the so-called subcritical collision systems have excluded an interpretation of the lines in the frame work of spontaneous positron emission .
In follow-up extended e<sup>+</sup>e<sup>-</sup>-coincidence experiments electrons were measured in coincidence with positrons, and the measured e<sup>+</sup>e<sup>-</sup>–sum-energy spectra revealed even narrower lines . First attempts were made to interpret these lines in the context of the e<sup>+</sup>e<sup>-</sup>–decay of a previously unknown neutral particle with a mass around 1.8 MeV/c<sup>2</sup> . But soon this hypothesis was definitively ruled out by subsequent Bhabha-scattering experiments .
Only in recent experimental and theoretical studies, extended investigations of Internal Pair Conversion (IPC) as potential source of the observed e<sup>+</sup>e<sup>-</sup>–sum-energy lines in very heavy collision systems at the Coulomb barrier have been carried out. To be more specific, the IPC scenario has been systematically addressed in our investigations by studying the collision systems <sup>238</sup>U + <sup>206</sup>Pb and <sup>238</sup>U + <sup>181</sup>Ta . It could be shown that electromagnetic transitions in one of the colliding nuclei which de-excites by emission of IPC e<sup>+</sup>e<sup>-</sup>–pairs can lead to in principle observable narrow lines in the corresponding Doppler-shift corrected e<sup>+</sup>e<sup>-</sup>–sum-energy spectra and in back-to-back e<sup>+</sup>e<sup>-</sup>–coincidence spectra, as observed with the double-Orange spectrometer . The cross sections of the observed $`\gamma `$–lines are typically of the order of some mb up to several 10 mb leading to weak e<sup>+</sup>e<sup>-</sup>–sum-energy lines with cross sections of 0.1 $`\mu `$b to several $`\mu `$b. The experimental sensitivity for the detection of e<sup>+</sup>e<sup>-</sup>–lines was limited to cross sections of this order. Based on this new experience, we could show moreover that the general features of most of the previously reported weak e<sup>+</sup>e<sup>-</sup>–sum-energy lines resemble conspicuously to the IPC process.
In the context of the investigation of IPC processes also the energy distributions for positrons emitted after IPC was reconsidered. Our focus was particularly on positron spectra from several $`\gamma `$–transitions with energies between 1250 keV and 1600 keV observed with high resolution $`\gamma `$–spectroscopy in the Ta–like nucleus . The motivation of the present work was to find out if e<sup>+</sup>–emission from discrete nuclear transitions could give rise to lines in the e<sup>+</sup>–energy spectra with characteristics similar to previously reported e<sup>+</sup>–lines. In these former experiments line structures with differential cross sections of the order of 0.2 $`\mu `$b/sr up to about 1 $`\mu `$b/sr in Doppler-shift uncorrected e<sup>+</sup>–energy spectra were reported in several collision systems . Particularly in the system <sup>238</sup>U + <sup>181</sup>Ta , two very weak e<sup>+</sup>–lines at energies of $``$230 keV and of $``$310 keV have been observed at a beam energy of 5.9$`\times `$A MeV with production probabilities of $`(3.2\pm 0.8)\times 10^7/`$collision \[$`(0.40\pm 0.10)\mu `$b/sr\] and of $`(1.5\pm 0.5)\times 10^7/`$collision \[$`(0.23\pm 0.06)\mu `$b/sr\], respectively. The FWHM of the lines was (35 $`\pm `$ 11) keV and (13 $`\pm `$ 4) keV, respectively. Both lines appeared in coincidence with heavy ions scattered into a rather broad angular range ($`15^{}\theta _{ion}50^{}`$.
The recent systematic investigations allowed to study the response and sensitivity of our experimental setup to e<sup>+</sup>–spectra from IPC processes. Additionally, extensive Monte Carlo simulations, which take into account the progress made in the theoretical treatment of IPC in very heavy nuclei during the last few years , were performed to support our understanding of the IPC process. In these simulations we consider the complete kinematics of the collision system <sup>238</sup>U + <sup>181</sup>Ta at a beam energy of 6.3$`\times `$A MeV, as also studied in our experiments. Lepton pairs from an IPC transition are generated in the rest frame of the emitting ion with energy distributions taken from theoretical calculations. The energies of the leptons are transformed into the laboratory system and the simulated events are then analyzed with the same analysis program used to analyze the experimental data. The experimental acceptance of the setup is also considered in the simulation program. In a second step, the laboratory energies of the leptons are corrected event by event by taking into account the angular resolution of our setup by a Monte Carlo procedure.
The experiments were performed at the UNILAC accelerator of GSI, using an improved experimental setup at the double-Orange spectrometer. The experimental setup, the methods used as well as the Doppler-shift technique exploited have been described in details in our previous publications (see e.g. in ). In particular, the recent investigations of the collision system <sup>238</sup>U + <sup>181</sup>Ta were first presented and discussed extensively in , with emphasis on the appearance of weak e<sup>+</sup>e<sup>-</sup>–sum-energy lines due to IPC processes. Here we report the results from a complementary analysis of this collision system with the objective to search for corresponding line structures in the measured e<sup>+</sup>–spectra.
## 2 Experimental and simulation results
The collision system <sup>238</sup>U + <sup>181</sup>Ta was investigated using beams of <sup>238</sup>U with an energy of 6.3$`\times `$A MeV and 1000 $`\mu `$g/cm<sup>2</sup> thick <sup>181</sup>Ta targets. The $`\gamma `$-ray spectrum, measured with a Ge(i) detector at $`90^{}`$ to the beam direction and obtained after an event-by-event Doppler-shift correction to the Ta–like recoiling ions, is shown in Fig. 1a for R<sub>min</sub> values between 21.4 and 24 fm. As can be seen, several electromagnetic transitions are excited with energies between 1300 keV and 1600 keV. Their total excitation cross sections have values between a few mb and some 10 mb. These $`\gamma `$–transitions are obviously hitherto unknown in the <sup>181</sup>Ta–nucleus , but they were also measured by the EPOS and APEX collaborations as well as by Ditzel et. al. . The excitation probability, P<sub>γ</sub>(R<sub>min</sub>), was determined as a function of the distance of closest approach, R<sub>min</sub>, by normalizing the $`\gamma `$–yield in certain R<sub>min</sub> intervals with the corresponding number of elastically scattered ions. R<sub>min</sub> was derived from the measured scattering angle $`\theta _{ion,CM}`$ assuming Rutherford trajectories.
Figure 1b shows the excitation probability P<sub>γ</sub>(R<sub>min</sub>) as a function of R<sub>min</sub> for the strongest $`\gamma `$–transition at E$`{}_{\gamma }{}^{}=`$ (1380 $`\pm `$2) keV, which is representative for all $`\gamma `$–lines observed between 1300 keV and 1600 keV. The dependence of the excitation probability on R<sub>min</sub> is nearly constant in the R<sub>min</sub> range between 20 and 25 fm. Hence it is significantly different from the corresponding behaviour of the well known low-energy $`\gamma `$–transitions in <sup>181</sup>Ta due to Coulomb excitation, as demonstrated by Fig. 1c for the case of the 718 keV E2 transition. This indicates a different excitation mechanism of the high-energy $`\gamma `$–lines.
All transitions shown in Fig. 1a are accompanied by e<sup>+</sup>e<sup>-</sup>–pair-emission after IPC with total cross sections expected between some 0.1 $`\mu `$b and several $`\mu `$b, assuming transitions with multipolarities $`l>0`$ and IPC coefficients of the order of $`10^4`$. The energy partition between the positron and the electron of an IPC pair is determined by the final state interaction with the Coulomb field of the emitting nucleus. The emission probability for an IPC positron with a definite kinetic energy is described by the energy-differential pair conversion coefficient d$`\beta `$/dE$`_{\text{e}\text{+}}`$ which is a function of the nuclear charge number, the energy and the multipolarity of the transition. Figure 2a shows a Monte Carlo simulation of the expected emission probabilities for an IPC positron and the partner electron as a function of the e<sup>+</sup>–energy for the strong 1380 keV transition in the Ta–like nucleus which is the best candidate for an observation in our experiments. The results in Fig. 2a are given in the rest frame of the emitting nucleus taking into account the latest theoretical calculations on IPC assuming E1 multipolarity . Due to the final state interaction the positron of an IPC pair is preferentially emitted with the highest available kinetic energy of E$`{}_{\text{e}\text{+}}{}^{}=`$ E$`{}_{\gamma }{}^{}`$2m<sub>e</sub>c$`{}_{}{}^{2}=`$ 358 keV (solid curve), while for the partner electron the energy is complementary and given by E$`{}_{\text{e}\text{-}}{}^{}=`$ E$`{}_{\gamma }{}^{}`$2m<sub>e</sub>c$`{}_{}{}^{2}`$ E$`_{\text{e}\text{+}}`$ (dotted curve).
After transformation of the e<sup>+</sup>–energies into the laboratory system we obtain the distribution shown in Fig. 2b for R$`{}_{min}{}^{}=21.424`$ fm. If we now take into account the acceptance of the Orange e<sup>+</sup>–spectrometer at the chosen current settings, the distribution of Fig. 2b is reduced to the broad (FWHM $``$80 keV) line structure at $``$280 keV shown in Fig. 2c. The e<sup>+</sup>–spectrum was scanned by ramping the spectrometer current up and down within a preselected current interval in the same manner as it was done in the experiment. The current settings for the present experiment were such chosen that the IPC positrons with energies $`>`$280 keV were detected with maximum efficiency, whereas those with laboratory energies larger than the maximum possible center-of-mass (CM) e<sup>+</sup>–energy of 358 keV cannot be detected. This is due to the fact that positrons are only detected in the backward hemisphere with emission angles between $`110^{}`$ and $`140^{}`$ relative to the beam direction (see e.g. in ). The shape of Fig. 2c at energies below 280 keV is affected by the significantly reduced efficiency of the spectrometer in this region and by the low–energy cut–off at 200 keV. It should be underlined here that the present experiment was optimized for the detection of narrow e<sup>+</sup>e<sup>-</sup>–lines with sum energies around 630 keV, and not for the measurement of single positron spectra.
An event-by-event Doppler shift correction to the events of Fig. 2c leads to the CM distribution shown in Fig. 2d. The correction was performed by taking into account the finite angular resolution of the positron and heavy-ion detectors. Particularly, the positron emission angle was set to a constant value of $`125^{}`$ as also used in the analysis of the real data (for details see Ref ). In this case we obtain a line with a FWHM of $``$60 keV, whose maximum is now shifted to an energy of $``$320 keV relative to the uncorrected distribution.
According to these simulation results the IPC positrons from the 1380 keV transition should give rise to a rather narrow peak-like contribution in the measured e<sup>+</sup>–spectra with a FWHM of 60 to 80 keV and energies of $``$320 keV and $``$280 keV, for the CM and laboratory e<sup>+</sup>–energies, respectively. The CM distribution can be reconstructed by means of an event-by-event Doppler-shift correction .
To simulate the total contribution of IPC positrons from the Ta–like ion to the measured e<sup>+</sup>–spectra we have still to take into account the remaining $`\gamma `$–transitions observed between 1300 keV and 1600 keV, for which pure E1 multipolarity is assumed. Additionally we consider for the weak $`\gamma `$-transitions around 1500 keV a possible admixture of an E0 contribution to an E2 transition. This possibility was discussed in where a line structure in the e<sup>+</sup>e<sup>-</sup>–sum-energy spectra, obtained in <sup>238</sup>U + <sup>181</sup>Ta collisions, was found after Doppler-shift correction to the Ta–like ions. It exhibits the characteristics of close-lying IPC lines originating from the electromagnetic transitions between 1500 and 1550 keV and appears 8 times stronger than expected from the corresponding $`\gamma `$–spectrum by assuming E2 IPC coefficients of the order of $`10^4`$. It should be made clear that this empirical assumption does not influence the final results at all, and has been considered here only for a consistent treatment with our previous results .
Including these considerations we obtain from our simulations the IPC e<sup>+</sup>–energy distributions shown in Fig. 3. Figure 3a shows the expected e<sup>+</sup>–spectrum when an event-by-event Doppler-shift correction to the Ta–like recoils is applied, whereas Fig. 3b shows the corresponding e<sup>+</sup>–spectrum in the laboratory system. As can be seen, in the Doppler-shift corrected spectrum (Fig. 3a) the 60 keV broad line-structure at $``$320 keV, resulting from the 1380 keV transition, is still clearly distinguishable. Without Doppler-shift correction both distributions are shifted and smeared out, such that only two low intensity narrow structures near 300 keV and 400 keV are still visible (Fig. 3b).
The measured e<sup>+</sup>–energy spectra together with the calculated IPC e<sup>+</sup>–energy distributions from the above discussed electromagnetic transitions are presented in Fig. 4. These spectra are shown in coincidence with only one ion (Ta–like), scattered in the R<sub>min</sub> range 21.4 to 24 fm, but without requiring a coincidence with an electron. The underlying R<sub>min</sub> range was chosen in order to optimize the signal-to-background ratio for the IPC positrons. The solid line in Fig. 4a shows the measured Doppler-shift corrected e<sup>+</sup>–energy distribution while the dotted line represents the expected distribution of the IPC positrons from the observed $`\gamma `$–transitions in the Ta–like nucleus(cf. Fig. 3a) scaled up by a factor of 10.
The expected narrow structure from the 1380 keV transition around 320 keV cannot be observed in the measured spectrum. The IPC production probability from this transition is $`(2.1\pm 0.5)\times 10^7`$ per elastically scattered ion for R$`{}_{min}{}^{}=`$ 21.4–24 fm, as calculated from the corresponding measured $`\gamma `$–transition probability and an IPC coefficient of $`10^4`$. Note here that the expected e<sup>+</sup>–line energy and its production probability are very close to the values, measured for an e<sup>+</sup>–line with an energy of $``$310 keV in the previous experiments .
For our latest measurements the statistical detection limit for IPC positrons from the 1380 keV transition is $`2.6\times 10^7/`$collision. It is extracted from the measured continuous spectra assuming a superimposed IPC e<sup>+</sup>–line with 60 keV FWHM and two standard deviations statistical significance. In this case the IPC production probability is lower than the statistical detection limit and thus consistent with the non-observation of the expected IPC line. The production probability of the continuous e<sup>+</sup>–distribution itself amounts to $`(4.20\pm 0.03)\times 10^6`$ per collision for e<sup>+</sup>–energies between 290 and 350 keV. It originates from unresolved $`\gamma `$–transitions and collision induced atomic positrons .
In the Doppler-shift uncorrected spectra (Fig. 4b) the possibility to observe lines from IPC positrons is even worse. The expected IPC e<sup>+</sup>–distribution represented by the dotted line in Fig. 4b (also scaled up by a factor of 10) shows only two slight structures around 300 keV and 400 keV which are too weak as to be detected in the measured spectra. They are dominated by the continuous positron distribution.
It is worth mentioning at this point that the $`\gamma `$-ray spectra, obtained after Doppler-shift correction to the Ta–like ion, show another pronounced structure composed of three rather close lying $`\gamma `$–lines at energies of $`1240`$ keV, $`1260`$ keV and $`1275`$ keV (not shown in Fig 1a). The positron production probability expected from these transitions is $`1.5\times 10^7/`$ collision using an IPC coefficient of $`0.5\times 10^4`$. In the present experiment one could expect from these transitions a $`50`$ keV broad e<sup>+</sup>–line structure centered around 235 keV in the laboratory energy spectra, very close to the previously reported e<sup>+</sup>–line at $``$230 keV. But due to the very low spectrometer efficiency ($`ϵ2\times 10^3`$) for low positron energies in this experiment the expected signature of a 235 keV line is far below the detection limit.
Figure 5a shows the $`\gamma `$-ray spectrum corrected for Doppler shifts, assuming an emission from the U–like ions, scattered in the R<sub>min</sub> range 19–29 fm. Only some very weak $`\gamma `$–lines with energies between 1300 keV and 1500 keV can be observed in the above R<sub>min</sub> range. For instance, one can mark two lines at (1364 $`\pm `$ 7) keV and (1402 $`\pm `$ 7) keV with differential cross sections close to 1 mb/sr. The corresponding IPC contribution from these two transitions to the e<sup>+</sup>–energy spectrum is very small, i.e. $``$$`6\times 10^9/`$collision for an IPC coefficient of 0.5$`\times 10^4`$. It is obvious that this weak IPC contribution is undetectable in the measured spectra. The latter are shown in Figs. 5b and Fig. 5c, without and after a Doppler-shift correction, respectively. The spectra were obtained in coincidence with only one scattered ion (U–like) for the R<sub>min</sub> range from 19 to 29 fm. From these spectra a detection limit of $`2.4\times 10^7/`$collision (2$`\sigma `$) is derived for an IPC line.
## 3 Summary and Conclusions
We studied positron emission after IPC processes of several high–lying $`\gamma `$–transitions in Ta– and U–like nuclei excited in <sup>238</sup>U + <sup>181</sup>Ta collisions at a bombarding energy of 6.3$`\times `$A MeV. Several $`\gamma `$–lines have been observed in the measured $`\gamma `$–spectra, taken after an event-by-event Doppler-shift correction to the Ta–like or U–like ions, with transition energies between 1300 keV and 1600 keV. These transitions can de-excite via the IPC branch with IPC coefficients of the order of $`10^4`$.
The goal of our investigations was to gain information about the shape and intensity of the IPC positron distributions, originating from the above discussed $`\gamma `$–transitions, which can be expected in the measured Doppler–shift corrected and uncorrected positron energy spectra. To accomplish it, extensive Monte Carlo simulations were in addition carried out, particularly for the strongest $`\gamma `$–transition at 1380 keV in the Ta–like nucleus. They revealed that, due to the angular and momentum acceptance of the Orange e<sup>+</sup>–spectrometer as used in this experiment, IPC from these transitions can lead to rather narrow peak-like contributions to the e<sup>+</sup>–energy spectra. The line-structures are expected at e<sup>+</sup>–energies between 280 and 350 keV and have widths of $``$60 keV up to $``$80 keV (FWHM), in the Doppler-shift corrected and uncorrected energy spectra, respectively. Only some $`\gamma `$–transitions in the Ta–like nuclei appear with sufficient strength to give rise to IPC positron lines with production probabilities close to the detection limit. While in the U–like nucleus no transition was found which could give rise to any detectable positron line.
The observations discussed above are of particular interest with respect to the appearance of weak positron lines in our previous experiments. More specifically, in the collision system <sup>238</sup>U + <sup>181</sup>Ta , measured at a somewhat lower bombarding energy of 5.9$`\times `$A MeV, two lines at $``$230 keV and $``$310 keV with production probabilities of $`3.2\times 10^7/`$collision and $`1.5\times 10^7/`$collision, respectively, were reported from the Doppler-shift uncorrected spectra . From the present findings, an IPC origin of these e<sup>+</sup>–lines cannot be excluded. Indeed, as shown above, the measured $`\gamma `$-transitions in the Ta-like nucleus would lead to IPC e<sup>+</sup>–lines with production probabilities ($`10^7/`$coll.) and energies ($``$230 keV and $``$320 keV) which are consistent with Monte Carlo Simulation results based on these IPC transitions. This does not apply, however, to the very first e<sup>+</sup>–lines found in the collision system U+U with rather high production probabilities ($`10^5/`$coll.) which are not consistent with the intensities expected from the $`\gamma `$-transitions measured in this experiment.
It should be noted again that a direct comparison between the previous and the present experiment is problematic due to the following facts: First in the present experiment, the spectrometer current settings were different and optimized for an electron-positron coincidence experiment. As mentioned above, in this case energies below 200 keV were completely suppressed and the detection efficiency for positrons with energies between 200 keV and 300 keV was reduced considerably in comparison to the former experiments. Second no high-resolution $`\gamma `$-ray spectra were available in the old experiments to reveal weak discrete IPC transitions, having probably underestimated their role at that time.
Our recent investigations of e<sup>+</sup>–emission with the Orange setup, equipped with a high-resolution $`\gamma `$-ray detector, revealed that discrete nuclear transitions, populated via Coulomb excitation or in nuclear transfer reactions, can indeed appear in the outgoing moving nuclei with energies above 1 MeV. IPC processes from these transitions can lead to rather weak and narrow line-like structures in the measured positron energy spectra. Due to the dominating continuous e<sup>+</sup>-spectral distributions appearing in heavy-ion collisions, the detection of these IPC e<sup>+</sup>–lines is rather difficult and depends strongly on the special features of the setup used as well as on the experimental sensitivity and techniques exploited.
Although the IPC scenario studied in the present experiment provides evidence for a rather simple explanation for some of our previously reported weak e<sup>+</sup>–lines, it is clear that a definite proof of this scenario as an overall explanation of all the previously reported e<sup>+</sup>\- and e<sup>+</sup>e<sup>-</sup>-lines would require new dedicated experiments, which to our knowledge are not planned neither at GSI nor elsewhere.
Acknowledgement: We would like to thank all the people of the UNILAC accelerator operating crew for their efforts in delivering stable <sup>238</sup>U beams with high intensities.
* present address: INFN Section of Turin, I–10125 Turin, Italy and CERN, CH–1211 Geneva 23, Switzerland
Figure Captions
Fig. 1. a) Doppler-shift corrected $`\gamma `$–ray spectrum observed in the collision system <sup>238</sup>U + <sup>181</sup>Ta at a beam energy of 6.3$`\times `$A MeV. Ta–like recoiling ions for rather peripheral collisions (R$`{}_{min}{}^{}=`$ 21.4 – 24 fm) are assumed to be the emitter. Several lines appear at energies below 1600 keV (see also Fig.8 of Ref. ).
b) Excitation probability of the strongest $`\gamma `$–line at $``$1380 keV as a function of R<sub>min</sub>. These data were first presented in Ref. .
c) Excitation probability for a well known E2 $`\gamma `$–transition at 718 keV in a lower-lying rotational band in $`{}_{}{}^{181}Ta`$ (not shown in 1a) as a function of R<sub>min</sub>.
Fig. 2. Monte Carlo simulation results for energy distributions expected for e<sup>+</sup> and e<sup>-</sup> emission after IPC in a nucleus with a charge number of Z$`=`$73 and an electromagnetic transition with an energy of E$`{}_{\gamma }{}^{}=`$ 1380 keV and E1 multipolarity.
a) Original e<sup>+</sup> (solid line) and e<sup>-</sup> (dotted line) energy distributions expected in the rest frame of the emitting nucleus at the emission time. The sum of their kinetic energies has always a constant value which is given by the energy of the $`\gamma `$–transition.
b) The original e<sup>+</sup>–energy distribution after transformation into the laboratory system to account for the Doppler shift due to the motion of the emitting ion.
c) The corresponding e<sup>+</sup>–energy distribution in the laboratory system by taking into account the chosen momentum acceptance of the spectrometer.
d) The same events as in 2c, but after an event-by-event Doppler-shift correction has been applied utilizing the angular resolution of the setup.
Fig. 3. Monte Carlo simulation of the e<sup>+</sup>–energy distribution expected from several IPC transitions in Ta (Z$`=`$73) with energies between 1250 keV and 1600 keV, as suggested by the $`\gamma `$-ray spectrum shown in Fig. 1a. All the $`\gamma `$–transitions are assumed to be of pure E1 multipolarity, with the exception of the small contribution of those appearing between 1500 keV and 1600 keV for which an admixture of E2 and E0 with a mixing ratio of 1:8 is taken into account.
a)The e<sup>+</sup>–energies are corrected for Doppler shifts, assuming an emission from the Ta–like ions.
b) The corresponding e<sup>+</sup>–energy distribution expected in the laboratory system is shown.
Fig. 4. a) Measured e<sup>+</sup>–energy spectrum from <sup>238</sup>U + <sup>181</sup>Ta collisions at a beam energy of 6.3$`\times `$A MeV. The dotted-line histogram indicates the e<sup>+</sup>–contribution, multiplied by 10, which can be expected from IPC of excited states in the Ta–like nucleus (cf. Fig. 3a). Both spectra are corrected for Doppler shifts, assuming an emission from the Ta–like recoils scattered in the R<sub>min</sub> range between 21.4 and 24 fm.
b) The corresponding spectra obtained in the laboratory system are shown.
Fig. 5. a) Measured $`\gamma `$-ray spectrum, corrected for Doppler shifts, assuming an emission from the U–like ions scattered in the R<sub>min</sub> range from 19 to 29 fm.
b) Measured e<sup>+</sup>–energy spectrum in the laboratory system for the above R<sub>min</sub> range.
c) The same as in Fig. 5b, but after Doppler-shift correction to the U–like scattered ions.
Figure 1
(S. Heinz et al., EPJ A)
Figure 2
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Figure 3
(S. Heinz et al., EPJ A)
Figure 4
(S. Heinz et al., EPJ A)
Figure 5
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# On the mass of moderately rotating strange stars in the MIT bag model and LMXBs
## 1 Kilohertz QPOs and the mass of LMXBs
The discovery of kHz QPOs (Strohmayer et al., 1996; van der Klis et al., 1996) in several low-mass X-ray binaries (LMXBs) has renewed interest in the maximum mass of neutron stars, as its value limits the maximum observable orbital frequency (Kluźniak et al., 1990). The maximum masses for neutron stars modeled with various equations of state (e.o.s.) are well established (Arnett and Bowers, 1977; Friedman et al., 1986; Cook et al., 1994; Salgado, et. al., 1994), and have been examined in the context of kHz QPOs (Kluźniak, 1998). In principle, LMXBs could contain strange stars instead (Cheng and Dai 1996), and it has been asked whether the observed values of kHz QPOs are compatible with the masses of such stars (Bulik et al., 1999a,b).
The existence of strange stars would shed light on the issue of stability of quark matter. It has been suggested that a quark fluid, composed of roughly equal number of up, down and strange quarks, may be the ground state of bulk matter (Bodmer, 1971; Witten, 1984) and a detailed descripion of such strange matter has been given by Farhi and Jaffe (1984). This idea has not been universally accepted, and it was argued that for realistic values of the QCD coupling constant, the phase transition to quark matter would occur at densities too high to be of interest (e.g., Bethe, Brown and Cooperstein, 1987).
Strange stars may be very hard to distinguish from neutron stars, particularly in LMXBs, as they may have a crust of normal matter (Alcock et al., 1986; Haensel et al., 1986). The crust will contribute little to the mass ($`10^5M_{}`$), but is expected to have a thickness sufficient to support nuclear burning, including helium flashes responsible for X-ray bursts, thus mimicking neutron stars. The heat released in conversion (at the bottom of the crust) of nuclei into strange matter is directed into the strange matter core (Miralda-Escudé et al., 1990; Haensel and Zdunik, 1990). At any rate, steady release of energy from nuclear conversion is very difficult to distinguish from the gravitational binding energy released in steady accretion.
However, in young radio pulsars the crust thickness does make a difference—the crust in strange stars would be far too thin to allow the redistribution of angular momentum necessary to explain the glitch phenomenon (sudden spin-up of a pulsar), therefore glitching pulsars are though to be neutron stars and not quark stars (Alpar, 1987). It is likely that the coalescence of two strange stars would lead to the contamination of our Galaxy with chunks of strange matter, and would thus preclude the formation of young neutron stars (Madsen, 1988; Caldwell and Friedman, 1991). These arguments make unlikely the presence of strange stars in the population of ordinary pulsars (including binary Hulse-Taylor type pulsars) or their accreting counterparts, the high-mass X-ray binaries. However, the presence of strange stars among the millisecond pulsars or LMXBs seems to be allowed (Kluźniak, 1994; Cheng and Dai, 1996). In these old systems, the stellar mass itself may yield clues as to the nature of the compact object.
We are led to consider the maximum mass of strange stars in the expectation that LMXB masses will become available in the near future, either through a better understanding of accretion phenomena, including the observed quasi-periodic oscillations (QPOs) in the X-ray light curve, or through classical optical determinations of binary parameters for newly discovered transient sources.
The maximum frequency of stable orbital motion is attained in the innermost (marginally) stable circular orbit (ISCO) allowed by general relativity, and if its value for some X-ray source were that of the observed maximum QPO frequency—ranging from 1.0 kHz to 1.2 kHz in twelve LMXBs—conclusions could be reached about the nature of the compact object. \[For a review of QPOs see e.g. van der Klis (1998).\]
The question, whether ISCO frequencies as low as the maximum observed QPO frequencies can be attained outside quark stars, has been answered in the affirmative for rapidly rotating strange stars, with an equation of state (e.o.s.) based on the MIT bag model of quark matter—strange stars rotating close to the equatorial mass-shedding limit can have ISCO frequencies below 1 kHz for masses as low as $`1.4M_{}`$ (Stergioulas et al., 1999). However, the same question has not yet been answered for moderately rotating strange stars, of periods $`P3`$ms or more, for which such ISCO frequencies would imply larger masses: $`M=2.2M_{}(1+0.75j)(1.0\mathrm{kHz}/f_{max})`$, where $`f_{max}`$ is the ISCO frequency, and $`j0.3`$ is the dimensionless angular momentum (Kluźniak et al., 1990; Kluźniak, 1998).
There is a compelling reason to consider stars rotating at relatively low rates. The period of the compact star in the the transient LMXB SAX J 1808.4-3658 (Wijnands and van der Klis, 1998), the first of possibly many such transients to be discovered, has been measured to be 2.5 ms. Whether or not kHz QPOs will be discovered in that source, its mass may eventually be determined by optical studies of the binary companion.
It has also been argued that the oscillations seen in some X-ray bursters imply a stellar rotational frequency of $`300`$Hz (Strohmayer et al., 1997). Finally, it is not yet clear whether strange stars can attain the equatorial mass-shedding limit because of the unusually high value of $`T/W>0.2`$ calculated for their models—note that for Newtonian stars the secular instability to a bar-mode deformation sets in already at $`0.1275T/W0.1375`$ (Bonazzola et al., 1996). For a discussion of other modes, also in general relativity, see, e.g., Gourgoulhon et al. (1999), and references therein.
For all these reasons, we considered exact numerical models of strange stars rotating with frequencies up to 700 Hz, and found that, to within a few per cent, the maximum mass of such strange stars is the same as that of the static configurations. Therefore, we investigate the maximum mass of static (non-rotating) strange stars, which till now has been discussed in the context of QPOs only in the simplest case of ideal quark gas in the bag model (Bulik et al., 1999a). Here, we use the e.o.s. of interacting, massive quarks within the MIT bag model of self-bound quark matter.
## 2 Strange stars
The maximum mass of quark stars was first derived (Brecher and Caporaso, 1976) for an unusually low value of the bag constant. Witten (1984) showed that the maximum mass of a static strange star is, within the MIT bag model (Farhi and Jaffe 1984),
$$M_{max}=1.98M_{}(B/B_1)^{1/2},$$
(1)
where $`B_1=58.9`$MeV$``$fm<sup>-3</sup>.
There seems to have been no systematic investigation of the maximum mass of the star as a function of the three basic parameters of quark matter in the MIT bag model: the mass of the strange quark, $`m_s`$, the bag constant $`B`$, and the strength of the QCD coupling constant, $`\alpha _c=g^2/4\pi `$. Detailed models of strange stars have been constructed and the structure of strange stars, including the mass–radius relationship has been discussed extensively in the literature (Alcock et al., 1986; Haensel et al., 1986; Glendenning, 1989; Frieman and Olinto, 1989; Prakash et al., 1990), however discussion of stellar parameters tended to concentrate on the maximum value of the bag constant, setting a lower bound on the maximum mass and an upper bound on the rotational frequency.
In this section, we discuss only static stellar models, constructed in general relativity by solving the TOV equations (Oppenheimer and Volkoff, 1939). We neglect the crust, whose contribution to the maximum mass is quite minor for stars composed mostly of self-bound quark matter (Haensel et al., 1986). Following Farhi and Jaffe (1984), we take $`m_s`$ and $`\alpha _c`$ to be renormalized at $`q=313`$ MeV, and as a model for quark matter consider a “bag” of positive vacuum energy density, $`B`$, filled with quarks (of two massless flavors and one massive) having interactions through first order in $`\alpha _c`$. The actual form of the thermodynamic expressions we use can be found in Haensel et al. (1986). As both the pressure and the density scale with the bag constant, TOV equations imply that the stellar mass and radius scale as $`MB^{1/2}`$, and $`RB^{1/2}`$ (Witten, 1984), provided that the masses of the quarks scale as $`mB^{1/4}`$.
In Figure 1, we plot the rescaled value of the maximum stellar mass, $`M(B/B_1)^{1/2}`$, as a function of the rescaled strange quark mass, $`m_s(B_1/B)^{1/4}`$, for various values of the QCD coupling constant. Note that the highest value of $`M_{max}`$, the maximum mass of the strange star, is independent of $`\alpha _c`$, if $`B`$ is fixed, and is obtained for $`m_s=0`$, i.e., for massless quarks. But in fact, the value of $`B`$ is not known, and as its lower bound does depend on the value of $`\alpha _c`$, the actual physical bounds on the maximum mass of a strange star will depend, through $`B`$, on the coupling constant. This is discussed in the next section.
## 3 The minimum value of $`B`$ and an upper bound to the mass of moderately rotating strange stars
To determine the highest possible value of the maximum mass of static strange stars in the MIT bag model, it is enough to consider the e.o.s. of an ulrarelativistic Fermi gas in a volume with vacuum energy density $`B>0`$. If the quarks are massive, the actual maximum mass of the star will be somewhat lower—as is evident from Figure 1, at fixed values of the other parameters, the maximum mass of a strange star decreases with increasing quark mass.<sup>1</sup><sup>1</sup>1The dependence of $`M`$ on $`m_s`$ is easily understood in the limit $`\alpha _c=0`$. The mass of a static star modeled with ideal Fermi gas e.o.s. scales with fermion mass as $`Mm^2`$ and attains a finite limit at $`m0`$ (Oppenheimer and Volkoff, 1939).
Currently, the actual value of $`B`$ cannot be reliably derived from fits to hadronic masses of the quark-model of nucleons. Instead, we must rely on a different argument to find $`B_{min}`$. We require that neutrons do not combine to form plasma of deconfined up and down quarks, or equivalently, that quark matter composed of up and down quarks in 1:2 ratio is unstable to emission of neutrons through the reaction $`u+2dn`$ (e.g., Haensel, 1987; Farhi, 1991), i.e., that the baryonic chemical potential at zero pressure of such quark matter satisfies
$$\mu _0^{u,d}>939.57\mathrm{MeV}.$$
(2)
This condition provides a lower bound on the bag constant $`B`$, and consequently, by eq. (1), an upper bound on the mass of static strange stars. We take the up and down quarks to be massless. Therefore, eq. (2) implies $`BB_1(12\alpha _c/\pi )`$. For noninteracting and massless quarks ($`\alpha =0`$, $`m_s=0`$), this value of the bag constant corresponds to a minimum density of strange matter (attained at zero pressure) of $`\rho _0(0)=4.20\times 10^{14}`$g/cm$`{}_{}{}^{3}4B_1/c^2`$, with a corresponding maximum mass of non-rotating strange stars of $`M_{max}(0)=1.98M_{}`$.
For interacting (but massless) quarks, through lowest order in gluon exchange, the equation of state is identical to that of non-interacting quarks (Chapline and Nauenberg, 1976), $`p=(\rho \rho _0)c^2/3`$, the only difference being in that the lower bound on the density at zero pressure, following from conditions of neutron stability (eq. ), is decreased with respect to the value for an ideal Fermi gas in a bag, by the same factor as the bag constant: $`\rho _0(\alpha _c)=\left(12\alpha _c/\pi \right)\rho _0(0)`$. Since the stellar mass scales as $`\rho _0^{1/2}`$, this implies that the least upper bound on the mass of the star as a function of the QCD coupling constant is given for non-rotating strange stars by
$$M_{max}(\alpha _c)=\left(1\frac{2\alpha _c}{\pi }\right)^{1/2}M_{max}(0),$$
$`(3)`$
through first order in $`\alpha _c`$ (Prakash et al., 1990). For $`\alpha _c=0.4`$, eqs. (2), (3) give a maximum strange star mass of $`2.29M_{}`$, higher by 16% than the maximum mass which is obtained for $`\alpha =0`$.
Stellar rotation at a frequency up to 700 Hz would increase the maximum mass by only a few percent. In Fig. 2, we present the maximum mass of a strange star as a function of the rotational frequency, $`f=1/P`$. This is not the frequency dependence of the mass of a single star—what is plotted is the termination point, at each frequency, of a sequence of stellar models varying in mass. Thus, the baryon number of the maximum-mass model varies with frequency. In fact, for all $`f>0`$, the maximum-mass stars are supramassive—if spun down at constant baryon number to $`f=0`$, they would become unstable to collapse (compare the discussion of neutron-star models in Cook et al., 1994). We obtained the results presented in Fig. 2 with an accurate code based on spectral methods, developed by, and described in, Gourgoulhon et. al. (1999).
## 4 The maximum mass of non-rotating strange stars as a function of $`m_s`$ and $`\alpha _c`$
To determine the maximum mass of such stars for specific values of model parameters ($`B,\alpha _c,m_s`$), the e.o.s. of matter composed of interacting, massive quarks must be considered, and it can only be determined numerically. The maximum stellar mass then following from the lower bound on $`B`$ implicit in eq. (2) is exhibited in Fig. 3 (continuous lines) as a function of the strange-quark mass and of the QCD coupling constant.
From Fig. 3 we can conclude for example, that if a strange star of two solar masses ($`2M_{}`$) were identified, and if $`\alpha _c`$ were less than 0.4, then $`m_s<200`$MeV (with both quantities renormalized at $`313`$MeV).
If the density at zero pressure of quark matter is indeed close to its lowest bound given by the stability limit of eq. (2), then for $`\alpha _c>0`$ it can attain a value lower than $`4\times 10^{14}`$g cm<sup>-3</sup> and the mass of even a non-rotating strange star may be sufficiently high to allow an orbital frequency of $`1.06`$kHz in the marginally stable orbit. It would appear that in the model of interacting quarks considered here, and in contrast to models in which $`\rho _0>4.2\times 10^{14}`$g cm<sup>-3</sup>, the possible presence of moderately rotating strange stars in LMXBs could be compatible with the keplerian model of kHz QPOs—compare Bulik et al. (1999a).
## 5 The maximum value of $`B`$
There is another stability constraint limiting the parameters of quark matter, if such a substance is indeed the ground state of matter. The chemical potential at zero pressure of quark matter in three flavours and electrons in beta equilibrium, should be less than the rest energy per nucleon of the iron nucleus $`\mu _0^{u,d,s}<\mu (^{56}\mathrm{Fe})=930.4\mathrm{MeV}`$ (Farhi and Jaffe 1984). For massless quarks this corresponds to $`B/(12\alpha _c/\pi )<91.49\mathrm{MeV}`$, but in general the constraint depends on the strange quark mass, and the maximum stellar mass corresponding to this upper bound on $`B`$ is exhibited in each panel of Fig. 3 as a dotted line \[for details see Prakash et al. (1990) and Zdunik et al. (2000)\]. Note that the upper and lower constraints on $`B`$, when taken together, exclude high masses of the strange quark for low values of $`\alpha _c`$, if strange quark gas is to be the stable form of matter.
## 6 Summary
We have investigated the question whether moderately rotating strange stars with masses somewhat higher than $`2M_{}`$ are allowed by relativistic equations of stellar equilibrium, and found that the answer could, in principle, be positive within the MIT bag model of beta-stable quark matter. An extension of this discussion to other models of strange matter will be given elsewhere.
The physical constraints on the bag constant following from the basic hypothesis of stability of self-bound quark matter (eq. ) allow $`B`$ to be so small, that the corresponding mass of the star could be as high as $`2.5M_{}`$ (eq. , Fig. 3). However, the strange star mass cannot be higher than $`2.6M_{}`$, even for stars of rotational period as short as 1.6 ms.
The perturbative approach used here is sensible only if the value of the QCD coupling constant $`\alpha _c`$ is small—we used $`0\alpha _c0.6`$. In this range, the direct dependence of $`M_{max}`$ on the coupling is practically negligible (Fig. 1), but the window of allowed values of $`B`$ does depend on $`\alpha _c`$ (Fig. 3). The actual value of $`B`$ is subject to a very large uncertainty. Fits to the hadronic mass spectrum (DeGrand et al., 1975) gave $`B=59\mathrm{MeV}\mathrm{fm}^3`$.
The stellar mass decreases with increasing mass, $`m_s`$, of the strange quark, and is lower by $`10\%`$ to $`20\%`$ than the one for massless-quark matter for the typical range considered in the literature, $`150m_sc^2/\mathrm{MeV}300`$ (Madsen 1999). Finally, the maximum stellar mass for strange stars rotating with a frequency up to $`700`$Hz, is larger than the one for non-rotating stars by less than 4% (Fig. 2).
The physical reason for which large values of mass are obtained for a low value of the bag constant, is that the latter is proportional to the density of quark matter at zero pressure, $`\rho _0`$, while the maximum mass of the star is proportional to $`\rho _0^{1/2}`$. Thus, within the MIT bag model of quark matter, strange stars with mass $`M>2.0M_{}`$ must have surface densities $`\rho _0<\rho _14.2\times 10^{14}`$g/cm<sup>3</sup> (unless rotating with periods $`1.5`$ms).
This research was supported in part by KBN grants 2P03D00418, 2P03D02117, 2P03D04013. The numerical calculations have been performed in part on DARC computers purchased thanks to a special grant from the SPM and SDU departments of CNRS.
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# Quantum Clock Synchronization Based on Shared Prior Entanglement
## Abstract
We demonstrate that two spatially separated parties (Alice and Bob) can utilize shared prior quantum entanglement, and classical communications, to establish a synchronized pair of atomic clocks. In contrast to classical synchronization schemes, the accuracy of our protocol is independent of Alice or Bob’s knowledge of their relative locations or of the properties of the intervening medium.
PACS: 03.67.-a, 03.67.Hk, 06.30.Ft, 95.55.Sh
In the Special Theory of Relativity, there are two standard methods for synchronizing a pair of spatially separated clocks, A and B, which are at rest in a common inertial frame. The usual procedure is Einstein Synchronization (ES), which involves an operational line-of-sight exchange of light pulses between two observers, say Alice and Bob, who are co-located with their clocks A and B, respectively . A less commonly used protocol is Eddington’s Slow Clock Transport. In this scheme, the two clocks A and B are first synchronized locally, and then they are transported adiabatically (infinitesimally slowly) to their final separate locations . A quantum algorithm for efficient clock transport has recently been proposed by Chuang .
In this paper we propose a third protocol that utilizes the resource of shared prior entanglement between the two synchronizing parties. Our proposed method of Quantum Clock Synchronization (QCS) has features in common with Ekert’s entanglement-based quantum key-distribution protocol in which Alice and Bob initially share only prior-entangled qubit pairs. The key does not exist initially but is created from the ensemble of entangled pairs through a series of measurements and classical messages. Similarly for our QCS protocol below, no actual clocks exist initially but rather only “entangled clocks” in a global state which does not evolve in time. The synchronized clocks are then extracted via measurements and classical communications performed by Alice and Bob. In this way our QCS scheme establishes synchrony without having to transport timing information between Alice and Bob. In contrast, in classical synchronization schemes, actual timing information must be transmitted from Alice to Bob over some channel, whose imperfections generally limit the accuracy of the synchronization.
We begin by reviewing the Ramsey temporal interferometer method for the construction of a quantum clock. A clock is constructed from an ensemble of two-level systems (qubits) whose temporal evolution properties will determine the time standard. In general any physical qubit may be used. For example, in the International System of Units (SI), the second is defined as the duration of exactly 9,192,631,770 periods of oscillation corresponding to the hyperfine (radio) transition frequency for the ground-state of the Cs<sup>133</sup> atom . Below we also consider other possible physical realisations of the qubit.
Consider a qubit with stationary states $`|0`$ and $`|1`$ having energy eigenvalues $`E_0<E_1`$ respectively. We introduce the dual basis $`|pos=\frac{1}{\sqrt{2}}(|0+|1)`$ and $`|neg=\frac{1}{\sqrt{2}}(|0|1)`$ and write $`\mathrm{\Omega }=\frac{1}{\mathrm{}}(E_1E_0)`$. The Hadamard transform $`H`$ is defined by the operation $`|0|pos`$ and $`|1|neg`$. Let us write $`\sigma _3`$ for the measurement in the $`\{|0,|1\}`$ basis and $`\sigma _1`$ for the measurement in the dual basis. Thus if the qubit is a spin $`\frac{1}{2}`$ particle in a $`z`$-oriented magnetic field then $`|0`$ and $`|1`$ are the $`z`$ spin eigenstates with $`\sigma _3`$ and $`\sigma _1`$ being $`\sigma _z`$ and $`\sigma _x`$ respectively. If the qubit comprises two hyperfine energy levels of a Cs<sup>133</sup> atom, then $`\sigma _3`$ measures population in these levels and $`\sigma _1`$ is measured by first applying $`H`$ and then measuring $`\sigma _3`$.
The Ramsey method for providing a time standard is based simply on the fact that the states $`|pos`$ and $`|neg`$ are not stationary states. They evolve in time as:
$$\begin{array}{ccc}|pos(t)\hfill & =& \frac{1}{\sqrt{2}}\left(e^{i\mathrm{\Omega }t/2}|0+e^{i\mathrm{\Omega }t/2}|1\right)\hfill \\ |neg(t)\hfill & =& \frac{1}{\sqrt{2}}\left(e^{i\mathrm{\Omega }t/2}|0e^{i\mathrm{\Omega }t/2}|1\right)\hfill \end{array}$$
(1)
At some time $`t=0`$ we apply $`H`$ to an ensemble of qubits in state $`|0`$ giving an ensemble of states $`|pos`$ which begin to evolve as in eq. (1). After a time $`t`$ we measure the observable $`\sigma _1`$ (either directly or by first applying $`H`$ and measuring $`\sigma _3`$, depending on the physical implementation of the qubits). A straightforward calculation shows that the probabilities for seeing outcomes 0 or 1 are given by
$$P(0)=\frac{1}{2}\left(1+\mathrm{cos}\left(\mathrm{\Omega }t\right)\right),P(1)=\frac{1}{2}\left(1\mathrm{cos}\left(\mathrm{\Omega }t\right)\right)$$
(2)
By monitoring the oscillations of either $`P(0)`$ or $`P(1)`$ as a function of time we get an estimate of the clock phase $`\mathrm{\Omega }t\text{mod}2\pi `$ and hence of $`t`$.
We now describe our proposed QCS scheme. We assume at the outset that Alice and Bob share an ensemble of singlet states $`|\psi ^{}=\frac{1}{\sqrt{2}}\left(|0_A|1_B|1_A|0_B\right)`$ where the subscripts refer to particles held by Alice and Bob. The pairs are labelled $`n=1,2,3,\mathrm{}`$ and the labels are known to both Alice and Bob. This singlet state is a “dark state” that does not evolve in time provided A and B undergo identical unitary evolutions. Indeed for any 1-qubit unitary $`U`$ we have $`(UU)|\psi ^{}=(detU)|\psi ^{}`$ so that $`|\psi ^{}`$ changes only by an overall unobservable phase. Our protocol below (slightly modified) would work equally well using the state
$$|\psi ^{}(\eta )=\frac{1}{\sqrt{2}}\left(|0_A|1_Be^{i\eta }|1_A|0_B\right)$$
(3)
for any fixed $`\eta `$. This state still has the essential property of being constant in time i.e. invariant under $`UU`$ where $`U`$ is time evolution, diagonal in the $`\{|0`$,$`|1\}`$ basis (but unlike the singlet, it is not invariant under $`UU`$ for more general $`U`$’s).
We will refer to a pair of clocks in the singlet state $`|\psi ^{}`$ as an entangled pair of pre-clocks. Since $`|\psi ^{}`$ is constant in time the pre-clock pairs could be said to be “idling” – they can provide no direct timing information. We may also write $`|\psi ^{}`$ in the $`\sigma _1`$ measurement basis as
$$|\psi ^{}=\frac{1}{\sqrt{2}}\left(|pos_A|neg_B|neg_A|pos_B\right)$$
(4)
Let $`t`$ be a time coordinate in the common rest frame of Alice and Bob. To start the clocks at some time $`t=0`$, Alice simultaneously measures all of her pre-clock pairs in the $`\sigma _1`$ basis $`\{|pos,|neg\}`$. Thus each pair collapses randomly and simultaneously at $`A`$ and $`B`$ into one of the following states:
$$\begin{array}{ccc}|\psi ^I\hfill & =& |pos_A|neg_B\hfill \\ |\psi ^{II}\hfill & =& |neg_A|pos_B\hfill \end{array}$$
(5)
with equal probability $`\frac{1}{2}`$. The $`A`$ and $`B`$ clocks begin to evolve in time, in accordance with Eq. (1) – all starting synchronously at a time $`t=0`$ in Alice and Bob’s shared inertial frame. Indeed Alice’s measurement effectively reproduces the result of the first one-clock Hadamard transform in the Ramsey scheme. However the result here is a mixture of two equally weighted sub-ensembles I and II. As a result of her measurement, Alice knows the labels belonging to the subensembles I and II but Bob is unable to distinguish them.
The density matrix of Bob’s overall ensemble is still $`\rho =\frac{1}{2}I`$, independent of $`t`$, so no measurement statistic can provide Bob with any timing information. For Bob to extract a clock, a classical message from Alice is required. Alice post-selects from her entire ensemble the sub-ensemble of Type-I qubits. Since the qubits are labelled, she can then tell Bob which subset of his qubits are also Type-I by broadcasting their labels via any form of classical communiqué. Bob is then able to extract his own Type-I and Type-II subensembles. Choosing the Type-II subensemble, Bob will have a clock ensemble exactly in phase with a Type-I clock that Alice started at $`t=0`$. To establish synchrony, Bob measures $`\sigma _1`$ on this ensemble (either directly or by applying $`H`$ and measuring $`\sigma _3`$) and monitors the oscillations of $`P(0)`$ as in eq. (2). Alice and Bob now have clocks that are ticking in unison.
The protocol as described above is still incomplete because of the following rather subtle point: there are extra hidden assumptions in the requirement that Alice and Bob are both able to perform the same $`H`$ operation and identify the same $`|pos`$ states. Indeed if we are given only $`|0`$ and $`|1`$ as physical states (i.e. normalised vectors up to overall phase) then the physical states $`|pos`$ and $`|neg`$ are not uniquely determined<sup>*</sup><sup>*</sup>*Note that if $`|0`$ and $`|1`$ are given as vectors then all other vectors such as $`|pos`$ are uniquely defined so our ambiguity depends essentially on the fact that a physical state is not just a (normalised) vector but rather, a set of all such vectors that differ by an overall phase. and so $`H`$ is also not uniquely determined (as, for example, it entails knowledge of $`|pos`$). A further arbitrary choice needs to be made, analogous to a choice of reference frame, to fix these further constructs.
The need for a further choice is most clearly seen by considering the spin $`\frac{1}{2}`$ qubit . The physical states $`|0`$ and $`|1`$ define a $`z`$ direction and $`|pos`$ defines a perpendicular $`x`$ direction. But given only a $`z`$ direction we are free to choose any orthogonal direction as $`x`$. On the Bloch sphere $`|0`$ and $`|1`$ are two given poles and $`|pos`$ may be arbitrarily chosen to be any point on the equator. Once $`|pos`$ is chosen, it must be consistently used in all future operations. Furthermore, there is then no further ambiguity in the identity of any state on the Bloch sphere e.g. $`|neg`$ and $`H`$ are then uniquely fixed.
The same remarks apply to the Cs atom qubit but the physical interpretation is quite different. The operation $`H`$ (and hence $`|pos`$) is physically defined in terms of a $`\pi /2`$ pulse applied to the physical state $`|0`$. But this pulse has an origin of phase which must be chosen and then fixed (“locked”) for all future applications of $`H`$. Different choices of phase locking correspond to different choices of points on the Bloch sphere equator for $`|pos`$. Note that a choice of phase locking here corresponds physically to a choice of time origin in contrast to the spin $`\frac{1}{2}`$ qubit, where the choice was a spatial direction.
For our QCS protocol to work correctly, Alice and Bob must use the same choice of physical state $`|pos`$ (or equivalently use the same choice of Hadamard operation $`H`$). If they use two different choices (and use them locally consistently) then their clocks will not be ticking in synchrony, but be offset by an amount depending on the angle between the two choices of $`|pos`$ on the Bloch equator.
In the physical implementation given by the Cs atom qubit, a consistent choice of $`H`$ requires that Alice and Bob have mutually phase locked pulses. But this is equivalent to them having clocks ticking in synchrony thus defeating the purpose of the protocol! However the following extension of our protocol gets around this difficulty, allowing Alice and Bob to establish time synchrony without the resource of mutually phase locked pulses: we duplicate our protocol above for two different values $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ of $`\mathrm{\Omega }`$ e.g. we use two different species of atoms. Thus Alice and Bob will require two different kinds of pulses for the two frequencies. In his laboratory, Bob is able to lock the phases of his two pulses, and similarly for Alice, so there will be a common offset $`\delta `$ between the two locked settings of Alice and Bob. By measuring populations in state $`|0`$ as before, Alice will have the oscillations:
$$P_1^A=\frac{1}{2}(1+\mathrm{cos}\mathrm{\Omega }_1t)P_2^A=\frac{1}{2}(1+\mathrm{cos}\mathrm{\Omega }_2t)$$
(6)
and Bob’s oscillations will be offset by the constant (unknown) $`\delta `$:
$$P_1^B=\frac{1}{2}(1+\mathrm{cos}(\mathrm{\Omega }_1t+\delta ))P_2^B=\frac{1}{2}(1+\mathrm{cos}(\mathrm{\Omega }_2t+\delta ))$$
(7)
But now by observing the beats between the two oscillations $`P_1`$ and $`P_2`$ Alice and Bob are able to establish synchronously ticking clocks. Indeed we have
$$P_1^BP_2^B=\mathrm{sin}(\frac{1}{2}(\mathrm{\Omega }_1\mathrm{\Omega }_2)t)\mathrm{sin}(\frac{1}{2}(\mathrm{\Omega }_1+\mathrm{\Omega }_2)t+\delta )$$
(8)
so that the envelope (given by the first term) oscillates independently of $`\delta `$ exactly in phase with Alice’s corresponding envelope.
It is interesting to consider the above problem, of locally consistent but different choices of $`|pos_A`$ and $`|pos_B`$, in the alternative physical scenario of clocks given by ensembles of spin $`\frac{1}{2}`$ qubits in a magnetic field. Although mathematically equivalent, we will see that the physical implications are quite different. In this scenario $`|0`$ and $`|1`$ are the $`z`$ spin eigenstates. We imagine that a third party (Clare) prepares an ensemble of pairs in the singlet state and simultaneously puts each spin in a labelled box containing a constant magnetic field $`B_z`$ in the $`z`$ direction. She then distributes the boxes (complete with their magnetic fields) to Alice and Bob (appropriately for each pair). Note that Alice and Bob may determine the $`z`$ direction (if they do not already know it) by measuring the direction of the (classical) magnetic field in a box (without disturbing the particle).
Alice now chooses an $`x`$ direction (perpendicular to $`z`$) and at some time $`t=0`$ she measures $`\sigma _x`$ on all her particles. Then, just as before, Bob may establish synchrony by monitoring the oscillations of $`\sigma _x`$ measurement outcomes on a sub-ensemble of his particles (selected by classical information from Alice). The previous problem of consistent phase locked pulses becomes the problem of Bob choosing the same $`x`$ direction that Alice used. Previously, the problem was equivalent to the original goal of the protocol (time synchrony) but here it is different i.e. a requirement of space parallelism (“space synchrony”). This allows the possibility of new physical resolutions of the problem, not available for Cs atom qubits. For example, Alice and Bob may have a prior agreement to use the direction to the pole star as their $`x`$ direction (which would be parallel to high accuracy for any two locations on Earth) i.e. $`x`$-“space synchrony” may be given for free, whereas time synchrony is not.
The idea of the previous resolution – using two frequencies – may be used for spin $`\frac{1}{2}`$ qubits as well (e.g. if Alice and Bob are unable to see any fixed stars.) Clare sets up boxes with two different magnetic fields (both in the $`z`$ direction) giving the two different frequencies. Bob chooses his $`x`$ axis randomly (perpendicular to $`z`$) and the constant phase offset $`\delta `$ now arises from the fixed angle between Alice’s and Bob’s chosen $`x`$ directions. An important point here is that different physical realisations of a qubit – although mathematically equivalent – lead to quite different avenues for getting around limitations of a (mathematically) given protocol.
For some applications, such as satellite-based Very Long Baseline Interferometry (VLBI) , the fact that Alice and Bob’s clocks are phase locked up to only modulo $`2\pi `$ is sufficient. However, there are other applications, such as the synchronization of satellite-borne atomic clocks in the Global Positioning System (GPS) , where it is important to have a shared origin of time. For such applications, we may adapt our QCS protocol to construct a common temporal point of reference as follows. Using the protocol Alice and Bob set up clocks ticking synchronously for two different frequencies $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_1+\mathrm{\Delta }\mathrm{\Omega }`$. The envelope of beats between these frequencies oscillates with frequency $`\frac{1}{2}\mathrm{\Delta }\mathrm{\Omega }`$. If the protocol for establishing the two ticking synchronisations is completed in time $`T`$ and $`\mathrm{\Delta }\mathrm{\Omega }`$ is chosen so that $`\mathrm{\Delta }\mathrm{\Omega }T<\frac{\pi }{2}`$ then Alice and Bob may determine a unique common time origin as the first maximum of the beat oscillations.
There are several immediate applications and advantages of our QCS protocol. For example, in the GPS satellite constellation, the ability of the space-borne atomic clocks to synchronize with a master atomic clock on the ground is affected by the fluctuating refractive index of the atmosphere, causing random variations in the speed of light and limiting the accuracy of the classical ES protocol. This index fluctuation error is the current limiting factor of GPS precision . With our QCS scheme, the properties of the atmosphere have no effect. In fact, Alice and Bob need not even have exact knowledge of their relative locations.
Also classical ES requires the exchange and timing of light pulses, but light is actually a quantum field. Hence the arrival time of a light pulse is itself subject to quantum fluctuations, limiting the accuracy of the ES protocol. In contrast, our QCS scheme is unaffected by this kind of noise.
The Ramsey two-pulse temporal interferometer is isomorphic, via the SU(2) algebra, to an optical or matter-wave Mach-Zehnder interferometer. Hence, the QCS protocol may be readily adapted to the task of phase locking a pair of spatially separated optical or atom interferometers, with applications to various forms of interferometry such as VLBI.
A shortcoming of our QCS protocol that it does not specify a method by which the shared prior entanglement between Alice and Bob may be established. One possibility is for Alice and Bob to meet at a common location, create an ensemble of $`N`$ identical EPR pairs each in the state $`\frac{1}{\sqrt{2}}\left(|0_A|1_B|1_A|0_B\right)`$ and then go their separate ways. But then, could not Alice and Bob just establish time synchrony at their meeting and retain accurate clocks for future use, instead of carrying the entanglement? In practice, clocks drift and periodic corrections of synchronisation will be necessary, suffering from the limitations of the classical schemes. It is not clear whether the task of carrying and maintaining the required entanglement is equivalent to the task of carrying and maintaining accurately running clocks. One difference is that in the former case, the time synchrony does not initially exist but is set up only when required, which may have applications for security.
An alternative scheme for establishing the shared prior entanglement would not require Alice and Bob to meet at all. Instead, it would involve Alice and Bob each receiving corresponding members of EPR pairs from some common source and then using entanglement purification to distill them into an ensemble of singlet states as required by QCS. Unfortunately, there is a hidden assumption of simultaneity in the actions to be performed by Alice and Bob in the current entanglement purification protocols when the states $`|0`$ and $`|1`$ are non-degenerate in energy , as required in our protocol. This means ultimately that the existing entanglement purification schemes can only create states of the form $`|\psi ^{}(\eta )`$, where $`\eta `$ is unknown, rather than the true singlets (or states with known $`\eta `$) needed for QCS. We are currently investigating whether we can use such states in a modified version of QCS or indeed whether there are alternative (asynchronous) entanglement purification protocols that can produce pure singlets.
A second limitation of our protocol is the requirement that Alice and Bob be relatively at rest. In a more realistic scenario we would need to assess and correct the effects of relative motions and accelerations, especially on the exact form of the entanglement existing between Alice and Bob.
In conclusion, we have presented a quantum protocol for synchronizing spatially separated atomic-clocks, which uses only shared prior entanglement and a classical channel. The two synchronizing parties may be at far-distant and unknown relative locations and the accuracy of the time synchronisation is not affected by the distance of separation or by noise on the classical channel. Our protocol has direct applications for use in very long baseline interferometry and also provides a means for phase locking remote optical or matter-wave interferometers.
###### Acknowledgements.
We would like to acknowledge valuable discussions with S. L. Braunstein, N. Cerf, J. I. Cirac, H. J. Kimble, N. Linden, H. Mabuchi, L. Maleki, S. Popescu, J. P. Preskill, and N. Yu.
D.S.A., J.P.D. and C.P.W. are supported by a contract with the National Aeronautics and Space Administration. R. J. is supported by the U.K. Engineering and Physical Sciences Research Council.
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# Hydration interactions: aqueous solvent effects in electric double layersAccepted to Phys. Rev. E.
## I Introduction
Electrolytes, in contact with charged surfaces or macro-ions, play an important role in determining the properties of many biological and chemical systems. One of the most widely used tools for studying ions in aqueous solutions is the Poisson-Boltzmann (PB) theory . The mathematical and conceptual simplicity of this theory makes it very appealing both for numerical computation, and for gaining insight into the underlying physical principles. Although the theory contains important simplifications, it has proven to be a useful and accurate tool in the study of systems such as colloidal dispersions, biological membranes, polyelectrolytes and complex systems, e.g., proteins or DNA interacting with charged membranes.
The Poisson-Boltzmann theory is obtained by making two simplifying approximations. The first approximation is the treatment of the electrostatic interactions on a mean-field level. The ions are treated as independent charged particles interacting with an external electrostatic potential, derived self-consistently from the mean charge density distribution. Thus, correlations between the ion positions are not taken into account. The second approximation is the treatment of the ions as point-like objects, interacting only through the electrostatic interaction in a dielectric medium. In reality, ions in aqueous solutions have more intricate interactions . These include a non-Coulombic interaction between ion pairs, which is mainly a short-range steric repulsion, interactions with the polar solvent molecules and short-range interactions with the confining charged surfaces.
Various models have been proposed for the inclusion of effects not accounted for by the PB theory. These include liquid state theory approaches , field theory expansions , computer simulations and other modifications to the PB theory . Most of these models remain within the framework of the so-called “primitive model”, in which the interaction between the ions is modeled as a purely repulsive hard-core interaction. On the other hand, relatively few works have addressed explicitly the discrete nature of the solvent molecules . Clearly, the replacement of the solvent by a continuous medium cannot be precise when the inter-ion distance is comparable to the solvent molecular size. Therefore, when the ions reach high densities the discreteness of the solvent is expected to have an important effect on the ionic distribution. This is of particular importance for water. Due to its high polarity, the strong screening of the electrostatic interaction (represented by the dielectric constant) is modified at small ion separations.
Using the surface force apparatus, it is possible to measure precisely the force between charged mica plates. These measurements supply evidence for the importance of the solvent structure in aqueous solutions . At inter-plate separations below approximately $`20\text{Å}`$ significant deviations are found from the prediction of the Derjaguin-Landau-Verwey-Overbeek (DLVO) theory . The measured force is oscillatory or consists of a series of steps, with a period corresponding to the water molecular size. Oscillatory forces are known to arise as a result of the solvent structuring in layers between surfaces . However, a repulsive contribution is found in addition to the oscillatory force at plate separations below several nanometers . This repulsive force is often referred to as the “hydration force” , and its origin is not yet completely understood .
### A Aqueous pair potential model
Recently , an aqueous pair potential model has been proposed for electrolytes, in which the effect of the solvent on the ions is described as a short-range two-body interaction between the ions. The solvent is replaced by a continuum dielectric medium as in PB theory, but the ions also interact through a two-body short-range hydration interaction . This is shown schematically in Fig. 1.
This aqueous pair potential model , involves several simplifying assumptions. One is that the effect of the solvent can be represented as a linear superposition of two-body potentials between all ion pairs. Another simplification is that the effective potential between the ions is taken as the effective potential in the bulk, regardless of the ion concentration, and of the geometry imposed by the charged surfaces. Finally, a short-range surface-ion effective potential should be included in addition to the ion-ion effective potential. Despite of the simplifications made in the aqueous pair potential model, it offers a first step towards a qualitative understanding of solvent effects on the ion distribution, in particular near highly charged surfaces.
### B Effective ion pair interaction
For the short-range ion-ion interaction, the so-called potential of mean force between ions in solution can be used. Potentials of mean force are defined as $`k_\mathrm{B}T\mathrm{log}g_{ij}(𝐫)`$ where $`g_{ij}(𝐫)`$ are the ion-ion radial distribution functions for ion pairs of species $`i`$ and $`j`$. The radial distribution functions have been calculated numerically for a single ion pair immersed in an aqueous solution using molecular dynamics techniques .
An alternative approach has been proposed in Refs. . In this approach, a Hamiltonian consisting of a pairwise effective potential between the ions is obtained using the so-called “reverse Monte-Carlo” approach. The ion-ion radial distribution functions are first calculated using a molecular dynamics simulation for a system including solvent molecules and a finite concentration of ions. The ion-ion effective potential in the system without the solvent is then adjusted iteratively until the same distribution functions are obtained using Monte-Carlo simulations.
The different available calculations of potentials of mean force differ in their quantitative predictions. This may be a result of high sensitivity of the models to detailed features used for the water molecules and for the inter-molecular interactions . However, all the potentials of mean force as well as the effective potentials are qualitatively similar . Thus, for the purpose of the present work, aiming at a qualitative understanding of solvent effects, any one of these potentials may be used.
At large ionic separations the ion-ion effective potential is well approximated by a screened electrostatic interaction, with the water dielectric constant in the continuum limit. At short ionic separation, the difference between the total effective potential and the screened electrostatic interaction is a short-range potential reflecting the structure of water molecules in the ion vicinity. Fig. 2 shows the short-range contribution (excluding the screened electrostatic part) to the effective potential calculated between $`\mathrm{Na}^+`$ \- $`\mathrm{Na}^+`$ pairs in the reverse Monte-Carlo approach . Below about $`3`$ Å, the electrostatic repulsion between the ions becomes unscreened. Therefore it is much larger than the screened repulsion in the dielectric medium, and the effective potential is strongly repulsive. The unscreened electrostatic potential leads to an effectively enlarged hard-core separation between the ions, relative to a hard-core diameter of about $`2.3\text{Å}`$ used in the short-range part of the bare ion-ion potential. At larger separations, the effective potential is oscillatory, and mainly attractive. It has a distinct minimum at an ion-ion separation of about $`3.6`$ Å, followed by a maximum and a second minimum at approximately $`6`$ Å.
### C The present work
The replacement of the discrete solvent by a continuum medium, with electrostatic and short range interactions between the ions, is a considerable simplification. Still, the statistical mechanical treatment of an electrolyte solution in this model is difficult, and requires the use of further approximations, or simulations.
The Anisotropic Hyper-Netted Chain approximation (AHNC) was previously used to calculate the effects of hydration interactions in the aqueous pair potential model . When the ion concentration is large enough, e.g., near a highly charged surface, the hydration interaction is found to have a significant effect on the distribution of ions in the solution. It was also proposed that the so-called repulsive “hydration forces” between surfaces arise from the ionic structure near highly charged surfaces. According to this description, at large distances from the plate, the ion distribution follows a PB profile with a reduced effective surface charge. When two plates approach each other, the ions near the surfaces come into contact giving rise to an apparent new repulsive force.
In the present work a simple description for ions interacting through electrostatic and short-range attractive interactions as mediated by the solvent molecules is introduced. We apply this description to the aqueous pair potential model. Our aim is limited to describe the important effects of the short-range interaction, and not to provide an accurate tool for their calculation. Therefore our model follows the PB theory as closely as possible, and describes the short-range interaction using a simplified term added to the free energy. The advantage of this approach over more elaborate treatments such as the AHNC , is that it provides relatively simple equations that can be treated numerically and analytically with relative ease as well as allowing extensions to non-planar geometries. In a planar geometry, we show that the effect of the ion-ion hydration interaction can be understood as a perturbation over the PB results. An increase in the concentration of counter-ions near the charged surface is found, and it results in an apparent surface charge which is reduced relatively to the PB theory.
In addition to the calculation of charge distributions, the effect of the hydration interaction on the force between charged particles or surfaces can be studied and will be presented elsewhere . Since we will not discuss inter-surface forces in this paper, it is worthwhile to mention that the results obtained using our model are in good agreement with AHNC calculations and may provide an explanation to the hydration forces as observed in surface force measurements . For high surface charge and plate separations up to approximately $`20\text{Å}`$, important modifications of the PB predictions are found.
The outline of the paper is as follows. Section II presents the model. In section III we apply the model to a single charged plate, present numerical results for the ion density profile and discuss the modifications to the PB theory due to the addition of short-range interactions. In section IV we present analytical results in the low salt limit. We calculate the effective PB surface charge and the effect of the hydration interaction on the density profile of counter-ions in a system with no added salt.
## II The Model
### A Free energy
We start from an approximated free energy, written as a functional of the various ion densities. We choose the electrostatic boundary conditions to be of fixed surface charge densities and write the free energy as a sum of the usual PB term and a correction term, due to hydration, as will be explained below:
$$\mathrm{\Omega }=\mathrm{\Omega }_{\mathrm{PB}}+\mathrm{\Delta }\mathrm{\Omega }$$
(1)
We discuss first how the PB free energy is obtained, and then generalize this result to include the short-range hydration interaction.
#### 1 Poisson-Boltzmann free energy
The Hamiltonian of the system is:
$$H=\frac{1}{2}\underset{i}{}_V\mathrm{d}^3𝐫e_i\rho _i(𝐫)\varphi (𝐫)+\frac{1}{2}_V\mathrm{d}^2𝐫_\mathrm{s}\sigma (𝐫_\mathrm{s})\varphi (𝐫_\mathrm{s})$$
(2)
where $`V`$ is the the volume occupied by the electrolyte solution, $`\sigma (𝐫_\mathrm{s})`$ is the surface charge density of immobile charges on the boundaries $`V`$, and $`e_i`$ is the charge of the $`i`$th ion species. The ion densities $`\rho _i(𝐫)`$ are:
$$\rho _i(𝐫)\underset{j=1}{\overset{N_i}{}}\delta (𝐫𝐫_i^j)$$
(3)
where $`𝐫_i^j`$ is the position of the $`j`$th ion of the $`i`$th species, and the electrostatic potential $`\varphi (𝐫)`$ is a function the different ion positions:
$$\varphi (𝐫)=\underset{j}{}_V\mathrm{d}^3𝐫^{}\frac{e_j\rho _j(𝐫^{})}{\epsilon \left|𝐫𝐫^{}\right|}+_V\mathrm{d}^2𝐫_\mathrm{s}\frac{\sigma (𝐫_\mathrm{s})}{\epsilon \left|𝐫𝐫_\mathrm{s}\right|}$$
(4)
where $`\epsilon =78`$ is the dielectric constant of water. The PB theory is obtained by using a mean-field approximation for the electrostatic interaction. The Hamiltonian (2) is first replaced by a mean-field Hamiltonian, where the electrostatic potential $`\varphi (𝐫)`$ is replaced by an external field $`\mathrm{\Psi }(𝐫)`$. In the thermodynamic limit the free energy can then be written as a functional of the mean densities of the ion species, $`c_i(𝐫)=\rho _i(𝐫)_{\mathrm{MF}}`$, as follows:
$`\mathrm{\Omega }_{\mathrm{MF}}`$ $`=`$ $`k_BT{\displaystyle _V}{\displaystyle \underset{i}{}}c_i\left[\mathrm{log}{\displaystyle \frac{c_i}{\zeta _i}}1\right]\mathrm{d}^3𝐫`$ (7)
$`+{\displaystyle \frac{1}{2}}{\displaystyle _V}{\displaystyle \underset{i}{}}e_i\mathrm{\Psi }(𝐫)c_i(𝐫)\mathrm{d}^3𝐫`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle _V}\sigma (𝐫_s)\mathrm{\Psi }(𝐫_s)`$
where $`k_\mathrm{B}T`$ is the thermal energy, and $`\zeta _i`$ is the fugacity of the $`i`$th ion species. The mean-field approximation is obtained by requiring that the external potential $`\mathrm{\Psi }(𝐫)`$ is the thermodynamical average of (4) in the system with the mean-field Hamiltonian, $`\mathrm{\Psi }(𝐫)=\varphi (𝐫)_{\mathrm{MF}}`$, i.e.:
$$\mathrm{\Psi }(𝐫)=\underset{i}{}_V\mathrm{d}^3𝐫^{}\frac{e_ic_i(𝐫^{})}{\epsilon \left|𝐫𝐫^{}\right|}+_V\mathrm{d}^2𝐫_\mathrm{s}\frac{\sigma (𝐫_\mathrm{s})}{\epsilon \left|𝐫𝐫_\mathrm{s}\right|}$$
(8)
This relation is equivalent to the Poisson equation:
$$^2\mathrm{\Psi }=\frac{4\pi }{\epsilon }\underset{i}{}e_ic_i$$
(9)
supplemented by the boundary condition:
$$\mathrm{\Psi }\widehat{𝐧}|_{𝐫_𝐬}=\frac{4\pi }{\epsilon }\sigma (𝐫_𝐬)\text{on the charged surfaces}$$
(10)
where the normal vector $`\widehat{𝐧}`$ points away from the charged surfaces into the volume occupied by the ionic solution. Using this boundary condition and Eq. (9), the second and third terms of Eq. (7) can be re-expressed as:
$`{\displaystyle \frac{1}{2}}{\displaystyle _V}{\displaystyle \underset{i}{}}e_i\mathrm{\Psi }(𝐫)c_i(𝐫)\mathrm{d}^3𝐫+{\displaystyle \frac{1}{2}}{\displaystyle _V}\sigma (𝐫_s)\mathrm{\Psi }(𝐫)`$ (11)
$`={\displaystyle \frac{\epsilon }{8\pi }}{\displaystyle _V}(\mathrm{\Psi })^2\mathrm{d}^3𝐫`$ (12)
Substituting this relation in Eq. (7) we obtain the PB free energy:
$`\mathrm{\Omega }_{\mathrm{PB}}`$ $`=`$ $`{\displaystyle \frac{\epsilon }{8\pi }}{\displaystyle (\mathrm{\Psi })^2\mathrm{d}^3𝐫}`$ (13)
$`+`$ $`k_BT{\displaystyle \underset{i}{}c_i\left[\mathrm{log}\frac{c_i}{\zeta _i}1\right]\mathrm{d}^3𝐫}`$ (14)
$`+`$ $`{\displaystyle \lambda (𝐫)\left(^2\mathrm{\Psi }+\frac{4\pi }{\epsilon }\underset{i}{}e_ic_i\right)\mathrm{d}^3𝐫}`$ (15)
The first term in $`\mathrm{\Omega }_{\mathrm{PB}}`$ is the electrostatic free energy and the second term is the entropy of the ions. The fugacity $`\zeta _i`$, in the second term, is equal in PB theory to the bulk concentration $`c_{\mathrm{b},i}`$ of the $`i`$th ion species, $`\zeta _i=c_{\mathrm{b},i}`$, as for an ideal gas. For more generalized free energies, a different relation may exist between the fugacity of each ion species and its respective bulk concentration. The electrostatic potential $`\mathrm{\Psi }`$ is a functional of the ion densities $`c_i`$, and is determined by the Poisson equation (9) and the boundary conditions (10) imposed by the surface charges. Alternatively, in Eq. (15) $`\mathrm{\Psi }`$ is regarded as an independent field and a third term containing a Lagrange multiplier $`\lambda (𝐫)`$ is added to $`\mathrm{\Omega }_{\mathrm{PB}}`$. The PB equilibrium mean densities $`c_i(𝐫)`$ result from minimizing $`\mathrm{\Omega }_{\mathrm{PB}}`$. With the introduction of $`\lambda (𝐫)`$ the minimization is equivalent to requiring an extremum of $`\mathrm{\Omega }_{\mathrm{PB}}`$ with respect to the three fields $`c_i`$, $`\mathrm{\Psi }`$ and $`\lambda `$, subject to the boundary condition (10). By requiring first an extremum of $`\mathrm{\Omega }_{\mathrm{PB}}`$ with respect to $`\mathrm{\Psi }`$ and $`c_i`$ the following relations are obtained:
$$\lambda =\frac{\epsilon }{4\pi }\mathrm{\Psi }$$
(16)
and:
$$c_i=\zeta _i\mathrm{exp}\left(\beta e_i\mathrm{\Psi }\right)$$
(17)
Where $`\beta =1/(k_\mathrm{B}T)`$ and $`\zeta _i=c_{\mathrm{b},i}`$. The extremum condition with respect to $`\lambda `$ gives the Poisson equation:
$$^2\mathrm{\Psi }=\frac{4\pi }{\epsilon }\underset{i}{}e_ic_i$$
(18)
Combining these relations we obtain the PB equation:
$$^2\mathrm{\Psi }=\frac{4\pi }{\epsilon }\underset{i}{}\zeta _ie_i\mathrm{e}^{\beta e_i\mathrm{\Psi }}$$
(19)
Alternatively, the first two relations, obtained from the extremum condition with respect to $`\mathrm{\Psi }`$ and $`c_\mathrm{i}`$, can be substituted into Eq. (15). Formally, this gives $`\mathrm{\Omega }_{\mathrm{PB}}`$ as a functional of $`\lambda `$. Using Eq. (16), the expression obtained for $`\mathrm{\Omega }_{\mathrm{PB}}`$ can be written as a functional of $`\mathrm{\Psi }`$:
$`\mathrm{\Omega }_{\mathrm{PB}}`$ $`=`$ $`{\displaystyle \frac{\epsilon }{8\pi }}{\displaystyle _V}(\mathrm{\Psi })^2\mathrm{d}^3𝐫+{\displaystyle _V}\sigma \mathrm{\Psi }\mathrm{d}^2𝐫_\mathrm{s}`$ (21)
$`k_\mathrm{B}T{\displaystyle _V}\mathrm{d}^3𝐫{\displaystyle \underset{i}{}}\zeta _i\mathrm{e}^{\beta e_i\mathrm{\Psi }}`$
where the second integration is over the charged surfaces. Requiring an extremum of this functional with respect to $`\mathrm{\Psi }`$ is another way to obtain the Poisson-Boltzmann equation (19).
A more formal derivation of the mean field, PB free energy, and a discussion on its generalization to systems with non-electrostatic interactions is presented in Ref. . The PB free energy (15) can also be derived by formulating the problem using field theory methods. In this approach the mean-field approximation is obtained as the saddle point of the functional integral, and corrections due to ion-ion correlations can be obtained in a systematic expansion .
#### 2 Inclusion of the hydration interaction
As discussed in the introduction, our starting point is a model in which the hydration interaction, arising from solvent effects, is described as an effective ion-pair interaction. We denote this short-range potential between ions of species $`i`$ and $`j`$ at distance r as $`u_{ij}(𝐫)`$. The potential is taken as the short-range effective potential between ions immersed in a bulk ionic solution having a specific, constant concentration. Therefore, $`u_{ij}(𝐫)`$ is assumed to be isotropic and does not depend on the ion positions or the confining geometry.
Our aim is to treat the long range electrostatic interaction on the mean-field level, as in PB theory. Thus, we begin by considering the free energy of a system placed in some arbitrary field $`\mathrm{\Psi }(𝐫)`$, where the ions interact with each other only through the two-body potential $`u_{ij}(𝐫)`$. Due to the short-range nature of the hydration interaction, the free energy can be obtained from a virial expansion of the grand canonical partition function. Since we will be interested in highly inhomogeneous systems, we perform an expansion in the inhomogeneous ion density. The derivation is given in Appendix A. Including terms up to the quadratic order in the expansion we obtain:
$`\mathrm{\Omega }_\mathrm{h}`$ $`=`$ $`k_\mathrm{B}T{\displaystyle \underset{i}{}c_i\left[\mathrm{log}\frac{c_i}{\zeta _i}1\right]\mathrm{d}^3𝐫}`$ (24)
$`+{\displaystyle \underset{i}{}e_ic_i\mathrm{\Psi }\mathrm{d}^3𝐫}`$
$`+{\displaystyle \frac{k_BT}{2}}{\displaystyle \underset{i,j}{}}{\displaystyle c_i(𝐫)U_{ij}(𝐫𝐫^{})c_j(𝐫^{})\mathrm{d}^3𝐫\mathrm{d}^3𝐫^{}}`$
where $`\mathrm{\Psi }(𝐫)`$ is an external field, coupled to the $`i`$th ion charge density $`e_ic_i`$. The short-range weighted potential $`U_{ij}`$ in the third term of $`\mathrm{\Omega }_\mathrm{h}`$ is defined as:
$$U_{ij}=1\mathrm{e}^{\beta u_{ij}(|𝐫𝐫^{}|)}$$
(25)
where $`u_{ij}`$ is the nominal short-range interaction potential between ions of species $`i`$ and $`j`$. This form of describing the short-range interaction is a rather crude approximation, valid only in the low density limit. Its advantage is its simplicity. The free energy $`\mathrm{\Omega }_\mathrm{h}`$ amounts to setting the direct correlation function $`c_2(|𝐫𝐫^{}|)`$ to be equal to $`U(|𝐫𝐫^{}|)`$, and all higher order direct correlation functions to zero .
Having found the hydration free energy $`\mathrm{\Omega }_\mathrm{h}`$, the electrostatic interaction can be treated on the mean-field level. This is done by considering $`\mathrm{\Psi }(𝐫)`$ as the electrostatic potential and imposing the self-consistency requirement of the Poisson equation (9). This is essentially the approximation we used to derive the PB equation (19), with the difference that the free energy of a dilute, non-interacting ion distribution is replaced by the free energy $`\mathrm{\Omega }_\mathrm{h}`$ of Eq. (24). The result is the free energy of Eq. (1), with $`\mathrm{\Delta }\mathrm{\Omega }`$ defined as follows:
$$\mathrm{\Delta }\mathrm{\Omega }=\frac{k_BT}{2}\underset{i,j}{}c_i(𝐫)U_{ij}(𝐫𝐫^{})c_j(𝐫^{})\mathrm{d}^3𝐫\mathrm{d}^3𝐫^{}$$
(26)
We conclude this section with some remarks on the approach presented above. Important solvent effects are already introduced in the PB theory by using an electrostatic interaction with a dielectric constant $`\epsilon =78`$ of water, instead of the bare electrostatic interaction. In the modified model a more precise effective potential between the ions is used. The separation of this potential into a long-range electrostatic term and a short-range hydration term allows each of these two interactions to be treated in a simple though approximated form. The virial expansion is a standard choice for approximating short-range interactions. Such an expansion fails for the electrostatic interaction due to its long-range . On the other hand, the wide success of PB theory demonstrates that the electrostatic interaction can be treated quite well in the mean-field approximation. Therefore we use this approximation for the long-range part of the interaction, and in this respect we remain within the framework of PB theory.
The free energy (1) can also be obtained by rewriting the grand canonical partition function as a field-theory partition function. The short-range part of the interaction can be separated from the electrostatic interaction and a different expansion can be performed for each of these interactions. By using a density expansion for the short-range interaction and a loop expansion for the electrostatic interaction, Eq. (1) is obtained up to second order in the density expansion and first order in the electrostatic potential .
The simplicity of our approach can lead to elegant analytical results, but has several limitations. The use of only the second term in the virial expansion implies that we are using a low density approximation. The validity of such an approximation for a bulk fluid can be assessed by considering $`k_\mathrm{B}TB_2c`$, where $`B_2`$ is the second virial coefficient in the expansion of the pressure, and $`c`$ is the ion density. Qualitatively, if $`B_2c`$ is small compared to unity, the correction to the ideal gas behavior is small and truncating the virial expansion after the second term is sensible. For non-homogeneous cases, the corresponding quantity is $`(1/2)_jd𝐫^{}c(𝐫^{})U_{ij}(𝐫𝐫^{})`$. For relatively high surface charges considered here this integral approaches values of order unity near the charged surfaces, indicating that the approximation should only be expected to give qualitative results. Another deficiency of the virial expansion to second order can be seen from the fact that the direct correlation function is simply $`U_{ij}(𝐫)`$. This implies that the hard core interaction is not described accurately in our treatment. A faithful description would require the vanishing of the pair correlation function $`h_2(𝐫)`$ for separations smaller than the hard-core diameter. Hence, it should be kept in mind that our main concern is to study the effects of a short-range interaction with a dominant attractive part. Finally, the fact that we describe the electrostatic interaction in the mean-field approximation implies that ion-ion correlations are ignored, as they are in PB theory. When our approach is applied for the aqueous pair potential model, these approximations should also be kept in mind. In particular, we follow Ref. and do not include an effective ion-surface potential .
### B Density equations
The mean density distribution is obtained by minimizing the total free energy $`\mathrm{\Omega }=\mathrm{\Omega }_{\mathrm{PB}}+\mathrm{\Delta }\mathrm{\Omega }`$. From equations (1), (15) and (26) we have:
$`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{\epsilon }{8\pi }}{\displaystyle (\mathrm{\Psi })^2\mathrm{d}^3𝐫}+k_BT{\displaystyle \underset{i}{}c_i\left(\mathrm{log}\frac{c_i}{\zeta _i}1\right)\mathrm{d}^3𝐫}`$ (29)
$`+{\displaystyle \frac{k_BT}{2}}{\displaystyle \underset{i,j}{}}{\displaystyle c_i(𝐫)U_{ij}(𝐫𝐫^{})c_j(𝐫^{})\mathrm{d}^3𝐫\mathrm{d}^3𝐫^{}}`$
$`+{\displaystyle \lambda (𝐫)\left(^2\mathrm{\Psi }+\frac{4\pi }{\epsilon }\underset{i}{}c_ie_i\right)\mathrm{d}^3𝐫}`$
where $`\mu _i`$ and $`\zeta _i=\mathrm{exp}(\beta \mu _i)/\lambda _\mathrm{T}^3`$ are the chemical potential and the fugacity of the ion species $`i`$, respectively. The thermal de Broglie wavelength, $`\lambda _\mathrm{T}`$, is equal to $`h/(2\pi mk_\mathrm{B}T)^{1/2}`$, where $`h`$ is the Planck constant and $`m`$ is the ion mass. Requiring an extremum of $`\mathrm{\Omega }`$ with respect to $`\mathrm{\Psi }`$ gives: $`\lambda =(\epsilon /4\pi )\mathrm{\Psi }`$ as in Eq. (16). Taking the variation with respect to $`c_i`$ then gives:
$$\mathrm{log}\frac{c_i(𝐫)}{\zeta _i}+\underset{j}{}c_j(𝐫^{})U_{ij}(𝐫𝐫^{})\mathrm{d}^3𝐫^{}+\beta e_i\mathrm{\Psi }(𝐫)=0$$
(30)
This equation is supplemented by the Poisson equation (9). Since Eq. (30) is an integral equation, the $`c_i`$ cannot be written as a simple function of $`\mathrm{\Psi }`$ as in the PB case. Therefore, a single equation for $`\mathrm{\Psi }`$, analogous to the PB equation, cannot be obtained, and we are left with the two coupled integro-differential equations (30) and (9). These equations should be solved together to obtain the electrostatic potential and density profiles. In the case $`U0`$, Eq. (30) reduces to the Boltzmann relation $`c_i=\zeta _i\mathrm{exp}(\beta e_i\mathrm{\Psi })`$ with $`\zeta _i=c_{\mathrm{b},i}`$. Combining this relation with Eq. (9) reproduces the PB equation (19).
In order to simplify the set of equations, we assume the same short-range interaction between the different pairs of ion species. Assuming that the charged surfaces are negatively charged, we choose: $`u_{ij}(𝐫)=u_{++}(𝐫)u(𝐫)`$, where $`u_{++}(𝐫)`$ is the short-range effective potential between the (positive) counter-ions. This assumption is not exact for the effective potentials of ions in water . However, since only the counter-ions reach high densities, close to the oppositely charged surfaces, and the co-ions are repelled from the surface neighborhood, the exact choice of the potentials $`u_+(𝐫)`$ and $`u_{}(𝐫)`$ is expected to be of only minor significance.
We now consider an electrolyte of valency $`z_+`$:$`z_{}`$, i.e., a solution of positive and negative ions of charges $`e_\pm =\pm z_\pm e`$, where $`e`$ is the electron charge. We designate the surface charge density on the plate as a constant $`\sigma `$ and the bulk densities of the positive and negative ions as $`c_\mathrm{b}c_{\mathrm{b},+}`$ and $`c_{\mathrm{b},}`$, respectively. Due to charge neutrality in the bulk, $`c_{\mathrm{b},}=(z_+/z_{})c_\mathrm{b}`$ and similarly, $`\zeta _{}=(z_+/z_{})\zeta `$ where $`\zeta \zeta _+`$. Equation (30) can then be written as follows:
$$c_\pm (𝐫)=\zeta _\pm e^{\beta ez_\pm \mathrm{\Psi }}\mathrm{exp}\left[c(𝐫^{})U(𝐫𝐫^{})\mathrm{d}^2𝐫^{}\right]$$
(31)
where $`c(𝐫)=c_+(𝐫)+c_{}(𝐫)`$ is the total ion density, and $`U(𝐫)=U_{++}(𝐫)`$ is obtained from $`u(𝐫)`$ using Eq. (25). From the Poisson equation (9) we obtain:
$`^2\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \frac{4\pi e}{\epsilon }}\left(z_+c_+z_{}c_{}\right)`$ (32)
$`=`$ $`{\displaystyle \frac{4\pi e}{\epsilon }}\zeta z_+\left(\mathrm{e}^{\beta ez_{}\mathrm{\Psi }}\mathrm{e}^{\beta ez_+\mathrm{\Psi }}\right)`$ (34)
$`\times \mathrm{exp}\left[{\displaystyle c(𝐫^{})U(𝐫𝐫^{})\mathrm{d}^3𝐫^{}}\right]`$
Note that in addition to the explicit dependence on the ion valencies $`z_\pm `$ in equations (31) and (34), in a more realistic model the details of the potential $`u(𝐫)`$ should also depend on the type of counter-ion species present in the problem.
## III Single charged plate
### A Density equations
After presenting the general formalism let us consider, as an example, a single negatively charged planar surface (Fig. 1). The charged surface is in contact with an electrolyte of valency $`z_+`$:$`z_{}`$. We designate the axis perpendicular to the plate as the $`z`$ axis, and consider the ion solution in the region $`z>0`$. For simplicity we consider positive and negative ions of the same hard-core diameter $`d_{\mathrm{hc}}`$. The coordinate of closest approach of the ions to the plate is designated as $`z=0`$. Hence the “real” surface lies at a distance of one ion radius $`d_{\mathrm{hc}}/2`$ from the actual $`z=0`$ plate position, as shown in Fig. 1a. When we refer to conventional PB results, however, the ions are point-like and the plate should be understood to be positioned exactly at $`z=0`$.
Due to the one-dimensional symmetry imposed by the uniformly charged planar plate, the integration in Eq. (31) can be performed over the $`xy`$ plane, leaving us with profiles depending only on $`z`$, the distance from the plate:
$$c_\pm (z)=\zeta _\pm e^{\beta ez_\pm \mathrm{\Psi }}\mathrm{exp}\left[_0^{\mathrm{}}c(z^{})B(zz^{})dz^{}\right]$$
(35)
where $`c(z)=c_+(z)+c_{}(z)`$ is the total ion density and $`B(z)`$ is the integral of $`U(𝐫)`$ in the plane of constant $`z`$. Using cylindrical coordinates:
$$B(z)=2\pi _0^{\mathrm{}}\rho d\rho U\left(\sqrt{z^2+\rho ^2}\right)$$
(36)
and the Poisson equation (34) reads:
$`{\displaystyle \frac{\mathrm{d}^2\mathrm{\Psi }}{\mathrm{d}z^2}}`$ $`=`$ $`{\displaystyle \frac{4\pi e}{\epsilon }}\zeta z_+\left(\mathrm{e}^{\beta ez_{}\mathrm{\Psi }}\mathrm{e}^{\beta ez_+\mathrm{\Psi }}\right)`$ (38)
$`\times \mathrm{exp}\left[{\displaystyle _0^{\mathrm{}}}c(z^{})B(zz^{})dz^{}\right]`$
Equations (35) and (38) are supplemented by the boundary conditions:
$$\frac{\mathrm{d}\mathrm{\Psi }}{\mathrm{d}z}|_{z=0}=\frac{4\pi }{\epsilon }\sigma ;\frac{\mathrm{d}\mathrm{\Psi }}{\mathrm{d}z}|_z\mathrm{}=0$$
(39)
Finally, the relation between $`\zeta `$ and the bulk density $`c_\mathrm{b}`$ can be obtained from Eq. (35). As $`z\mathrm{}`$, $`\mathrm{\Psi }`$ becomes zero, and $`c_\pm `$ assume their asymptotic constant, bulk values. Thus the integrand inside the exponential can be replaced by $`(1+z_+/z_{})c_\mathrm{b}B(zz^{})`$. Recalling that $`c_+=c_\mathrm{b}`$ and $`c_{}=(z_+/z_{})c_\mathrm{b}`$, we obtain:
$$c_\mathrm{b}=\zeta \mathrm{exp}\left[\left(1+\frac{z_+}{z_{}}\right)B_\mathrm{t}c_\mathrm{b}\right]$$
(40)
where:
$$B_\mathrm{t}_{\mathrm{}}^{\mathrm{}}dzB(z)=\mathrm{d}^3𝐫U(𝐫)$$
(41)
is also equal to $`2B_2`$, the second virial coefficient. Note that $`B(z)`$ and $`B_\mathrm{t}`$ are negative for an attractive interaction. The limit $`B_\mathrm{t}c_\mathrm{b}0`$ is the limit in which the short-range interaction becomes negligible in the bulk. In this limit the relation between the bulk density and fugacity of Eq. (40) tends to the ideal gas relation $`c_\mathrm{b}=\zeta =\mathrm{exp}(\beta \mu )/\lambda _\mathrm{T}^3`$.
Two special cases will be of particular interest in the following sections. The first is the case of a monovalent 1:1 electrolyte, where we have:
$`c_\pm (z)`$ $`=`$ $`\zeta \mathrm{e}^{\beta e\mathrm{\Psi }}\mathrm{exp}\left[{\displaystyle _0^{\mathrm{}}}c(z^{})B(zz^{})dz^{}\right]`$ (42)
$`{\displaystyle \frac{\mathrm{d}^2\mathrm{\Psi }}{\mathrm{d}z^2}}`$ $`=`$ $`={\displaystyle \frac{4\pi }{\epsilon }}c(z)`$ (43)
and:
$$c_\mathrm{b}=\zeta \mathrm{exp}\left(2B_\mathrm{t}c_\mathrm{b}\right)$$
(44)
The second case is that of no added salt. The solution contains only monovalent counter-ions ($`z_+=1`$, $`z_{}=0`$). This case can be obtained by taking formally the limit $`\zeta 0`$ of Eq. (43), or by repeating the derivation from Eq. (29) with only one type of ions, of charge $`e`$. The term $`k_\mathrm{B}T\mathrm{d}^3\mathrm{r}c\mathrm{log}(\zeta )`$ in $`\mathrm{\Omega }`$ is then a Lagrange multiplier added to impose the condition: $`_0^{\mathrm{}}dzec(z)=|\sigma |`$. The following equations are then obtained:
$`c(z)`$ $`=`$ $`\zeta _0\mathrm{e}^{\beta e\mathrm{\Psi }}\mathrm{exp}\left[{\displaystyle _0^{\mathrm{}}}c(z^{})B(zz^{})dz^{}\right]`$ (45)
$`{\displaystyle \frac{\mathrm{d}^2\mathrm{\Psi }}{\mathrm{d}z^2}}`$ $`=`$ $`{\displaystyle \frac{4\pi e}{\epsilon }}c(z)`$ (46)
where $`\zeta _0`$ is an arbitrary reference fugacity. The choice of $`\zeta _0`$ determines the (arbitrary) position in which $`\mathrm{\Psi }`$ is zero. Note that the electrostatic potential $`\mathrm{\Psi }`$ diverges in the bulk. This divergence exists also in the usual PB theory, because the system is effectively one dimensional with no screening by added salt. Although $`\mathrm{\Psi }(z)`$ has a weak logarithmic divergence, the density of counter-ions decays to zero, $`lim_z\mathrm{}c(z)=0`$ as it should.
### B Parameters and length scales
For the ion-ion potential $`u(𝐫𝐫^{})`$ we use an effective potential between $`\mathrm{Na}^+`$ \- $`\mathrm{Na}^+`$ ion pairs. The potential was calculated using a Monte-Carlo simulation for an NaCl ionic solution of concentration $`0.55`$ M, at room temperature. The electrostatic interaction between the ions is subtracted, and the net short-range potential is shown in Fig. 2. For ion-ion separations below $`2.9`$ Å a hard core interaction is assumed. Fig. 3 shows the function $`B(z)`$, derived from this potential, using Eq. (36). Note that $`B(z)`$ has several local maxima and minima. These correspond to the local maxima and minima of $`u(𝐫)`$. Thus the structure of $`B(z)`$ reflects the oscillatory behavior of the effective potential.
We use the effective potential calculated for $`c_\mathrm{b}=0.55`$ M, regardless of the actual bulk ion concentration in the system. Since the important effects occur near the charged surface, where the ion concentration is much larger than $`c_\mathrm{b}`$, it seems reasonable to use an effective potential calculated in the presence of a rather high salt concentration. The choice of $`c_\mathrm{b}=0.55`$ M is still somewhat arbitrary, and we rely on the fact that the dependence of the effective potential on the ion concentration is weak .
It is useful to employ two length scales that characterize the PB density profiles . The Gouy-Chapman length, defined as $`b=\epsilon k_\mathrm{B}T/(2\pi e|\sigma |`$), characterizes the width of the diffusive counter-ion layer close to a single plate charged with a surface charge $`\sigma `$, in the absence of added salt. The Debye-Hückel screening length, $`\lambda _\mathrm{D}=(8\pi c_\mathrm{b}e^2/\epsilon k_\mathrm{B}T)^{1/2}`$, equal to $`19.6\text{Å}`$ for $`c_\mathrm{b}=0.025`$M at room temperature characterizes the decay of the screened electrostatic interaction in a solution with added salt. The strength of the electrostatic interaction can also be expressed using the Bjerrum length, $`l_\mathrm{B}=e^2/(\epsilon k_\mathrm{B}T)`$. This is the distance at which the electrostatic interaction between two unit charges becomes equal to the thermal energy. The Bjerrum length is equal to about $`7`$ Å in water at room temperature.
The inclusion of hydration interactions introduces additional length scales in the system. For the interaction shown in Figs. 2 and 3, the range of the interaction $`d_{\mathrm{hyd}}`$ can be seen to be approximately $`7`$ Å, over twice the hard core diameter $`d_{\mathrm{hc}}=2.9`$ Å. The strength of the hydration interaction is characterized by $`B_\mathrm{t}(7.9\text{Å})^3`$, as is calculated from Eq. (41).
### C Numerical results
Equations (35) and (38) are a set of three nonlinear integro-differential equations. We treat them numerically using an iterative scheme, based on the assumption that the positive ion density profile is dominated by the electrostatic interaction. We start with the analytically known PB profile close to a single charged plate and calculate iteratively corrections to this profile, as result from equations (35) and (38). For a 1:1 electrolyte we iteratively solve the equation:
$`{\displaystyle \frac{\mathrm{d}^2\mathrm{\Psi }^{(n)}}{\mathrm{d}z^2}}`$ $`=`$ $`{\displaystyle \frac{8\pi e}{\epsilon }}\zeta \mathrm{sinh}\left(\beta e\mathrm{\Psi }^{(n)}\right)`$ (48)
$`\times \mathrm{exp}\left[{\displaystyle _0^{\mathrm{}}}c^{(n1)}(z^{})B(zz^{})dz^{}\right]`$
where $`c(z)=c_+(z)+c_{}(z)`$ is the total ion density and the superscript $`n`$ stands for the $`n`$th iteration. For $`n>0`$:
$`c_\pm ^{(n)}(z)`$ $``$ $`\zeta \mathrm{e}^{\beta e\mathrm{\Psi }^{(n)}}`$ (50)
$`\times \mathrm{exp}\left[{\displaystyle _0^{\mathrm{}}}c^{(n1)}(z^{})B(zz^{})dz^{}\right]`$
and the zeroth order densities $`c_\pm ^{(0)}`$ are taken as the density profiles generated by the PB equation (19). The boundary conditions (39) are satisfied by the electrostatic potential $`\mathrm{\Psi }^{(n)}`$ in all the iterations. Note that using our iterative scheme, Eq. (48) is an inhomogeneous differential equation, because the integral in the exponential is a known function of $`z`$, calculated numerically in the $`(n1)`$ iteration. A similar iterative scheme, based on Eq. (46) can be used when only counter-ions are present in the solution.
Figure 4 shows the calculated density profile of the counter-ions on a semi-logarithmic scale, for a charged plate with a surface charge, $`|\sigma |=0.333\mathrm{C}/\mathrm{m}^2`$, corresponding to an area of approximately $`48\text{Å}^2`$ per unit charge. This is a typical high surface charge obtained with mica plates. It corresponds to a Gouy-Chapman length $`b=1.06`$ Å, at a temperature of $`298`$K, with $`\epsilon =78`$. No salt is present in the solution. The calculated density profile (solid line) is compared to the PB prediction (dotted line). The short-range attraction favors an increased concentration of counter-ions in the vicinity of the charged plate. This results in an increase of the concentration relative to the PB prediction. For a surface charge as in Fig. 4, an increase of the concentration is seen at distances from the plate up to approximately $`4.5\text{Å}`$. The overall number of counter-ions is fixed by the requirement of charge neutrality. Therefore, the increase in the density of counter-ions near the plate is balanced by a reduced concentration further away.
When salt is present in the solution, the short-range attraction draws additional ions from the bulk solution to the diffuse electrical layer near the plate. This can be seen in Fig. 5, in a comparison of counter-ion profiles for different values of the bulk concentration $`c_\mathrm{b}`$. For each salt concentration, the figure shows the ratio between the counter-ion density and the density predicted by PB theory, as a function of the distance from the plate. The dotted line shows the result in the no-salt limit. As the salt concentration increases, the counterion concentration increases relative to the PB concentration at all distances from the charged plate. Qualitatively, however, the hydration effect on the counter-ion profile is similar in all the curves. As long as the Debye-Hückel screening length is large compared to the Gouy-Chapman length, $`b=1.06\text{Å}`$, the density profile in the vicinity of the plate is dominated by the balancing counter-ions and the salt has only a small effect.
The effect of the hydration interaction is strongly dependent on the surface charge $`\sigma `$. As $`\sigma `$ is increased, the ion density near the surface increases too. The exponential in Eq. (46) deviates more strongly from unity, leading to a larger deviation from PB theory. The dependence on $`\sigma `$ is demonstrated in Fig. 6. The ratio of the positive ion density to its PB value is shown for three values of the surface charge. The effect of the hydration potential is very minor for small surface charge, $`|\sigma |=0.0333\mathrm{C}/\mathrm{m}^2`$ (dotted line), where the deviation from PB is less than $`2\%`$ at its maximum, and considerable for a surface charge of $`0.333\mathrm{C}/\mathrm{m}^2`$ (dashed line), where the deviation from PB reaches almost $`40\%`$.
The numerical scheme, described above, requires several iterations to converge fully. It is interesting to note, though, that the first iteration captures most of the effect of the short range interaction. This indicates that the density profile is dominated, as we assumed, by the electrostatic interaction, and assures that the convergence of the iterative scheme is good with the PB density profile as the zero-th order approximation. On the theoretical level it indicates that the effect of the hydration interaction can be seen as a perturbation over the PB results. The fact that the first iteration provides a good approximation to the full iterative result can lead to further analytical approximations. For example, the corrections to the density profiles, in the no added salt limit, are studied analytically in the next section, based on this observation.
As an example for the results of the first iteration, we compare, in Fig. 7, the correction to the counter-ion density profile obtained in the first iteration (dashed line), with the full iterative result (solid line). We use a high surface charge of $`0.333\mathrm{C}/\mathrm{m}^2`$, where the differences between the exact profile and that of the first iteration are relatively pronounced. The two density profiles differ by at most $`3.2`$ percent, where the ion density deviates from the PB value by $`30`$ percent. For smaller surface charge the results obtained in the first iteration are even better.
### D Contact density and the contact theorem
The contact density of the ions is barely modified as compared with the PB prediction. This is evident in Figs. 4-6. As long as the Debye-Hückel screening length is large compared to the Gouy-Chapman length, or the hydration interaction is negligible in the bulk, the modification remains small. This result can be obtained from a generalization of the PB contact theorem :
$$\underset{i}{}c_i(0)\frac{2\pi \beta }{\epsilon }\sigma ^2=P_{\mathrm{bulk}}$$
(51)
where $`P_{\mathrm{bulk}}`$ is the bulk pressure of the ionic solution. Equation (51) is derived in detail for the free energy used in our model in Ref. . It is obtained from the equality of the internal pressure in the electrolyte solution at different distances from the charged plate. Far away from the charged plate the pressure must be equal to the bulk pressure of the ionic solution, because the densities approach their bulk values and the electrostatic potential becomes constant. At the contact plane between the plate and the solution, the pressure involves only an electrostatic contribution and an osmotic contribution, as in PB theory. This is due to the fact that in our model no short range interaction between the plate and the ions is included. Equating the pressure at the contact plane and far away from the plate results in Eq. (51).
The contact density, as expressed by Eq. (51), differs from the PB prediction only due to the change in the actual value of $`P_{\mathrm{bulk}}`$. This change is negligible if the short-range interaction is not of importance in the bulk. In addition, if the surface charge is high, such that $`b\lambda _\mathrm{D}`$, $`P_{\mathrm{bulk}}`$ is negligible compared to the second term in the left hand side of Eq. (51). Thus the contact density remains very close to the PB prediction. In the no-salt limit $`P_{\mathrm{bulk}}`$ is zero and the contact density coincides exactly with the PB result, $`c_+(0)=(2\pi \beta /\epsilon )\sigma ^2`$.
## IV Analytical solutions
The simplicity of the model makes it possible to obtain various analytical results. The effect of the hydration on the ion distribution can be characterized by several quantities, such as the magnitude of the deviation from the PB result and the effective PB surface charge density seen at a distance from the plate. Using several simplifying assumptions it is possible to obtain analytical expressions for these quantities.
First we assume that the hydration interactions can be neglected in the bulk, i.e., $`B_\mathrm{t}c_\mathrm{b}1`$. In this case, the effect of the hydration potential is significant only in the vicinity of the charged surface, where the ion density becomes large. In addition, the Debye-Hückel screening length, $`\lambda _\mathrm{D}`$, is taken to be large compared to the Gouy-Chapman length $`b=e/(2\pi l_\mathrm{B}|\sigma |)`$. Since $`\lambda _\mathrm{D}b`$, the negative co-ion density near the negatively charged surface can be neglected compared to the positive counter-ion density. Far away from the charged plate, the system is well described using the PB equation, with an effective surface charge density $`\sigma _{\mathrm{eff}}`$ different from the actual charge density $`\sigma `$. The result of the above two simplifying assumptions is that the salt is of minor importance in the region where the effective surface charge is determined. The effective surface charge can then be inferred by considering the case in which only counter-ions are present in the solution (no added salt).
Equation (46) can now be recast in a simpler form, by considering $`\eta \mathrm{log}(c/\zeta _0)`$, as expressed by Eq. (46), and taking its second derivative:
$$\frac{\mathrm{d}^2\eta }{\mathrm{d}z^2}=\frac{4\pi }{\epsilon }\beta e^2\zeta _0\mathrm{e}^\eta _0^{\mathrm{}}\zeta _0\mathrm{e}^{\eta (z^{})}\frac{\mathrm{d}^2B(zz^{})}{\mathrm{d}z^2}dz^{}$$
(52)
The PB density profile, $`c_{\mathrm{PB}}(z)\zeta _0\mathrm{e}^{\eta _0(z)}`$, for the same surface charge, satisfies the equation $`d^2\eta _0/dz^2=(4\pi \beta e^2\zeta _0/\epsilon )\mathrm{exp}(\eta _0)`$. Its exact solution is known to be:
$$c_{\mathrm{PB}}(z)=\zeta _0\mathrm{e}^{\eta _0(z)}=\frac{1}{2\pi l_\mathrm{B}}\frac{1}{(z+b)^2}$$
(53)
Note that only in the PB equation $`\eta (z)`$ is the reduced electrostatic potential $`e\mathrm{\Psi }(z)/k_\mathrm{B}T`$. From the generalized contact theorem (51), the surface density in the no added salt case and in the presence of one plate is $`c(0)=2\pi \beta \sigma ^2/\epsilon `$, as in PB theory. Therefore:
$$\eta (z=0)=\eta _0(z=0)$$
(54)
From the derivative of $`c(z)`$, Eq. (46), we find:
$$\frac{\mathrm{d}\eta }{\mathrm{d}z}=\beta e\frac{\mathrm{d}\mathrm{\Psi }}{\mathrm{d}z}_0^{\mathrm{}}dz^{}c(z^{})\frac{\mathrm{d}B(zz^{})}{\mathrm{d}z}$$
(55)
and using the boundary condition (39):
$$\frac{\mathrm{d}\eta }{\mathrm{d}z}|_{z=0}=\frac{\mathrm{d}\eta _0}{\mathrm{d}z}|_{z=0}+_0^{\mathrm{}}dz^{}c(z^{})\frac{\mathrm{d}B(z^{})}{\mathrm{d}z}$$
(56)
where the odd parity of $`\mathrm{d}B/\mathrm{d}z`$ has been used. This relation can be used together with Eq. (54) as a second boundary condition at $`z=0`$, instead of the boundary condition of vanishing $`\mathrm{d}\eta /\mathrm{d}z`$ at infinity.
Linearizing Eq. (52) with respect to:
$$w\eta \eta _0=\mathrm{log}(c/c_{\mathrm{PB}})$$
(57)
which is valid for relatively small deviations from the PB profile, results in the following equation:
$`{\displaystyle \frac{\mathrm{d}^2w}{\mathrm{d}z^2}}{\displaystyle \frac{4\pi }{\epsilon }}\beta e^2c_{\mathrm{PB}}(z)w(z)`$ (58)
$`={\displaystyle _0^{\mathrm{}}}dz^{}(1+w(z^{}))c_{\mathrm{PB}}(z^{}){\displaystyle \frac{\mathrm{d}^2B(zz^{})}{\mathrm{d}z^2}}`$ (59)
This equation can be further simplified by omitting $`w(z^{})`$ from the integrand in the right hand side. This approximation was motivated in Sec. III C and is equivalent to stopping the iterative scheme (48) after the first iteration. The density profile is then replaced by the PB density profile in the term that involves the hydration interaction $`B(z)`$. This results in the equation:
$$\frac{\mathrm{d}^2w}{\mathrm{d}z^2}\frac{4\pi }{\epsilon }\beta e^2c_{\mathrm{PB}}(z)w(z)+\mathrm{\Gamma }(z)=0$$
(60)
where $`\mathrm{\Gamma }(z)`$ is the convolution integral:
$$\mathrm{\Gamma }(z)=\frac{1}{2\pi l_\mathrm{B}}_0^{\mathrm{}}dz^{}\frac{1}{(z^{}+b)^2}\frac{\mathrm{d}^2B(zz^{})}{\mathrm{d}z^2}$$
(61)
The corresponding boundary conditions, obtained from equations (54) and (56) using the same approximations, are:
$`w(z=0)`$ $`=`$ $`0`$ (62)
$`{\displaystyle \frac{\mathrm{d}w}{\mathrm{d}z}}|_{z=0}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}dz^{}c_{\mathrm{PB}}(z^{}){\displaystyle \frac{\mathrm{d}B(z^{})}{\mathrm{d}z}}`$ (63)
Equation (60) is a second order linear differential equation for $`w(z)`$ and can be solved analytically. The solution, given in detail in Appendix B, is expressed in terms of the convolution integral $`\mathrm{\Gamma }(z)`$ of Eq. (61). The effective surface charge and the effect of the hydration on the density profile can then be calculated in several limits, described in detail in Appendix B. Here we outline the main results.
### A Slowly varying density: $`bd_{\mathrm{hyd}}`$
In the limit $`bd_{\mathrm{hyd}}`$, the PB distribution varies slowly on the scale of the hydration interaction, described by $`B(z)`$, and the theory becomes effectively a local density functional theory. The specific form of $`B(z)`$ is not important, and all the results simply depend on $`B_\mathrm{t}=_{\mathrm{}}^{\mathrm{}}B(z)dz`$. The deviation of the effective Gouy-Chapman length $`b_{\mathrm{eff}}`$ from the actual Gouy-Chapman length $`b`$ depends linearly on $`B_t`$ and on the surface charge $`\sigma 1/b`$. This can be expected since we use a linearized equation. Thus we have, on dimensional grounds, $`b_{\mathrm{eff}}bB_t/l_Bb`$. The detailed calculation gives the numerical prefactor:
$$b_{\mathrm{eff}}b\frac{B_\mathrm{t}}{4\pi l_\mathrm{B}}\frac{1}{b}$$
(64)
Since $`B_\mathrm{t}`$ is negative $`b_{\mathrm{eff}}`$ is larger than $`b`$ and the effective surface charge, $`\sigma _{\mathrm{eff}}`$, is smaller than the actual surface charge $`\sigma `$. This result should be expected. The short range interaction attracts counterions to the vicinity of the charged plate and the surface charge is screened more effectively than in the PB equation.
The correction to the counter-ion density profile, described by $`w(z)=\mathrm{log}[c(z)/c_{\mathrm{PB}}(z)]`$, is found to be:
$$w(z)=\frac{B_\mathrm{t}}{2\pi l_\mathrm{B}}\left\{\frac{3}{2(z+b)^2}\frac{1}{b(z+b)}\right\}$$
(65)
The density profile is increased relative to PB theory for distances smaller than $`b/2`$, and decreased for larger distances. The deviation from PB, $`w(z)`$, is maximal at $`z=0`$, where it is equal to $`B_\mathrm{t}/(4\pi l_\mathrm{B}b^2)`$, and minimal at $`z=2b`$, where it is equal to $`B_\mathrm{t}/(12\pi l_\mathrm{B}b^2)`$.
Figure 8 shows the approximated function $`w(z)`$ of Eq. (65) for $`b=21.2`$ Å, corresponding to $`b/d_{\mathrm{hyd}}3`$ (dotted line). The approximation is compared with the function $`w(z)`$ obtained from the exact solution of equation (46) for the case of no added salt (solid line). Although $`b`$ is not much larger than $`d_{\mathrm{hyd}}`$, the approximation describes well the correction to the PB profile. Note that $`w(z)`$, as expressed by Eq. (65) is maximal at $`z=0`$, whereas according to the contact theorem $`w(0)`$ should be zero. This apparent inconsistency results from neglecting the range of the hydration potential relative to $`b`$. In the precise solution of Eq. (46) $`w(0)`$ is zero, as it should be. The prediction of Eq. (65) is valid only for distances $`zd_{\mathrm{hyd}}`$, as can be seen in Figure 8.
The range of validity of the linearization procedure can be found by requiring that the minimal and maximal values of $`w(z)`$ are small compared to unity:
$$\frac{B_\mathrm{t}}{4\pi l_\mathrm{B}b^2}1$$
(66)
### B Surface layer limit: $`bd_{\mathrm{hyd}}`$
In the limit in which $`bd_{\mathrm{hyd}}`$, the ion density effectively becomes a dense layer concentrated at $`z=0`$ on the scale of the hydration interaction. The effective Gouy-Chapman length has the same form as in the limit of slowly varying density, $`bd_{\mathrm{hyd}}`$, but having a different prefactor:
$$b_{\mathrm{eff}}b\frac{B_\mathrm{t}}{12\pi l_\mathrm{B}}\frac{1}{b}$$
(67)
The effective surface charge is, therefore, smaller than the actual surface charge. Note that $`b_{\mathrm{eff}}`$ depends on $`B(z)`$, in this limit, only through $`B_\mathrm{t}`$. The linear dependence on $`\sigma 1/b`$ follows from the linearization leading to Eq. (60), as in the previous limit.
It should be stressed that although $`b`$ is small compared to $`d_{\mathrm{hyd}}`$ we still assume that $`b`$ is large enough for the linearization to be valid, i.e., we assume that $`w(z)`$ is small compared to unity. Furthermore, the counter-ion density should be small enough that we can sensibly use only the quadratic term in the virial expansion. To check the validity of these assumptions, the correction to the density profile should be considered.
The form of $`w(z)`$ depends, in the surface layer limit, on the specific form of $`B(z)`$. In order to study $`w(z)`$ analytically, we use an approximated form of $`B(z)`$, described in Appendix B. A typical form of the approximated $`w(z)`$, obtained using this approximation \[Eq. (B21)\], is shown in Fig. 9 (dotted line). The Gouy-Chapman length is $`b=1.06`$ Å, corresponding to $`b/d_{\mathrm{hyd}}0.15`$. In addition, the function $`w(z)`$ obtained from the exact solution of equation (46) is shown for comparison (solid line). The approximated curve captures well the qualitative behavior of the correction to the PB profile. Note that the discrepancy between the approximated and actual profiles results not only from the linearization and small $`b`$ limit, but also from the loss of detail due to the use of an approximated form for $`B(z)`$.
The deviation from the PB profile, $`w(z)`$, can be qualitatively described as follows. For $`z<d_{\mathrm{hc}}`$, $`w(z)`$ increases from zero quadratically (with an additional term of the form $`z^2\mathrm{log}z`$) to its value at $`z=d_{\mathrm{hc}}`$. It then decreases from its maximum positive value to a minimum, negative value, on a scale of the range of the attractive part of $`B(z)`$. This minimum value is equal to approximately $`B_\mathrm{t}/6\pi l_\mathrm{B}bd_{\mathrm{hyd}}`$. For distances larger than the interaction range, $`w(z)`$ assumes the form $`w(z)1/z`$, characterizing a PB profile with a modified, effective surface charge. For finite values of $`b`$, we can expect the above behavior to be smoothed over a scale of order $`b`$.
The validity of the linearization can be found by requiring that $`|w(z)|1`$. This requirement results in the following condition:
$$\frac{B_\mathrm{t}}{6\pi l_\mathrm{B}bd_{\mathrm{hyd}}}1$$
(68)
The validity of stopping the virial expansion at the quadratic order can be shown to have the same condition. For the hydration potential of Fig. 2, the condition expressed in Eq. (68) implies that the various approximations we use start to break down when $`b`$ becomes smaller than approximately $`1`$ Å, or $`\sigma 0.022e/\text{Å}^2`$. When $`b`$ is of this order, it is well below $`d_{\mathrm{hyd}}`$, making the surface layer limit a sensible approximation.
### C Effective surface charge
In the two limits described above, the effective Gouy-Chapman length was found to be of the form $`b_{\mathrm{eff}}bB_t/l_\mathrm{B}b`$, with different prefactors in the two limits. For intermediate values of $`b`$, the effective charge depends on the specific structure of the function $`B(z)`$. In order to study this dependence, we use a simple approximated form for $`B(z)`$, described in Appendix B. Using this approximation, an analytical expression can be obtained for the effective Gouy-Chapman length for all values of $`b`$.
Figures 10a and 10b show the predicted $`b_{\mathrm{eff}}`$ and $`b_{\mathrm{eff}}b`$, respectively (solid lines) as a function of $`b`$, together with the asymptotic limits (64) and (67) (dotted lines). As the surface charge increases from zero (and $`b`$ decreases from infinity), the effective charge $`|\sigma _{\mathrm{eff}}|`$ increases too (but is always smaller than the actual surface charge). When $`b`$ reaches a certain value $`b^{\mathrm{min}}`$, $`b_{\mathrm{eff}}`$ starts increasing with further reduction of $`b`$, i.e., the effective charge decreases with increasing surface charge above $`|\sigma |^{\mathrm{max}}=e/(2\pi l_\mathrm{B}b^{\mathrm{min}})`$. The value of $`b^{\mathrm{min}}`$ depends on the structure of the function $`B(z)`$, but can be estimated to be between the values predicted by the asymptotic expressions (64) and (67). From the condition $`\mathrm{d}b_{\mathrm{eff}}/\mathrm{d}b|_{b=b^{\mathrm{min}}}=0`$ we find:
$$\sqrt{\frac{B_\mathrm{t}}{12\pi l_\mathrm{B}}}<b^{\mathrm{min}}<\sqrt{\frac{B_\mathrm{t}}{4\pi l_\mathrm{B}}}$$
(69)
and:
$$b_{\mathrm{eff}}^{\mathrm{min}}2b^{\mathrm{min}}$$
(70)
For the hydration interaction of Fig. 2, $`B_\mathrm{t}`$ is approximately $`500\text{Å}^3`$. The value of $`b^{\mathrm{min}}`$ is then between $`1.36`$ Å and $`2.35`$ Å, corresponding to a surface charge density between $`0.15\mathrm{C}/\mathrm{m}^2`$ and $`0.26\mathrm{C}/\mathrm{m}^2`$. The values obtained from the approximated curve, shown in Fig. 10, are $`b^{\mathrm{min}}1.5\text{Å}`$ and $`b_{\mathrm{eff}}^{\mathrm{min}}3.4`$ Å.
For small enough values of $`b`$, the effective surface charge $`|\sigma _{\mathrm{eff}}|`$ should increase again with an increase of $`|\sigma |`$ and become larger than $`|\sigma |`$. This effect cannot be predicted by our model because of the low density approximation used for the hard core interaction. In particular, the hard core of the ions should cause the density to saturate at the close packing density, leading to a reduced screening of the surface charge relative to PB theory . In our model, as in the PB theory, the counterion density near the surface is not bounded, and increases indefinitely as $`\sigma `$ is increased. Although our model includes the steric repulsion between ions, this repulsion is “softened”, and is always outweighed by the attractive part of the ion-ion interaction.
In addition to the prediction obtained using the linearized approximation, Figure 10 shows values of $`b_{\mathrm{eff}}`$ extracted from numerical solutions of the full equation (43), using the original interaction $`B(z)`$. The equation was solved with two different salt concentrations: $`10^7`$ M (circles) and $`0.1`$ M (crosses). The value of $`b_{\mathrm{eff}}`$ was estimated from the positive ion density at large distances from the plate, by finding the value of $`b`$ that would result in the same calculated values of the density in a solution of the PB equation. Note that for both salt concentrations, $`b_{\mathrm{eff}}`$ is very close to its predicted value, meaning that the salt has a very small effect on $`\sigma _{\mathrm{eff}}`$. This result is not obvious for the high salt concentration of $`0.1\mathrm{M}`$. The Debye-Hückel screening length is approximately $`9.6`$ Å, not much larger than the range of the hydration interaction, $`d_{\mathrm{hyd}}7\text{Å}`$, and comparable to the Gouy-Chapman length at the large $`b`$ region of the plot.
## V Conclusions and Outlook
In this work we have studied the effects due to the discreteness of the solvent in aqueous ionic solutions. Hydration interactions are found to have a significant effect on the structure of the diffusive layer near highly charged surfaces. The counter-ion density is increased in the vicinity of the charged surface, relative to the PB prediction, and decreased further away. The distance from the charged plate in which the density is increased, and the magnitude of the deviation from the PB density, depend strongly on the surface charge, and on the parameters of the short-range hydration interaction between ion pairs.
The ion-ion hydration interaction can be described roughly using two parameters. The first parameter is the range of the hydration interaction, $`d_{\mathrm{hyd}}`$, equal to approximately $`7\text{Å}`$ for Na<sup>+</sup>-Na<sup>+</sup> pairs. The second parameter, $`B_\mathrm{t}`$ has dimensions of volume and characterizes the strength of the hydration interaction. It is equal to approximately $`500\text{Å}^3`$ for Na<sup>+</sup>-Na<sup>+</sup> pairs. Two limits can be considered, where the Gouy-Chapman length, $`b1/\sigma `$, is small or large compared to the range of the hydration interaction $`d_{\mathrm{hyd}}`$. In both of these limits we assume that the Debye-Hückel screening length, $`\lambda _\mathrm{D}`$, is large compared to $`b`$ and $`d_{\mathrm{hyd}}`$.
In the limit $`bd_{\mathrm{hyd}}`$, the counter-ion density becomes depleted, relative to the PB prediction, starting at a distance $`zb/2`$ from the charged plate. The maximum absolute value of $`w(z)=\mathrm{log}[c(z)/c_{\mathrm{PB}}(z)]`$ scales as $`B_\mathrm{t}/l_\mathrm{B}b^2`$. In the limit $`bd_{\mathrm{hyd}}`$, the distance from the plate, where the counter-ion density becomes lower than the PB prediction, is between $`z=d_{\mathrm{hc}}`$ and $`z=d_{\mathrm{hyd}}`$. The maximum absolute value of $`w(z)`$ scales as $`B_\mathrm{t}/l_\mathrm{B}d_{\mathrm{hyd}}b`$.
Far away from the charged plate, the density profile can be well described using the PB theory with an effective surface charge that can be calculated analytically. The correction to the Gouy-Chapman length in the two limits $`bd_{\mathrm{hyd}}`$ and $`bd_{\mathrm{hyd}}`$ is always positive and scales as $`B_\mathrm{t}/l_\mathrm{B}b`$, but has different numerical prefactors. When the surface charge on the plate is increased, the effective surface charge, $`\sigma _{\mathrm{eff}}`$, is found to reach a certain maximal value. Above this maximal value $`\sigma _{\mathrm{eff}}`$ decreases with further increase of the actual $`\sigma `$ on the plate. The various approximations we use start to break down when $`b`$ is smaller than approximately $`B_\mathrm{t}/6\pi l_\mathrm{B}d_{\mathrm{hyd}}`$, corresponding to $`b1\text{Å}`$.
An important outcome of this work is that the correction of the PB ion density due to the hydration interaction is significant near highly charged surfaces. The electrostatic interaction dominates the ionic distribution and the hydration interaction can be seen as a perturbation. For a high surface charge density of, say, one unit charge per $`48\text{Å}^2`$ the counterion density deviates from its Poisson Boltzmann value by at most $`30`$ percent. The effective change in the surface charge is more significant, from $`1e/48\text{Å}^2`$ to about $`1e/13\text{Å}^2`$.
The hydration effect on inter-surface forces can be very pronounced, as opposed to the effect on the ion distribution. This result will be presented elsewhere . Our model predicts an attractive contribution to the pressure between two parallel charged plates. At distances below several nanometers this contribution can outweigh the electrostatic repulsion and lead to an overall attraction between the plates. Our two-plate findings can also be compared with available AHNC results , showing good qualitative agreement both for the ion density profile and pressure.
The formalism we present can be readily generalized to other geometries. This could lead to an estimation of the aqueous solvent effects on phenomena such as the Manning condensation on cylindrical polyions , and charge renormalization of spherical mycelles or colloids . In this respect our formalism offers an advantage over the AHNC approximation which was applied so far only in a planar geometry. Another interesting extension of this work would be to consider the combination of fluctuation and hydration effects. This is particularly important for ionic solutions with divalent counter-ions, where fluctuation effects become large .
###### Acknowledgements.
We wish to thank S. Marc̆elja for introducing us to the subject of solvent effects in aqueous ionic solutions, and for valuable discussions and suggestions. We would like to thank R. Netz, H. Orland and R. Podgornik for useful discussions. Partial support from the U.S.-Israel Binational Foundation (B.S.F.) under grant No. 98-00429, and the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities - Centers of Excellence Program is gratefully acknowledged.
## A Inhomogeneous virial expansion
We consider an inhomogeneous system of particles with a short-range two-body interaction, and aim to express the free energy of the system in the low density limit as a functional of the density distribution. For simplicity we consider only one species of particles. The inhomogeneity of the system arises from the inclusion of an external field $`\phi (𝐫)`$, or from the boundary conditions imposed on the system. We begin by considering the grand canonical ensemble. The grand canonical partition function is:
$$Z_\mathrm{G}=\underset{N}{}\frac{1}{N!}\left(\frac{\mathrm{e}^{\beta \mu }}{\lambda _\mathrm{T}^3}\right)^NQ_N$$
(A1)
where $`\mu `$ is the chemical potential, $`\lambda _\mathrm{T}`$ is the de Broglie thermal wavelength and $`Q_N`$ is:
$$Q_N=\underset{i=1}{\overset{N}{}}\mathrm{d}^3𝐫_i\mathrm{e}^{\beta U_N(\{𝐫_i\})}$$
(A2)
$$U_N(\{𝐫_i\})=\underset{i}{}\phi (𝐫_i)+\frac{1}{2}\underset{i}{}\underset{ji}{}u\left(\left|𝐫_i𝐫_j\right|\right)$$
(A3)
We proceed on similar lines as the usual virial expansion in a bulk fluid, expanding $`\mathrm{log}Z_\mathrm{G}`$ in powers of the activity. Up to second order we have:
$`\mathrm{log}Z_\mathrm{G}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{e}^{\beta \mu }}{\lambda _\mathrm{T}^3}}\right)Q_1+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\mathrm{e}^{\beta \mu }}{\lambda _\mathrm{T}^3}}\right)^2\left(Q_2Q_1^2\right)=\left({\displaystyle \frac{\mathrm{e}^{\beta \mu }}{\lambda _\mathrm{T}^3}}\right){\displaystyle \mathrm{d}^3𝐫\mathrm{e}^{\beta \phi (𝐫)}}`$ (A5)
$`+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\mathrm{e}^{\beta \mu }}{\lambda _\mathrm{T}^3}}\right)^2{\displaystyle \mathrm{d}^3𝐫\mathrm{d}^3𝐫^{}\mathrm{e}^{\beta \left(\phi (𝐫)+\phi (𝐫^{})\right)}\left(\mathrm{e}^{\beta u\left(\left|𝐫𝐫^{}\right|\right)}1\right)}`$
This can be seen as an expansion in powers of the field $`\mathrm{exp}\left[\beta \left(\mu \phi (𝐫)\right)\right]/\lambda _\mathrm{T}^3`$. The local density $`c(𝐫)`$ can be expressed in a similar expansion:
$`c(𝐫)`$ $`=`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \frac{\delta \mathrm{log}Z_\mathrm{G}}{\delta \phi (𝐫)}}=\left({\displaystyle \frac{\mathrm{e}^{\beta \mu }}{\lambda _\mathrm{T}^3}}\right)\mathrm{e}^{\beta \phi (𝐫)}`$ (A7)
$`+\left({\displaystyle \frac{\mathrm{e}^{\beta \mu }}{\lambda _\mathrm{T}^3}}\right)^2\mathrm{e}^{\beta \phi (𝐫)}{\displaystyle \mathrm{d}^3𝐫^{}\mathrm{e}^{\beta \phi (𝐫^{})}\left(\mathrm{e}^{\beta u\left(\left|𝐫𝐫^{}\right|\right)}1\right)}`$
This relation can be inverted to obtain an expansion of $`\mathrm{exp}\left[\beta \left(\mu \phi (𝐫)\right)\right]/\lambda _\mathrm{T}^3`$ in powers of $`c(𝐫)`$. Up to the second order:
$$\frac{\mathrm{e}^{\beta (\mu \phi (𝐫))}}{\lambda _\mathrm{T}^3}=c(𝐫)+c(𝐫)\mathrm{d}^3𝐫^{}c(𝐫^{})\left(1\mathrm{e}^{\beta u\left(\left|𝐫𝐫^{}\right|\right)}\right)$$
(A8)
and by substituting this relation in Eq. (A5) $`\mathrm{log}Z_\mathrm{G}`$ can be expressed as an expansion in $`c`$. Up to the second order:
$$\mathrm{log}Z_\mathrm{G}=\mathrm{d}^3𝐫c(𝐫)+\frac{1}{2}\mathrm{d}^3𝐫\mathrm{d}^3𝐫^{}c(𝐫)c(𝐫^{})\left(1\mathrm{e}^{\beta u\left(\left|𝐫𝐫^{}\right|\right)}\right)$$
(A9)
The grand canonical potential can be obtained from the relation $`\mathrm{\Omega }=k_\mathrm{B}T\mathrm{log}Z_\mathrm{G}`$, with $`\mathrm{log}Z_\mathrm{G}`$ given by Eq. (A9). In this expression, $`c(𝐫)`$ is the mean density profile for the imposed external field $`\phi (𝐫)`$ and a given chemical potential $`\mu `$. We would like to express $`\mathrm{\Omega }`$ as a functional of a general ion density $`c(𝐫)`$, whose minimization with respect to $`c(𝐫)`$ would give the equilibrium mean density. Regarding $`k_\mathrm{B}T\mathrm{log}Z_\mathrm{G}`$ as a functional of $`\chi (𝐫)\phi (𝐫)\mu `$, we have:
$$k_\mathrm{B}T\frac{\delta \mathrm{log}Z_\mathrm{G}}{\delta \chi (𝐫)}=c(𝐫)$$
(A10)
The Legendre transform of this relation can be obtained by defining:
$$\mathrm{\Theta }=k_\mathrm{B}T\mathrm{log}Z_\mathrm{G}\mathrm{d}^3𝐫c(𝐫)\chi (𝐫)$$
(A11)
and expressing $`\mathrm{log}Z_\mathrm{G}`$ and $`\chi `$ as functionals of $`c(𝐫)`$. We have already expressed $`\mathrm{log}Z_\mathrm{G}`$ as a functional of $`c(𝐫)`$ in Eq. (A9). An expression for $`\chi (𝐫)`$ as a functional of $`c(𝐫)`$ can be obtained from Eq. (A8). Up to first order in $`c`$ we have :
$`\beta [\phi (𝐫)\mu ]`$ $`=`$ $`\mathrm{log}\left\{\lambda _\mathrm{T}^3c(𝐫)\left[1+{\displaystyle \mathrm{d}^3𝐫^{}c(𝐫^{})\left(1\mathrm{e}^{\beta u\left(\left|𝐫𝐫^{}\right|\right)}\right)}\right]\right\}`$ (A12)
$`=`$ $`\mathrm{log}\left[\lambda _\mathrm{T}^3c(𝐫)\right]{\displaystyle \mathrm{d}^3𝐫^{}c(𝐫^{})\left(1\mathrm{e}^{\beta u\left(\left|𝐫𝐫^{}\right|\right)}\right)}+\mathrm{O}(c^2)`$ (A13)
Using this relation and Eq. (A9) we obtain, up to second order in c:
$`\beta \mathrm{\Theta }(\{c(𝐫)\})`$ $`=`$ $`{\displaystyle \mathrm{d}^3𝐫c(𝐫)\left[\mathrm{log}(\lambda _\mathrm{T}^3c(𝐫))1\right]}`$ (A15)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \mathrm{d}^3𝐫\mathrm{d}^3𝐫^{}c(𝐫)c(𝐫^{})\left(1\mathrm{e}^{\beta u\left(\left|𝐫𝐫^{}\right|\right)}\right)}`$
The functional $`\mathrm{\Theta }`$ of $`c(𝐫)`$ has the property that:
$$\frac{\delta \mathrm{\Theta }}{\delta c(𝐫)}=\chi (𝐫)=\left[\phi (𝐫)\mu \right]$$
(A16)
or equivalently:
$$\frac{\delta }{\delta c(𝐫)}\left\{\mathrm{\Theta }+\mathrm{d}^3𝐫c(𝐫)\left[\phi (𝐫)\mu \right]\right\}=\frac{\delta \mathrm{\Omega }(\{c(𝐫)\})}{\delta c(𝐫)}=0$$
(A17)
Thus, using Eq. (A15), we obtain:
$`\mathrm{\Omega }(\{c(𝐫)\})`$ $`=`$ $`k_\mathrm{B}T{\displaystyle \mathrm{d}^3𝐫c(𝐫)\left(\mathrm{log}\frac{c(𝐫)}{\zeta }1\right)}+{\displaystyle \mathrm{d}^3𝐫c(𝐫)\phi (𝐫)}`$ (A19)
$`+{\displaystyle \frac{1}{2}}k_\mathrm{B}T{\displaystyle \mathrm{d}^3𝐫\mathrm{d}^3𝐫^{}c(𝐫)c(𝐫^{})\left(1\mathrm{e}^{\beta u\left(\left|𝐫𝐫^{}\right|\right)}\right)}`$
where $`\zeta =\mathrm{exp}(\beta \mu )/\lambda _\mathrm{T}^3`$.
The derivation of Eq. (A19) can be readily generalized to the case of several ion species of different charges and different pair interactions $`u_{ij}(𝐫)`$, resulting in Eq. (24).
A similar, more elaborate diagrammatic expansion of the thermodynamic potentials in the presence of an external field is presented in Ref. . A variational principal for the grand canonical potential $`\mathrm{\Omega }`$ is obtained in which $`\mathrm{\Omega }`$ is expressed as a functional of the mean density $`c(𝐫)`$ and the pair correlation function $`h_2(𝐫_1,𝐫_2)`$. This expression is equivalent to Eq. (A19) up to the second order in the cluster expansion.
## B Details of analytical results
In this appendix we present details of the analytical approximations of Sec. IV.
We consider first the analytical solution of Equation (60). This equation is a second order linear differential equation for $`w(z)`$. Note that the function $`c_{\mathrm{PB}}(z)`$ is a known function of $`z`$, given by Eq. (53). The solution of Eq. (60), with the boundary conditions of Eq. (63) is:
$$w(z)=\frac{1}{z+b}_0^zdz_2\left(z_2+b\right)^2_{z_2}^{\mathrm{}}\frac{\mathrm{d}z_1}{\left(z_1+b\right)}\mathrm{\Gamma }(z_1)$$
(B1)
where $`\mathrm{\Gamma }(z)`$ is the convolution integral, defined by Eq. (61). By writing $`\mathrm{\Gamma }(z)`$ as:
$$\mathrm{\Gamma }(z)=_0^{\mathrm{}}dz^{}\mathrm{\Gamma }(z^{})\delta (zz^{})$$
(B2)
$`w(z)`$ can be rewritten in the following form:
$`w(z)`$ $`=`$ $`{\displaystyle \frac{1}{z+b}}\left\{{\displaystyle \frac{b^3}{3}}{\displaystyle _0^{\mathrm{}}}dz^{}{\displaystyle \frac{\mathrm{\Gamma }(z^{})}{z^{}+b}}{\displaystyle \frac{1}{3}}{\displaystyle _0^z}dz^{}(z^{}+b)^2\mathrm{\Gamma }(z^{})\right\}`$ (B4)
$`+{\displaystyle \frac{(z+b)^2}{3}}{\displaystyle _z^{\mathrm{}}}dz^{}{\displaystyle \frac{\mathrm{\Gamma }(z^{})}{z^{}+b}}`$
The effective charge $`\sigma _{\mathrm{eff}}`$ (or equivalently, the effective Gouy-Chapman length $`b_{\mathrm{eff}})`$ can be calculated from the coefficient of $`z^1`$ in $`w(z)`$, as $`z`$ approaches infinity:
$$w(z)\frac{2\left(bb_{\mathrm{eff}}\right)}{z},z\mathrm{}$$
(B5)
We thus find:
$$b_{\mathrm{eff}}b=\frac{1}{6}_0^{\mathrm{}}dz\left[\frac{b^3}{z+b}(z+b)^2\right]\mathrm{\Gamma }(z)$$
(B6)
A simple form for the convolution integral $`\mathrm{\Gamma }(z)`$ can be obtained in the limits in which $`b`$ is small or large relative to $`d_{\mathrm{hyd}}`$, the characteristic range of the hydration potential.
### 1 Slowly varying density: $`bd_{\mathrm{hyd}}`$
In the limit $`bd_{\mathrm{hyd}}`$, the PB distribution varies slowly on the scale of the hydration interaction. The convolution integral $`\mathrm{\Gamma }(z)`$ of Eq. (61) can then be approximated in the following way:
$`\mathrm{\Gamma }(z)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi l_\mathrm{B}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}dz^{}{\displaystyle \frac{H(z^{})}{(z^{}+b)^2}}{\displaystyle \frac{\mathrm{d}^2B}{\mathrm{d}z^2}}(zz^{})`$ (B7)
$`=`$ $`{\displaystyle \frac{1}{2\pi l_\mathrm{B}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}dz^{}\left[{\displaystyle \frac{1}{b^2}}{\displaystyle \frac{\mathrm{d}\delta (z)}{\mathrm{d}z}}{\displaystyle \frac{2}{b^3}}\delta (z)+{\displaystyle \frac{6H(z)}{(z+b)^4}}\right]B(zz^{})`$ (B8)
$``$ $`{\displaystyle \frac{B_\mathrm{t}}{2\pi l_\mathrm{B}}}\left[{\displaystyle \frac{1}{b^2}}{\displaystyle \frac{\mathrm{d}\delta (z)}{\mathrm{d}z}}{\displaystyle \frac{2}{b^3}}\delta (z)+{\displaystyle \frac{6H(z)}{(z+b)^4}}\right]`$ (B9)
where $`H(z)`$ is the Heaviside function ($`H(z)`$ = 0 for $`z<0`$ and H(z) = 1 for $`z>0`$). Inserting this expression in Eq. (B6) we obtain Eq. (64) for the effective Gouy-Chapman length. By substituting equation (B9) in Eq. (B4), the form of $`w(z)`$, given in Eq. (65), is obtained.
### 2 Approximated form for $`B(z)`$
Some of the following results depend on the specific structure of the hydration interaction, characterized by the function $`B(z)`$. In order to obtain analytical expressions, we use a simple approximated form, $`B^{\mathrm{app}}(z)`$, instead of $`B(z)`$. Assuming that the hydration interaction consists of a hard core interaction and a short-range attractive part, the function $`B(z)`$ has some general characteristics that should be present in $`B^{\mathrm{app}}(z)`$. For $`z<d_{\mathrm{hc}}`$, $`B(z)`$ always has the parabolic form $`(B_0+\pi z^2)`$, where $`B_0=B(z=0)`$. We assume that the attractive part of the interaction dominates over the short-range repulsion so that $`B_0`$ is positive. For $`z`$ larger than some finite value $`d_{\mathrm{hc}}+\mathrm{\Delta }`$, $`B(z)`$ is practically zero due to the short range of the interaction. For $`d_{\mathrm{hc}}<z<d_{\mathrm{hc}}+\mathrm{\Delta }`$, $`B(z)`$ varies from $`(B_0+\pi d_{\mathrm{hc}}^2)`$ to zero in a functional form that depends on the details of the attractive potential. The most simple way to model this behavior of $`B(z)`$ is to have a linear increase of $`B^{\mathrm{app}}(z)`$ between $`z=d_{\mathrm{hc}}`$ and $`z=d_{\mathrm{hc}}+\mathrm{\Delta }`$, and to set $`B^{\mathrm{app}}`$ to be zero for $`z`$ larger than $`d_{\mathrm{hc}}+\mathrm{\Delta }`$:
$$B^{\mathrm{app}}(z)=\{\begin{array}{cc}(B_0+\pi z^2)\hfill & \left|z\right|d_{\mathrm{hc}}\hfill \\ \left(B_0+\pi d_{\mathrm{hc}}^2\right)\frac{\left(d_{\mathrm{hc}}+\mathrm{\Delta }z\right)}{\mathrm{\Delta }}\hfill & d_{\mathrm{hc}}<\left|z\right|d_{\mathrm{hc}}+\mathrm{\Delta }\hfill \\ 0\hfill & d_{\mathrm{hc}}+\mathrm{\Delta }<\left|z\right|\hfill \end{array}$$
(B10)
The parameters in this expression should be chosen to match, approximately, the form of $`B(z)`$. Setting $`d_{\mathrm{hc}}`$ to be the hard core diameter of the real potential and setting $`B_0=B(0)`$ ensures that $`B(z)`$ and $`B^{\mathrm{app}}(z)`$ are identical for $`z<d_{\mathrm{hc}}`$. The width $`\mathrm{\Delta }`$ can then be set such that $`B_\mathrm{t}^{\mathrm{app}}=B_\mathrm{t}`$:
$$2B_0d_{\mathrm{hc}}+\frac{2}{3}d_{\mathrm{hc}}^3+\mathrm{\Delta }\left(B_0+\pi d_{\mathrm{hc}}^2\right)=B_\mathrm{t}$$
(B11)
This is desirable in light of equations (64) and (67), since the effective surface charge depends only on $`B_\mathrm{t}`$ in these limits. Figure 11 shows $`B(z)`$ and $`B^{\mathrm{app}}(z)`$ for the hydration potential of Fig. 2.
### 3 Surface layer limit: $`bd_{\mathrm{hyd}}`$
In the limit $`bd_{\mathrm{hyd}}`$, the convolution integral in Eq. (61) becomes:
$$\mathrm{\Gamma }(z)\frac{|\sigma |}{e}\frac{\mathrm{d}^2B(z)}{\mathrm{d}z^2}=\frac{1}{2\pi l_\mathrm{B}b}\frac{\mathrm{d}^2B(z)}{\mathrm{d}z^2}$$
(B12)
The prefactor of $`\mathrm{\Gamma }(z)`$ in Eq. (B6) is -$`\frac{1}{6}z^2+O(b)`$ and therefore the effective Gouy-Chapman length is:
$$b_{\mathrm{eff}}b\frac{1}{12\pi l_\mathrm{B}b}_0^{\mathrm{}}dzz^2B^{\prime \prime }(z)=\frac{B_\mathrm{t}}{12\pi l_\mathrm{B}}\frac{1}{b}$$
(B13)
This result is independent on the specific form of $`B(z)`$.
To obtain $`w(z)`$, the deviation of the density profile relative to PB theory, Eq. (B12) can be substituted in Eq. (B4). Up to leading order in $`b`$ the following expression is obtained:
$$w(z)=\frac{1}{6\pi l_\mathrm{B}b}\frac{1}{z}_0^zdz^{}B^{\prime \prime }(z^{})z^2+\frac{1}{6\pi l_\mathrm{B}b}z^2_z^{\mathrm{}}dz^{}\frac{1}{z^{}}B^{\prime \prime }(z^{})$$
(B14)
Using $`B^{\mathrm{app}}(z)`$, the approximated form of $`B(z)`$ presented in the previous subsection, this gives:
$`w(z)=\{\begin{array}{cc}{\displaystyle \frac{1}{6\pi l_\mathrm{B}b}}z^2\left({\displaystyle \frac{4\pi }{3}}+{\displaystyle \frac{B_0+\pi d_{\mathrm{hc}}^2}{d_{\mathrm{hc}}(d_{\mathrm{hc}}+\mathrm{\Delta })}}2\pi \mathrm{log}{\displaystyle \frac{d_{\mathrm{hc}}}{z}}\right),\hfill & \hfill \left|z\right|d_{\mathrm{hc}}\\ & \\ {\displaystyle \frac{1}{6\pi l_\mathrm{B}b}}[d_{\mathrm{hc}}^2({\displaystyle \frac{4\pi }{3}}d_{\mathrm{hc}}+{\displaystyle \frac{B_0+\pi d_{\mathrm{hc}}^2}{\mathrm{\Delta }}}){\displaystyle \frac{1}{z}}\hfill & \\ {\displaystyle \frac{B_0+\pi d_{\mathrm{hc}}^2}{\mathrm{\Delta }(d_{\mathrm{hc}}+\mathrm{\Delta })}}z^2],\hfill & \hfill d_{\mathrm{hc}}<\left|z\right|d_{\mathrm{hc}}+\mathrm{\Delta }\\ & \\ {\displaystyle \frac{B_\mathrm{t}}{6\pi l_\mathrm{B}b}}{\displaystyle \frac{1}{z}},\hfill & \hfill d_{\mathrm{hc}}+\mathrm{\Delta }<\left|z\right|\end{array}`$ (B21)
The minimal, negative value of $`w(z)`$ is assumed at $`z=d_{\mathrm{hc}}+\mathrm{\Delta }`$ and is equal to:
$$w(d_{\mathrm{hc}}+\mathrm{\Delta })=\frac{B_\mathrm{t}}{6\pi l_\mathrm{B}b(d_{\mathrm{hc}}+\mathrm{\Delta })}\frac{B_\mathrm{t}}{6\pi l_\mathrm{B}bd_{\mathrm{hyd}}}$$
(B22)
This results in the condition (68) for the validity of the linearization in the surface layer limit.
Using only the quadratic term in the virial expansion is sensible if $`_0^{\mathrm{}}dz^{}c(z^{})B(zz^{})`$ is small compared to unity. In the surface layer limit, this integral is simply: $`(|\sigma |/e)B(z)=B(z)/(2\pi l_\mathrm{B}b)`$. Estimating the maximum value of $`|B(z)|`$ to be approximately $`B_\mathrm{t}/(2d_{\mathrm{hyd}})`$ we obtain the requirement: $`B_\mathrm{t}/(4\pi l_\mathrm{B}bd_{\mathrm{hyd}})1`$ , which is analogous to (68).
### 4 Effective Gouy-Chapman length
Using $`B^{\mathrm{app}}(z)`$ in equations (61) and (B6) we find the following approximation for the effective Gouy-Chapman length:
$`b_{\mathrm{eff}}b`$ $`=`$ $`{\displaystyle \frac{1}{12\pi l_\mathrm{B}b}}\{B_\mathrm{t}^{\mathrm{app}}\pi d_{\mathrm{hc}}^2b+2\pi d_{\mathrm{hc}}b^2+2B_0\mathrm{log}\left({\displaystyle \frac{b+d_{\mathrm{hc}}+\mathrm{\Delta }}{b}}\right)b`$ (B25)
$`{\displaystyle \frac{2}{\mathrm{\Delta }}}\left(\pi d_{\mathrm{hc}}^2\mathrm{\Delta }+\pi d_{\mathrm{hc}}^3+B_0d_{\mathrm{hc}}\right)\mathrm{log}\left({\displaystyle \frac{b+d_{\mathrm{hc}}}{b+d_{\mathrm{hc}}+\mathrm{\Delta }}}\right)b`$
$`2\pi \mathrm{log}\left({\displaystyle \frac{b+d_{\mathrm{hc}}}{b}}\right)b^3\}`$
This expression is shown in Fig. 10 and discussed in section IV. In the limits $`bd_{\mathrm{hyd}}`$ and $`bd_{\mathrm{hyd}}`$ it reduces to the asymptotic expressions (64) and (67), respectively.
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# Nucleon-nucleon elastic scattering to 3 GeV
## I Introduction
This analysis of elastic nucleon-nucleon scattering data updates our previous analyses to 1.6 GeV and 2.5 GeV in the laboratory kinetic energy. The present analysis has been extended to 3 GeV in order to include both the elastic $`pp`$ polarized measurements $``$ from SATURNE II at Saclay, and $`pp`$ differential cross sections measured by the EDDA collaboration using the cooler synchrotron at COSY. A detailed description of our database is given in Section II.
As discussed in Ref., the region beyond 2 GeV is interesting for several reasons. These include the suggestion of a narrow dibaryon resonance, corresponding to a center-of-mass energy of 2.7 GeV. Near this energy, a sharp structure was found in preliminary $`A_{yy}`$ measurements and this was taken as support for such a resonance. The authors of Ref. considered this possibility but found no evidence in their differential cross section measurements. No significant anomaly was seen in the angular and energy dependence of detailed analyzing power and correlated spin measurement at Saclay.
The possible onset of behavior suggested by dimensional counting at high energy and fixed center-of-mass angle is also intriguing. In Ref., we noted that $`d\sigma /dt`$ appeared to be approaching $`s^{10}`$, as expected within perturbative QCD. Thus, an extended energy range is needed to verify this trend. An extension of the $`np`$ analysis beyond 1.3 GeV would also benefit those studying the two-body photodisintegration of the deuteron at Jefferson Lab, which shows an interesting scaling behavior at some scattering angles . Unfortunately, the $`np`$ data base remains too sparse to support a reliable analysis beyond 1.3 GeV.
In the present work, we have focused mainly on the influence of new polarization data at higher energies, and on the behavior of $`ϵ_1`$ at low energies. The Saclay group has recently performed a single-energy phase-shift analysis of elastic $`pp`$ scattering data to 2.7 GeV and $`np`$ elastic scattering data to 1.1 GeV. In this study, a second set of amplitudes was obtained through a direct reconstruction of the scattering amplitudes at fixed energies and angles. We have compared our results to these and, in some cases, find evidence (complimentary to that given by the Saclay group) for non-uniqueness at higher energies. At lower energies, where the behavior of $`ϵ_1`$ has been a source of controversy, we compare our results to those of several other groups and suggest there is little evidence for an anomalously large tensor interaction. Results of our analyses are displayed in Section III. In Section IV, we summarize our findings and conclusions.
## II Database
Our two previous $`pp`$ scattering analyses extended to 1.6 and 2.5 GeV, respectively. In both cases, the associated $`np`$ analysis was restricted to 1.3 GeV. The present data base is considerably larger due both to an expanded energy range for the $`pp`$ system and the addition of new data at lower energies. The full data base has increased by 30% since the publication of Ref., and is about 70% larger than the set available for the analysis of Ref.. (The total data base has doubled over the last decade (see Table I).) The distribution of recent (post-1997) $`pp`$ and $`np`$ data is given in Fig. 1.
In the full data base, one will occasionally find experiments which give conflicting results. Some of these have been excluded from our fits. We have, however, retained all available data sets so that comparisons can be made through our on-line facility . Below, we list recent additions to our data base. Some of the data listed as new were available, in unpublished form, at the time of our previous analysis. A complete description of the data base and those data not included in our analyses is available from the authors.
Two thirds of the 4802 new $`pp`$ polarization data were produced at Saclay using the SATURNE II accelerator $``$ . These measurements of 9 spin-dependent quantities have increased our data base by a factor of two over the energy range from 1.6 to 2.5 GeV and accounted for a third of the data from 2.5 to 3.0 GeV. Many of the new $`pp`$ polarization measurements below 450 MeV were made at the Indiana cooler (A<sub>y</sub>, A<sub>yy</sub>, A<sub>xx</sub>, A<sub>zz</sub>, and A<sub>xz</sub>) $``$ . Also, in this energy range are 11 new unpolarized cross sections (at 398 MeV) measured at the Osaka facility .
The $`np`$ data base has not increased significantly. As a result, we have not extended the range of our analyses for the $`I=0`$ system beyond 1.3 GeV. The Geneva group, working at PSI, has provided 247 new $`np`$ spin measurements. From this source, we have also added spin observables A<sub>y</sub>, A<sub>t</sub>, D<sub>t</sub>, and R<sub>t</sub> from 260 to 538 MeV . About 60% of the recent SATURN II $`np`$ polarized measurements fall within our energy range; the full range extends from about 1.1 to 2.4 GeV (182 data points) . A few measurements of $`\mathrm{\Delta }\sigma _L`$ and $`\mathrm{\Delta }\sigma _T`$ were produced by TUNL and Charles University at Prague . A single measurement of D<sub>t</sub>, at 16 MeV, was provided by the ISKP cyclotron at Bonn. We have added 12 unpolarized measurements, between 29 and 73 MeV, from the Louvain-la-Neuve Cyclotron . Finally, we have retained in the analysis a set of 82 total cross section measurements, between 4 and 231 MeV , which had earlier been removed in order to have a low-energy data base identical to that used by the Nijmegen group.
## III Partial-Wave Analyses
Fits to the expanded data base and extended energy range were found to be possible within the formalism used and described in our previous analyses . Both energy-dependent and single-energy solutions were obtained from fits to the combined $`pp`$ and $`np`$ data bases to 1.3 GeV, and from fits to the $`pp`$ data base alone from 1.3 to 3 GeV. In Table II, we compare the energy-dependent and single-energy results over the energy bins used in the single-energy analyses. Also listed are the number of parameters varied in each single-energy solution. A total of 147 parameters were varied in the energy-dependent analysis (SP00).
Our single-energy and energy-dependent results for the isovector and isoscalar partial-wave amplitudes are displayed in Figs. 2 and 3. Here, we also compare with our previous fit (SM97) and a much older fit (FA91). Partial waves with $`J<6`$ are displayed, whereas the analysis fitted waves up to $`J=7`$. In most cases, SP00 and SM97 are in good agreement. Somewhat larger changes are seen in comparisons with FA91. Differences are generally largest, as one would expect, near the energy upper limits for the various solutions and in the smaller partial waves. Figs. 2 and 3 therefore show that a doubling of the data base has led to a refinement of the amplitudes, but has not required a dramatic change in their behavior.
Single-energy solutions were produced up to 2825 MeV (for $`pp`$ scattering). In these fits, initial values for the partial-wave amplitudes and their (fixed) energy derivatives were obtained from the energy-dependent solution. A comparison of global and single-energy solutions then serves as a check for structures that could have been “smoothed over” in the energy-dependent analysis. However, structures with widths less than 10 MeV would be very difficult to detect.
In order to ascertain that the extension to 3 GeV (1.3 GeV for $`np`$) did not seriously degrade our low-energy results, a 0 $``$ 400 MeV fit was also developed. The resultant solution, SP40, used 30 $`I=1`$ and 27 $`I=0`$ variable parameters to give a $`\chi ^2`$/datum of 4398/3454 ($`pp`$), and 5415/3831 ($`np`$). The global fit, SP00, produced, for the same energy range, a $`\chi ^2`$/datum of 4593/3454 ($`pp`$) and 5371/3831 ($`np`$). We consider this quite reasonable given that the number of variable parameters per datum is nearly three times larger for SP40 than for SP00.
In Figs. 4 and 5, we compare our results with the Saclay single-energy analyses for isovector waves below 2.7 GeV and isoscalar waves below 1.1 GeV. In the isoscalar case, the agreement is quite good, given the overall scatter of single-energy fits around the global result. More substantial differences are seen in the Saclay results for $`I=1`$ partial waves above 1 GeV.
A possible explanation for this discrepancy is given in the recent Saclay amplitude-reconstruction analysis. In Fig. 6, we compare the Saclay results to curves generated from our single-energy and energy-dependent solutions. (A similar comparison was made in the $`I=0`$ case, with good overall agreement between the three solutions.) Here, we are using the notation of Ref. and write the scattering matrix, $`M`$, as
$`M(\stackrel{}{k}_f,\stackrel{}{k}_i)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(a+b)+(ab)(\stackrel{}{\sigma }_1\stackrel{}{n})(\stackrel{}{\sigma }_2\stackrel{}{n})+(c+d)(\stackrel{}{\sigma }_1\stackrel{}{m})(\stackrel{}{\sigma }_2\stackrel{}{m})`$ (1)
$`+`$ $`(cd)(\stackrel{}{\sigma }_1\stackrel{}{l})(\stackrel{}{\sigma }_2\stackrel{}{l})+e(\stackrel{}{\sigma }_1+\stackrel{}{\sigma }_2)\stackrel{}{n}],`$ (2)
where $`\stackrel{}{k}_f`$ and $`\stackrel{}{k}_i`$ are the scattered and incident momenta in the center-of-mass system, and
$`\stackrel{}{n}={\displaystyle \frac{\stackrel{}{k}_i\times \stackrel{}{k}_f}{|\stackrel{}{k}_i\times \stackrel{}{k}_f|}},\stackrel{}{l}={\displaystyle \frac{\stackrel{}{k}_i+\stackrel{}{k}_f}{|\stackrel{}{k}_i+\stackrel{}{k}_f|}},\stackrel{}{m}={\displaystyle \frac{\stackrel{}{k}_f\stackrel{}{k}_i}{|\stackrel{}{k}_f\stackrel{}{k}_i|}}.`$In Refs. , multiple solutions were found at most energy-angle points. The $`pp`$ amplitudes are plotted together in Fig. 6 where we can see that, in some cases, our single-energy results favor one branch, while the energy-dependent fit follows another. This feature was also evident in Ref. , where it was shown that the Saclay partial-wave analyses followed a branch different from our preliminary energy-dependent fits.
Finally, in Fig. 7, we return to the low-energy region which has been controversial mainly due to recent determinations of $`ϵ_1`$. In Ref. , the trend of recent determinations was used to argue for an $`NN`$ tensor interaction stronger than predicted by all meson-exchange-based potential models, and in conflict with values found in both the Nijmegen and VPI partial-wave analyses. This trend is absent in our figure, where we have compared two of our energy-dependent fits, SP00 and SP40, and the Nijmegen potential, to a selection of recent single-energy fits. Clearly there is considerable scatter in the single-energy results. However, given this variation, the overall agreement with energy-dependent fits is quite good.
## IV Conclusions and Future Prospects
In our previous analysis , an extension of the energy range for $`pp`$ elastic scattering, from 1.6 to 2.5 GeV, was mainly motivated by the addition of precise new (unpolarized) cross section measurements from the EDDA collaboration . Given the sparse polarization data in this region, the fit was expected to change significantly with the addition of Saclay $``$ and future COSY (polarized) measurements. It is therefore somewhat surprising how little our new solution (SP00) has changed from SM97 .
We have seen that the $`I=0`$ amplitudes are generally in good agreement with the Saclay results. This holds true in the low-energy region as well, the Saclay value for $`ϵ_1`$ being consistent with our result. However, the agreement for $`pp`$ ($`I=1`$) amplitudes above 1 GeV is less impressive.
As mentioned above, this difference in partial-wave solutions may be a reflection of the non-uniqueness seen in the Saclay amplitude reconstruction. The selection of data included in these analyses could also be a factor. Clearly, this serves as further motivation for the polarization measurements being performed at COSY and JINR .
###### Acknowledgements.
The authors express their gratitude to J. Ball, Z. Dolezal, F. Dohrmann, W. Haeberli, H. O. Klages, C. Lechanoine-Leluc, F. Lehar, B. Lorentz, H.-O. Meyer, N. Olsson, B. von Przewoski, H. M. Spinka, W. Tornow, S. W. Wissink, and K. Yasuda for providing experimental data prior to publication or for clarification of information already published. Also, we thank F. Lehar for fruitful discussions and PWA predictions. This work was supported in part by the U. S. Department of Energy Grant DE–FG02–99ER41110. The authors gratefully acknowledges a contract from Jefferson Lab under which this work was done. Jefferson Lab is operated by the Southeastern Universities Research Association under the U. S. Department of Energy Contract DE–AC05–84ER40150.
FIGURE CAPTIONS
* Energy-angle distribution of recent (post-1997) (a) $`pp`$ and (b) $`np`$ data. The $`pp`$ data contribution below 3000 MeV is 30% and most new data are A<sub>y</sub> (47%) or A<sub>yy</sub> (29%). The $`np`$ data contribution below 1300 MeV is 6% and most new data are A<sub>y</sub> (24%). Total cross sections are plotted at zero degrees.
* Isovector partial-wave amplitudes from 0 to 3 GeV in the proton kinetic energy. Solid (dashed) curves give the real (imaginary) parts of amplitudes corresponding to the SP00 solution. The real (imaginary) parts of single-energy solutions are plotted as filled (open) circles. The SM97 solution is plotted with long dash-dotted (real part) and short dash-dotted (imaginary part) lines. FA91 solution is shown by dashed lines. The dotted curve gives the unitarity limit Im$`T`$ \- $`T^2`$ \- $`T_{sf}^2`$ from SP00, where $`T_{sf}`$ is the spin-flip amplitude. All amplitudes are dimensionless.
* Isoscalar partial-wave amplitudes from 0 to 1.2 GeV. Notation as in Fig. 2.
* Phase-shift parameters for isovector partial-wave amplitudes from 0 to 3000 MeV. The SP00 and SM97 solutions are plotted as solid and dash-dotted curves, respectively. Our single-energy solutions and those from Saclay ) are given by filled and open circles, respectively.
* Phase-shift parameters for isoscalar partial-wave amplitudes from 0 to 1200 MeV. Notation as in Fig. 6.
* Direct-reconstruction scattering amplitudes at (a) 1.80 GeV, (b) 2.10 GeV, (c) 2.40 GeV, and (d) 2.70 GeV. The real (imaginary) parts of amplitudes $`a`$ to $`e`$ are shown in $`\sqrt{mb/sr}`$ as a function of the c. m. scattering angle and plotted as filled (open) circles. Our SP00 (single-energy) solution is plotted with solid (dashed) lines.
* Summary of analyses giving $`ϵ_1`$ in the energy range up to 80 MeV. The solid (dashed) curve gives our SP00 (SP40) PWA results. Nijmegen potential results are plotted as a dash-dotted line. Filled circles (diamonds) give GW (Saclay ) single-energy PWA results. Open squares denote the single-energy PWA from PSI . Other results are from TUNL (star) , Bonn (filled box) , Prague (open diamond) , Erlangen (open circle) , PSI (filled triangle) , and Karlsruhe (open triangle) .
Table I. Comparison of present (SP00 and SP40) and previous (SM97, SM94, VZ40, FA91, SM86, and SP82) energy-dependent partial-wave analyses. The $`\chi ^2`$ values for the previous solutions correspond to our published results (, , and $``$ ).
Solution Range (MeV) $`\chi ^2`$/$`pp`$ data Range (MeV) $`\chi ^2`$/$`np`$ data Ref. SP00 $`03000`$ 36617/21796 $`01300`$ 18693/11472 Present SP00 $`(02500)`$ 34277/20947 $`01300`$ 17693/11330 Present SP00 $`(01600)`$ 23927/15766 $`01300`$ 17693/11330 Present SP00 $`(0400)`$ 4593/ 3454 $`(0400)`$ 5371/ 3831 Present SP40 $`0400`$ 4398/ 3454 $`0400`$ 5415/ 3831 Present SM97 $`02500`$ 28686/16994 $`01300`$ 17437/10854 SM94 $`01600`$ 22371/12838 $`01300`$ 17516/10918 VZ40 $`0400`$ 3098/2170 $`0400`$ 4595/3367 FA91 $`01600`$ 20600/11880 $`01100`$ 13711/7572 SM86 $`01200`$ 11900/7223 $`01100`$ 8871/5474 SP82 $`01200`$ 9199/5207 $`01100`$ 9103/5283
Table II. Comparison of the single-energy (SES) and energy-dependent (SP00) fits to $`pp`$ and $`np`$ data. Values of $`\chi ^2`$ are given for the single-energy and SP00 fits (evaluated over the same energy bins). Also listed is the number of parameters varied in each single-energy solution.
Energy Range (MeV) $`\chi ^2`$ SES(SP00)/$`pp`$ data $`\chi ^2`$ SES(SP00)/$`np`$ data Parameters 4-6 22(39)/28 78(83)/63 6 7-12 84(134)/88 254(333)/101 6 11-19 17(47)/27 205(455)/247 8 19-30 123(268)/114 292(321)/316 8 32-67 282(354)/224 809(879)/548 10 60-90 48(63)/72 514(629)/355 10 80-120 152(156)/154 465(453)/382 10 125-174 313(336)/287 603(653)/333 11 175-225 494(542)/435 701(734)/504 13 225-270 222(246)/228 299(345)/278 13 276-325 771(802)/740 628(680)/564 17 325-375 460(474)/406 416(460)/353 17 375-425 738(758)/607 804(870)/599 17 425-475 1055(1156)/803 828(870)/682 18 475-525 1311(1565)/1081 1248(1404)/787 30 525-575 858(956)/754 672(694)/488 31 575-625 1039(1112)/760 423(484)/367 34 625-675 908(842)/773 1270(1611)/873 36 675-725 860(923)/797 404(468)/386 37 725-775 1007(1311)/827 518(556)/381 37 775-824 1690(1840)/1301 1550(1861)/948 38 827-874 1155(1330)/914 388(467)/365 39 876-924 342(475)/389 752(905)/625 41 926-974 762(992)/679 363(512)/353 43
Table II. (continued)
Energy Range (MeV) $`\chi ^2`$ SES(SP00)/$`pp`$ data $`\chi ^2`$ SES(SM97)/$`np`$ data Parameters 976-1020 917(1177)/708 284(425)/328 43 1078-1125 815(1128)/573 519(846)/427 47 1261-1299 691(1006)/505 $``$ 30 1481-1521 140(307)/149 $``$ 30 1590-1656 505(892)/460 $``$ 31 1685-1724 174(309)/116 $``$ 31 1778-1818 625(1097)/506 $``$ 33 1929-1975 377(463)/366 $``$ 33 2065-2120 1173(1938)/829 $``$ 33 2175-2225 1476(2046)/758 $``$ 33 2330-2470 1013(1808)/713 $``$ 33 2500-2600 250(523)/311 $``$ 33 2642-2714 302(1016)/307 $``$ 33 2792-2869 148(405)/153 $``$ 33
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# Statistical mechanics of the shallow water system
## 1 Introduction
Two-dimensional flows with high Reynolds numbers have the striking property of organizing spontaneously into large scale coherent vortices . These vortices are common features of geophysical and astrophysical flows with the well-known example of Jupiter’s Great Red Spot, a huge vortex persisting for more than three centuries in a turbulent shear between two zonal jets. These vortices share some common features with stellar systems like elliptical or spherical galaxies that form after a phase of “violent relaxation” . They can also have applications in the process of planet formation which may have begun inside persistent gaseous vortices born out of the protoplanetary nebula. Understanding and predicting the structure of these organized states is still a challenging problem.
Since the dynamics of these systems is highly nonlinear, a deterministic description of the flow for late times is impossible and one must recourse to statistical methods. The statistical mechanics of 2D flows was first considered by Onsager (1949), followed by Joyce & Montgomery (1973), in the point vortex approximation. The more realistic case of continuous vorticity fields was later on treated by Kuzmin (1982), Miller (1990) and Robert & Sommeria (1991).
The statistical theory of the Euler equation provides a systematic framework to tackle the problem of self-organization in 2D flows. At a microscopic level, complex, inviscid, deformation of the vorticity field creates an intricate filamentation; however, if we introduce a macroscopic level of description (a “coarse-graining”) it can be shown that an overwhelming majority of these microscopic configurations are close to a macroscopic state (the statistical equilibrium or Gibbs state) obtained by maximizing a “mixing entropy” while accounting for all the constraints of the Euler equation (the conservation of energy and of the detailed distribution of vorticity). The maximum entropy theory generally predicts a nonlinear relationship between vorticity and stream function which respects the properties of the inviscid dynamics.
The predictions of the statistical theory ($`\nu =0`$) have been tested in a large number of numerical simulations and fluid laboratory experiments at high Reynolds numbers ($`\nu 0`$). Good agreement is obtained in many cases, as long as stirring is sufficiently intense to lead to equilibrium before significant influence of viscous effects .
However the Euler equation itself is of limited scope for applications to atmospheric or oceanic motion, where the Coriolis force and density stratification have a strong influence. A first step in this direction is provided by the quasi-geostrophic model. The application of the statistical theory to this case is straightforward as the flow is still assumed non-divergent, and the vorticity is just replaced by a potential vorticity. This has been discussed in several papers .
We here extend the theory to the more general case of shallow water system, a “compressible” 2D flow, whose properties are recalled in section 2. Several numerical computations, for instance , indicate inertial organisation of vortical motion into coherent vortices, like with the incompressible Euler equations. In these flows dominated by vortical motion, the influence of density waves is weak, and the free surface tends to be controled by the vortical motion through a balance condition generalizing the quasi-geostrophic balance.
In section 3 the procedure of Robert and Sommeria (1991) is extended to the shallow water system with discussion of simplified cases in section 4. We still assume that the vorticity field creates intricate filamentation but the divergence and water height (surface density) fields are still smooth. These assumptions are justified for flows dominated by vortical motion at moderate Mach numbers, for which the generation of shocks is not effective. The relaxation toward equilibrium is discussed in section 5 providing practical methods for determining the statistical equilibrium as well as sub-grid scale modeling of turbulence in a shallow water system. Finally the case of particular geometries is discussed in section 6.
## 2 The shallow water equations
Consider a fluid layer with thickness $`h(x,y,t)`$ submitted to a gravity $`g`$ on a rotating planet. We assume that the layer is thin with respect to the characteristic length scale of the horizontal motion. In that case, the velocity field $`𝐮(x,y,t)`$ can be assumed two-dimensional and the vorticity $`\omega \omega \omega =\omega 𝐞_𝐳=𝐮`$ is directed along the vertical axis. We shall assume for simplicity a plane geometry, with rotation vector $`\mathrm{\Omega }`$$`\mathrm{\Omega }`$$`\mathrm{\Omega }`$ directed along the vertical, but extension to a spherical geometry would be straigthforward (we introduce the Coriolis effect but no centrifugal force as the latter is incorporate in the gravity of the planet). The time evolution of these quantities is determined by the shallow water equations (see, e.g., ):
$$\frac{h}{t}+(h𝐮)=0$$
(1)
$$\frac{𝐮}{t}+(\omega \omega \omega +2\mathrm{\Omega }\mathrm{\Omega }\mathrm{\Omega })𝐮=B$$
(2)
Here the usual advective term $`𝐮.𝐮\text{ }`$has been expressed using the well-known identity of vector analysis $`𝐮.𝐮=(𝐮^2/2)+\omega \omega \omega 𝐮`$, and the term $`𝐮^2/2\text{ }`$incorporated in the Bernouilli function
$$B=gh+\frac{𝐮^2}{2},$$
(3)
together with the pressure term $`gh`$. Note that the shallow water system can be viewed as a 2D flow of a compressible gas with density $`h`$ and equation of state $`p=gh^2/2`$, and our results could be readily generalized to 2D compressible adiabatic flows. We shall often refer to the shallow water system as the “compressible case”, by opposition with the “incompressible” case, for which (1) is replaced by $`.𝐮=0`$ .
One can easily check that the potential vorticity (PV)
$$q=\frac{\omega +2\mathrm{\Omega }}{h}$$
(4)
is conserved for each fluid parcel, i.e:
$$\frac{dq}{dt}\frac{q}{t}+𝐮q=0$$
(5)
Each mass element $`hd^2𝐫`$ is conserved during the course of the evolution. Together with (5), this implies the conservation of the Casimirs
$$C_f=f(q)hd^2𝐫$$
(6)
where $`f`$ is any continuous function of the potential vorticity. In particular, the moments $`\mathrm{\Gamma }_n`$ of the potential vorticity are conserved
$$\mathrm{\Gamma }_n=q^nhd^2𝐫$$
(7)
The moments $`n=0,1,2`$ are respectively the total mass $`M`$, the circulation $`\mathrm{\Gamma }`$ and the PV enstrophy $`\mathrm{\Gamma }_2`$. The energy
$$E=h\frac{𝐮^2}{2}d^2𝐫+\frac{1}{2}gh^2d^2𝐫$$
(8)
involving a kinetic and potential part, is also a conserved quantity. Note finally that for a multiply connected domain, like the annulus or channel discussed in section 6, the circulation along each boundary is conserved, in addition to $`\mathrm{\Gamma }`$ .
It will be convenient in the sequel to use a Helmholtz decomposition of the momentum $`h𝐮`$ into a purely rotational and a purely divergent part
$$h𝐮=𝐞_z\psi +\varphi $$
(9)
where $`\psi `$ and $`\varphi `$ are defined as solutions of the Poisson equations
$$\mathrm{\Delta }\psi =(h𝐮),\psi =const.\text{on each boundary}$$
(10)
$$\mathrm{\Delta }\varphi =(h𝐮),\varphi /\zeta =0\text{on each boundary}$$
(11)
where the conditions at the domain boundary (with normal coordinate $`\zeta \text{ }`$) are the consequences of the impermeability condition. We here consider a domain with a single (outer) boundary, so we can take $`\psi =0`$ at this boundary (as $`\psi `$ is defined within an arbitrary gauge constant).
For a steady solution, the mass conservation (1) reduces to $`(h𝐮)=0`$, so that $`\varphi =0`$ and
$$h𝐮=\text{e}_z\psi (\text{steady}),$$
(12)
and equation (5) reduces to $`𝐮q=0`$, which implies that $`q\text{ }`$is a function $`F\text{ }`$of the stream function $`\psi `$
$$q=F(\psi )$$
(13)
Finally, equation (2) reduces to
$$(\omega \omega \omega +2\mathrm{\Omega }\mathrm{\Omega }\mathrm{\Omega })𝐮=B$$
(14)
Taking the dot product with $`𝐮`$, we obtain $`𝐮B=0`$ or, equivalently, $`B=B(\psi )`$. This is known as the Bernouilli theorem. Substituting for (12) in equation $`(\text{14})`$, we obtain a specific relationship between the potential vorticity $`q\text{ }`$and the Bernouilli function $`B\text{ }`$in the form:
$$q=\frac{dB}{d\psi }$$
(15)
Small perturbations to a state of rest, with uniform thickness $`H\text{ }`$, satisfy linearized equations with two branches of solutions. For small scales, these are the usual surface waves on one hand, with purely divergent velocity and propagation speed $`c=(gH)^{1/2}`$, and steady vortical divergenceless modes on the other hand. In nonlinear regimes, these two modes interact. Vortical motion with scale $`l\text{ and }`$typical vorticity $`\omega `$ fluctuates on time scale $`\omega ^{1\text{ }}`$, so it emits waves at wavelength $`\lambda c/\omega `$, hence $`\lambda /l`$ is the inverse of the Mach number $`c/u`$ based on the local velocity $`u\omega l`$. Therefore we expect that for the case of small Mach numbers that we shall consider, velocity divergence and free surface deformation are much smoother than the vorticity field (their wavelength is much larger).
## 3 The maximum entropy theory
### 3.1 General principles and notations
For slow velocities, the shallow water system reduces to the quasi-geostrophic (QG) equations, with $`hcte\text{ }`$, such that (1) reduces to the incompressibility condition $`.𝐮=0\text{ }`$. Then the velocity field remains regular for any time, but it generally develops complex fine scale vorticity filaments so that a deterministic description of the flow would require a rapidly increasing amount of information as time goes on. The idea of the statistical theory is to give up *such a* deterministic description and refer to a probabilistic description. Therefore, the exact knowledge of the “fine-grained”, or microscopic potential vorticity field is replaced by the probability density (area fraction) $`\rho (𝐫,\sigma )`$ of finding the potential vorticity level $`\sigma `$ at position $`𝐫`$.
For the more general shallow water system, the inviscid dynamics generally leads to singularities (shocks), with associated energy dissipation (even in the absence of viscosity). This is a source of fundamental mathematical difficulty for the generalization of the equilibrium statistical mechanics initially developed for the Euler equations or QG system. Nevertheless, for the case of small Mach numbers that we consider, shocks occur only after a very long time (due to non-linear steepening of surface waves), and they involve a weak energy dissipation, since most of the energy remains in the vortical motion. Furthermore fine scale vorticity fluctuations behave like in the QG system, and only interact with surface waves and flow divergence at much larger scale, as discussed above. We shall therefore assume that vorticity fluctuates at small scale, but divergence is smooth as well as the height $`h`$.
We still describe the local vorticity fluctuations by the probability $`\rho (𝐫,\sigma )`$ of finding the potential vorticity level $`\sigma `$ in a small neighborhood of the position $`𝐫`$. The normalization condition yields at each point
$$\rho (𝐫,\sigma )𝑑\sigma =1$$
(16)
The locally averaged field of potential vorticity is expressed in terms of the probability density in the form
$$\overline{q}=\rho (𝐫,\sigma )\sigma 𝑑\sigma $$
(17)
This locally averaged field is called the macroscopic, or coarse-grained, potential vorticity. A macroscopic state is fully defined by $`\rho (𝐫,\sigma )`$, the height field $`h(𝐫)`$ and the flow divergence (assumed without microscopic fluctuations). The velocity field $`𝐮(𝐫)`$ can be also considered as smooth, as it integrates the vorticity fluctuations, and can be deduced by integration from its divergence and vorticity $`\overline{\omega }=\overline{q}h2\mathrm{\Omega }`$ . The energy (8) depends only on this smooth field, with negligible influence of local fluctuations, like in the incompressible case .
The conservation of the Casimirs (6) is equivalent to the conservation of the global distribution of potential vorticity (i.e the total area of each level of potential vorticity ponderated by $`h`$):
$$\gamma (\sigma )=\rho (𝐫,\sigma )h(𝐫)d^2𝐫$$
(18)
The microscopic moments of potential vorticity can be written:
$$\mathrm{\Gamma }_n=\gamma (\sigma )\sigma ^n𝑑\sigma =\overline{q^n}h(𝐫)d^2𝐫$$
(19)
where
$$\overline{q^n}=\rho (𝐫,\sigma )\sigma ^n𝑑\sigma $$
(20)
and are conserved during an inviscid evolution. Note that for $`n2`$, the macroscopic moments of potential vorticity $`\mathrm{\Gamma }_n^{c.g.}=\overline{q}^nh(𝐫)d^2𝐫`$ are not conserved, as there are partly transfered into fine-grained fluctuations.
The mixing entropy
$$S=\rho (𝐫,\sigma )\mathrm{ln}\rho (𝐫,\sigma )h(𝐫)d^2𝐫𝑑\sigma $$
(21)
measures the number of microscopic configurations associated with the same macroscopic field of potential vorticity. The dependence in $`\rho `$ is the same as for the incompressible case , and the factor $`h(𝐫)`$ is introduced to insure that entropy is conserved by mere macroscopic displacement of fluid parcels. Indeed the mass element $`h(𝐫)d^2𝐫`$ is conserved by fluid particle displacement instead of the surface element $`d^2𝐫`$ in the incompressible case. At equilibrium, the system is expected to be in the most probable (i.e. most mixed) state consistent with the constraints of the Euler equation. The entropy (21) has been justified by rigorous mathematical arguments (in the incompressible case) but other forms have been recently proposed (see discussion in ). Therefore, we shall consider a general form of entropy
$$S=s(\rho (𝐫,\sigma ))h(𝐫)d^2𝐫𝑑\sigma $$
(22)
and find that equation (21) is the only expression leading to a well defined entropy extremum.
### 3.2 First order variations
According to the previous discussion, the most probable macroscopic state is obtained by maximizing the entropy (22) with fixed energy (8), global vorticity distribution (18) and local normalization (16). This problem is treated by introducing Lagrange multipliers, so that the first variations satisfy
$$\delta S\beta \delta E\alpha (\sigma )\delta \gamma (\sigma )𝑑\sigma \zeta (𝐫)\delta \left(\rho (𝐫,\sigma )𝑑\sigma \right)hd^2𝐫=0.$$
(23)
The Lagrange multipliers are respectively the “inverse temperature” $`\beta ,`$ the “chemical potential” $`\alpha (\sigma )`$ of species $`\sigma ,`$ and $`\zeta (𝐫)`$.
We shall take $`h`$, $`\rho h`$ and $`𝐮`$ as independent variables characterizing the macroscopic state. Then, it is easy to establish, by differentiating respectively (21), (18) and (16), that:
$$\delta S=[\rho s^{}(\rho )s(\rho )]\delta hd^2𝐫𝑑\sigma s^{}(\rho )\delta (\rho h)d^2𝐫𝑑\sigma $$
(24)
$$\delta \gamma (\sigma )=\delta (\rho h)d^2𝐫$$
(25)
$$h\delta \left(\rho (𝐫,\sigma )𝑑\sigma \right)=\delta (\rho h)𝑑\sigma \rho \delta h𝑑\sigma $$
(26)
The variations of the energy are conveniently expressed in terms of the Bernouilli function $`B`$,
$$\delta E=B\delta hd^2𝐫+h𝐮\delta 𝐮d^2𝐫$$
(27)
Then, using the Helmholtz decomposition (9) for the momentum $`h𝐮`$ , the second integral can be rewritten
$$h𝐮\delta 𝐮d^2𝐫=(\psi \delta 𝐮)𝐳d^2𝐫+\varphi \delta 𝐮d^2𝐫$$
(28)
Integrating by parts with the identities $`(\psi \delta 𝐮)=\psi (\delta 𝐮)+\psi \delta 𝐮`$ and $`(\varphi \delta 𝐮)=\varphi (\delta 𝐮)+\varphi \delta 𝐮`$, and using the boundary conditions for $`\psi `$ and $`\varphi `$ , we obtain
$$h𝐮\delta 𝐮d^2𝐫=\psi \delta \overline{\omega }d^2𝐫\varphi (\delta 𝐮)d^2𝐫$$
(29)
Using (16) (4) and (17), we have finally
$$\delta E=B\rho \delta hd^2𝐫𝑑\sigma +\psi \sigma \delta (\rho h)d^2𝐫𝑑\sigma \varphi \delta (𝐮)d^2𝐫$$
(30)
The variation (23) vanishes for any small changes of the variables only if the coefficient of each independent variable vanishes:
$$\delta (\rho h):s^{}(\rho )=\beta \sigma \psi \alpha (\sigma )\zeta (𝐫)$$
(31)
$$\delta h:s^{}(\rho )\frac{s(\rho )}{\rho }=\beta B\zeta (𝐫)$$
(32)
$$\delta (𝐮):\varphi =0$$
(33)
Note that the right hand side of equation (32) is independent of $`\sigma `$. This implies that the term on the left hand side must be a constant (that we can take equal to $`1`$ without loss of generality):$`s^{}(\rho )s(\rho )/\rho =1`$. This equation is easily integrated in $`s(\rho )=A\rho +\rho \mathrm{ln}\rho `$ where $`A`$ is an integration constant. When substituted in equation (22), using (19), this yields
$$S=\rho (𝐫,\sigma )\mathrm{ln}\rho (𝐫,\sigma )h(𝐫)d^2𝐫𝑑\sigma AM$$
(34)
which is just the entropy (21) up to an additive constant term $`AM`$ (which we can take equal to zero without loss of generality). Therefore, the entropy (21) is the only functional of the form (22) for which the maximization problem has a solution. This result is astounding because it is obtained without any explicit reference to thermodynamical arguments.
### 3.3 The Gibbs states
Substituting explicitely $`s(\rho )=\rho \mathrm{ln}\rho `$ in equation (31), the optimal probability density can be expressed as
$$\rho (𝐫,\sigma )=\frac{1}{Z(\psi )}g(\sigma )e^{\beta \sigma \psi }$$
(35)
where $`Z(\psi )e^{\zeta (𝐫)+1}`$ and $`g(\sigma )e^{\alpha (\sigma )}`$. The normalization condition (16) leads to a value of the partition function $`Z`$ of the form
$$Z=g(\sigma )e^{\beta \sigma \psi }𝑑\sigma $$
(36)
and the locally averaged potential vorticity (17) is expressed as a function of $`\psi `$ according to:
$$\overline{q}=\frac{g(\sigma )\sigma e^{\beta \sigma \psi }𝑑\sigma }{g(\sigma )e^{\beta \sigma \psi }𝑑\sigma }=F(\psi )$$
(37)
This can be rewritten
$$\overline{q}=\frac{1}{\beta }\frac{d\mathrm{ln}Z}{d\psi }$$
(38)
This is the same functional relation as in the case of 2D incompressible Euler flows .
Differentiating equation (37) with respect to $`\psi `$, we check that the variance of the potential vorticity can be written
$$q_2\overline{q^2}\overline{q}^2=\frac{1}{\beta }F^{}(\psi )$$
(39)
or, alternatively (see equation (38)):
$$q_2=\frac{1}{\beta ^2}\frac{d^2\mathrm{ln}Z}{d\psi ^2}$$
(40)
Therefore the slope of the function $`\overline{q}=F(\psi )`$ is directly related to the variance of the vorticity distribution. Relation (39) has the same form and origin as the “fluctuation-dissipation” theorem in statistical field theory, where $`d\overline{q}/d\psi `$ is interpreted as a susceptibility. Since $`q_2>0`$, we find that the function $`\overline{q}=F(\psi )`$ is monotonic; it is decreasing for $`\beta >0`$ and increasing for $`\beta <0`$ (it is constant for $`\beta =0`$). Another proof of this result is given in .
Substituting explicitely $`s^{}(\rho )s(\rho )/\rho =1`$ in equation (32), we have
$$B=\frac{1}{\beta }\mathrm{ln}Z$$
(41)
This relation shows that the Bernouilli function plays the role of a free energy in the statistical theory. We note that both $`B`$ and $`\overline{q}`$ are functions of $`\psi `$, while $`\varphi =0`$ , from (33), as it should for steady flows. Furthermore, taking the derivative of (41) with respect to $`\psi `$ and using (38), we check that the relation $`\overline{q}=dB/d\psi `$ required for a steady solution of the shallow water equation is satisfied. Therefore the flow selected by the purely statistical entropy maximization procedure does not evolve anymore by the flow evolution, so the statistical theory is indeed consistent with the dynamics.
## 4 Properties of the Gibbs states in some particular cases:
### 4.1 Particular $`\overline{q}\psi `$ relationships:
The Gibbs states are characterized by the relation (37) between $`\overline{q}`$ and $`\psi `$, which is always monotonic, as shown in the previous section. It is determined by the conservation laws, but only indirectly through the set of Lagrange multipliers $`\beta `$ and $`\alpha (\sigma )`$. In practice we need to discretize the PV levels, and keeping only two levels, $`q=\sigma _0`$ and $`q=\sigma _1`$ is often representative of more general cases. Then the probability distribution $`\rho `$ just depends on a single probability $`p_1`$ of finding the level $`\sigma _1`$ (for instance), with the probability $`1p_1`$ of finding the complementary level $`\sigma _0`$ , i.e. $`g(\sigma )`$ is the sum of two Dirac function terms, $`g(\sigma )=g_1[\lambda \delta (\sigma \sigma _0)+\delta (\sigma \sigma _1)]`$. This probability $`p_1`$ is directly related to the PV average by $`\overline{q}=p_1\sigma _1+(1p_1)\sigma _0`$, or reversely $`p_1=(\overline{q}\sigma _0)/(\sigma _1\sigma _0)`$. The expression (37) for $`\overline{q}`$ reduces to
$$\overline{q}=\sigma _0+\frac{(\sigma _1\sigma _0)}{1+\lambda e^{\beta (\sigma _1\sigma _0)\psi }}$$
(42)
This relation corresponds to the Fermi-Dirac statistics. The two unknown parameters $`\lambda `$ and $`\beta `$ are indirectly determined by the conserved quantities. The associated Bernouilli function (41) becomes
$$B=\frac{1}{\beta }\mathrm{ln}g_1\sigma _0\psi +\frac{1}{\beta }\mathrm{ln}\left\{\lambda +e^{\beta (\sigma _0\sigma _1)\psi }\right\}$$
(43)
The problem is also greatly simplified in the alternative case for which $`g(\sigma )`$ is a Gaussian:
$$g(\sigma )=g_0e^{\frac{(\sigma \sigma _{})^2}{2\sigma _2}}.$$
(44)
Then the local probability distribution (35) is also a Gaussian, and the corresponding Bernouilli function (41) is
$$B=\frac{1}{\beta }\mathrm{ln}[g_0(2\pi \sigma _2)^{1/2}]+\frac{1}{2}\sigma _2\beta \psi ^2\sigma _{}\psi ,$$
(45)
corresponding to a linear relationship:
$$\overline{q}=\beta \sigma _2\psi +\sigma _{}$$
(46)
According to (39) the variance of the potential vorticity is then uniform, with value $`q_2=\sigma _2`$(more generally, all the even momenta of the Gaussian are related to $`\sigma _2`$ by $`\overline{(q\overline{q})^{2n}}=(2n1)!!\sigma _2^n`$ and the odd momenta cancel).
This Gaussian local probability distribution is obtained by maximizing the entropy (21), reducing the constraints of the global distribution $`\gamma (\sigma )`$ to its first moments $`\mathrm{\Gamma }_0M`$, $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$. This will be the true Gibbs state for a particular initial distribution $`\gamma (\sigma )`$ with higher order moments equal to the global moments of this simplified Gibbs state. A linear relationship between $`\overline{q}`$ and $`\psi `$ can also be obtained for any distribution $`\gamma (\sigma )`$ in the limit of strong mixing where $`\beta \sigma \psi 1`$, so that (37) can be linearized, as discussed in Chavanis & Sommeria (1996).
### 4.2 Unidirectional solutions
We consider here unidirectional solutions such that $`𝐮=u(y)𝐞_x`$. The equilibrium relation $`\overline{q}=dB/d\psi `$ then yields, multiplying each member by $`hu=d\psi /dy`$,
$$g\frac{dh}{dy}+\frac{2\mathrm{\Omega }}{h}\frac{d\psi }{dy}=0.$$
(47)
This condition of geostrophic balance can be readily integrated in
$$h=H\sqrt{1\frac{4\mathrm{\Omega }}{gH^2}\psi }$$
(48)
A second relation is provided by the expression of the Bernouilli function yielding
$$\frac{1}{2h^2}\left(\frac{d\psi }{dy}\right)^2=B(\psi )gh$$
(49)
Combined with the expression (48) of $`h`$ and the Gibbs state expression (41) of $`B(\psi )`$, this yields a first order ordinary differential equation for $`\psi .`$ This equation depends on the constant $`H`$ and Lagrange parameters, for instance $`g_1,`$ $`\beta `$ and $`\lambda `$ in the case with two PV levels. The constraints on total mass $`M`$, global mass of PV level $`\sigma _1`$ and total energy $`E`$, indirectly determine these parameters. Note that in reality this solution must be viewed in a x-wise translating frame of reference, as discussed in section 6, due to the additional conservation law for momentum.
### 4.3 Axisymmetric solutions
For axisymmetric solutions $`𝐮=u_\theta (r)𝐞_\theta `$ where $`(r,\theta )`$ are polar coordinates, and $`hu_\theta =d\psi /dr`$. Then Equation (47) is replaced by the cyclostrophic balance
$$gh\frac{dh}{dr}=\frac{1}{hr}\left(\frac{d\psi }{dr}\right)^22\mathrm{\Omega }\frac{d\psi }{dr}$$
(50)
When $`hu_\theta =d\psi /dr0`$ (cyclone), $`h`$ is an increasing function of $`r`$, so the vortex core is a trough. In the opposite case $`u_\theta 0`$ (anticyclone) , the vortex core is a bump in geostrophic regimes. However for large velocities (Rossby number larger than unity), the term $`\frac{u_\theta ^2}{r}`$ dominates, leading always to a trough.
The expression of the Bernouilli constant gives again the form (49), just replacing $`y`$ by the radius $`r`$ . Combining this relation with (50), one gets a couple of first order ordinary differential equations in the variables $`\psi `$ and $`h.`$ As in the unidirectional case, the solution depends on two constants of integration and the Lagrange parameters, which are $`g_1,`$ $`\beta `$ and $`\lambda `$ in the case with two PV levels. This solution must be viewed in general in a rotating frame of reference, due to the additional conservation of angular momentum , as discussed in section 6.
## 5 Relaxation equations
### 5.1 The Maximum Entropy Production Principle
Relaxation methods are convenient to compute the statistical equilibrium resulting from any initial condition. The aim is to increase the entropy in successive steps while keeping constant all the conserved quantities. Turkington & Whitaker (1996) have implemented a relaxation method to calculate the Gibbs states obtained with the Euler equations. Robert and Sommeria (1992) had previously proposed relaxation equations in the form of a parameterization of sub-grid scale eddies which drives the system toward statistical equilibrium by a continuous time evolution. Such relaxation equations can be used both as a realistic coarse resolution model of the turbulent evolution, and as a method of determination of the statistical equilibrium resulting from this evolution (see for a short review). We here generalize this approach to the shallow water system.
We first decompose the vorticity $`\omega `$ and velocity $`𝐮`$ into a mean and fluctuating part, namely $`\omega =\overline{\omega }+\stackrel{~}{\omega }`$, $`𝐮=\overline{𝐮}+\stackrel{~}{𝐮}`$, keeping $`h`$ smooth. Taking the local average of the shallow water equations (1)(2), we get
$$\frac{h}{t}+(h\overline{𝐮})=0$$
(51)
$$\frac{\overline{𝐮}}{t}+(\overline{\omega \omega \omega }+2\mathrm{\Omega }\mathrm{\Omega }\mathrm{\Omega })\overline{𝐮}=B𝐞_z𝐉_\omega $$
(52)
$$B=gh+\frac{\overline{𝐮}^2}{2}$$
(53)
where the current $`𝐉_\omega =\overline{\stackrel{~}{\omega }\stackrel{~}{𝐮}}`$ represents the correlations of the fine-grained fluctuations. Although we have neglected the fluctuation energy $`\stackrel{~}{𝐮}^2`$ in front of $`\overline{𝐮}^2`$ (as well as fluctuations of $`h`$), we keep the correlations $`𝐉_\omega =\overline{\stackrel{~}{\omega }\stackrel{~}{𝐮}}`$, which represent the PV transport by sub-grid-scale eddies. This assumption is justified since, denoting $`ϵ`$ the typical scale of vorticity fluctuations, we have $`\stackrel{~}{𝐮}^2ϵ^2\overline{\omega }^2`$ and $`\overline{\stackrel{~}{\omega }\stackrel{~}{𝐮}}ϵ\overline{\omega }^2\stackrel{~}{𝐮}^2`$ (while $`\stackrel{~}{\omega }\overline{\omega }`$).
We deduce an equation for the evolution of the potential vorticity (4), taking the curl of (52) and using (51),
$$\frac{}{t}(h\overline{q})+.(h\overline{q}\overline{𝐮})=.𝐉_\omega $$
(54)
This equation can be viewed as a local conservation law for the circulation $`\mathrm{\Gamma }=\overline{q}hd^2𝐫`$. We shall determine the unknown current $`𝐉_\omega `$ by the thermodynamic principle of Maximum Entropy Production (MEP) . For that purpose, we need to consider not only the locally average PV field $`\overline{q}`$, but the whole probability distribution $`\rho (𝐫,\sigma ,t)`$ now evolving with time $`t`$. The conservation of the global vorticity distribution $`\gamma (\sigma )=\rho hd^2𝐫`$ can be written in the local form
$$\frac{}{t}(h\rho )+.(h\rho \overline{𝐮})=.𝐉$$
(55)
where $`𝐉(𝐫,\sigma ,t)`$ is the (unknown) current associated with the level $`\sigma `$ of potential vorticity. Integrating equation (55) over all the PV levels $`\sigma `$, using (16), and comparing with (51), we find the constraint
$$𝐉(𝐫,\sigma ,t)𝑑\sigma =0$$
(56)
Multiplying (55) by $`\sigma `$, integrating over all the PV levels, using (17) and comparing with (54), we get
$$𝐉(𝐫,\sigma ,t)\sigma 𝑑\sigma =𝐉_\omega $$
(57)
We can express the time variation of the energy $`\dot{E}dE/dt`$ in terms of $`𝐉`$, using (8) and (52), leading to the energy conservation constraint
$$\dot{E}=𝐉\sigma h\overline{𝐮}_{}d^2𝐫𝑑\sigma =0,$$
(58)
where $`\overline{𝐮}_{}𝐞_z\overline{𝐮}`$ . Using (21) and (55), we similarly express the rate of entropy production as
$$\dot{S}=𝐉(\mathrm{ln}\rho )d^2𝐫d\sigma .$$
(59)
According to the Maximum Entropy Production Principle, we determine the current $`𝐉`$ which maximizes the rate of entropy production $`\dot{S}`$ respecting the constraints $`\dot{E}=0`$, (56) and $`\frac{J^2}{2\rho }𝑑\sigma C(𝐫,\sigma )`$. The last constraint expresses a bound (unknown) on the value of the diffusion current. Convexity arguments justify that this bound is always reached so that the inequality can be replaced by an equality. The corresponding condition on first variations can be written at each time $`t`$:
$$\delta \dot{S}\beta (t)\delta \dot{E}\zeta (𝐫,t)\delta \left(𝐉𝑑\sigma \right)d^2𝐫D^1(𝐫,t)\delta \left(\frac{J^2}{2\rho }\right)𝑑\sigma d^2𝐫=0$$
(60)
and leads to a current of the form
$$𝐉=D(𝐫,t)(\rho +\beta (t)\sigma \rho h\overline{𝐮}_{}\zeta (𝐫,t)\rho )$$
(61)
The Lagrange multiplier $`\zeta (𝐫,t)`$ is determined by the constraint (56), which leads to
$$𝐉=D(𝐫,t)\left[\rho +\beta (t)\rho (\sigma \overline{q})h\overline{𝐮}_{}\right]$$
(62)
This optimal current is similar to the one obtained in ordinary incompressible 2D turbulence except that the term $`h\overline{𝐮}_{}`$ now replaces $`\psi .`$ The impermeability condition imposes that the normal component of the velocity and of the current vanishes at the wall. We therefore have the boundary conditions
$$\text{n}\overline{.\text{u}}=0(\text{on }\text{each }\text{boundary})$$
(63)
$$\text{n}.\rho =\beta (t)\rho (\sigma \overline{q})h\text{n}\overline{𝐮}_{}(\text{on }\text{each }\text{boundary})$$
(64)
where n is a unit vector normal to the wall.
The diffusion coefficient $`D`$ is not determined by the MEPP as it depends on the unknown bound $`C`$ on the current. It can be determined by more systematic procedures inspired from kinetic theories of plasma physics like in for the incompressible case. For the purpose of reaching the Gibbs state, it can simply be chosen arbitrarily. However, we shall show below that $`D`$ must be positive so as to insure entropy increase.
The conservation of energy (58) at any time determines the evolution of the Lagrange multiplier $`\beta (t)`$ according to
$$\beta (t)=\frac{D\overline{q}h\overline{𝐮}_{}d^2𝐫}{D(\overline{q^2}\overline{q}^2)(h\overline{𝐮}_{})^2d^2𝐫}$$
(65)
We can now provide an explicit form for the vorticity current $`𝐉_\omega \text{ }`$to introduce back in the shallow water equation (52). Indeed, using (62) and (17), we find
$$𝐉_\omega =D\left[\overline{q}+\beta (t)(\overline{q^2}\overline{q}^2)h\overline{𝐮}_{}\right]$$
(66)
Substituting (66) in equation (52), we obtain
$$\frac{\overline{𝐮}}{t}+(\overline{\omega \omega \omega }+2\mathrm{\Omega }\mathrm{\Omega }\mathrm{\Omega })\overline{𝐮}=B+D\left[𝐞_z\overline{q}\beta (t)(\overline{q^2}\overline{q}^2)h\overline{𝐮}\right]$$
(67)
Since $`\beta (t)0`$ in relevant situations, the last term in (67) represents a forcing proportional to $`\overline{𝐮}`$ which compensates the diffusion caused by the term $`\widehat{𝐞}_z\overline{q}\mathrm{\Delta }\overline{𝐮}`$. This additional term depends on the local PV variance $`\overline{q^2}\overline{q}^2\text{ }`$, related to the probability distribution $`\rho `$, and we need to keep track of it by solving the probability transport equations (55) in addition to the modified shallow water system (51) and (67). This set of equations increases the entropy (at an optimal rate), while preserving all the conservation laws of the initial inviscid shallow water system. We now check that the steady solutions reached by these relaxation equations is indeed the Gibbs state.
### 5.2 Relaxation towards the statistical equilibrium
The entropy production (59) can be written
$$\dot{S}=\frac{𝐉}{\rho }(\rho +\beta \rho (\sigma \overline{q})h𝐮_{})d^2𝐫𝑑\sigma +\beta 𝐉(\sigma \overline{q})h𝐮_{}d^2𝐫𝑑\sigma $$
(68)
Using (56) and (58), the second integral is seen to cancel out. Substituting for (62) in the first integral, we have
$$\dot{S}=\frac{J^2}{D\rho }d^2𝐫𝑑\sigma $$
(69)
which is positive provided that $`D0`$ , and this is clearly a necessary condition to assure entropy increase in all cases. A stationary solution $`\dot{S}=0`$ is such that $`𝐉=0`$ yielding, together with (9),
$$(\mathrm{ln}\rho )+\beta (\sigma \overline{q})\psi =0$$
(70)
For any reference PV level $`\sigma _0`$, it writes
$$(\mathrm{ln}\rho _0)+\beta (\sigma _0\overline{q})\psi =0$$
(71)
Substracting (70) and (71), we obtain $`\mathrm{ln}(\rho /\rho _0)+\beta (\sigma \sigma _0)\psi =0`$, which is immediately integrated into
$$\rho (𝐫,\sigma )=\frac{1}{Z(𝐫)}g(\sigma )e^{\beta \sigma \psi }$$
(72)
where $`Z^1(𝐫)\rho _0(𝐫)e^{\beta \sigma _0\psi (𝐫)}`$ and $`g(\sigma )e^{A(\sigma )}`$, $`A(\sigma )`$ being a constant of integration. Therefore, entropy increases until the Gibbs state (35) is reached, with $`\beta =lim_t\mathrm{}\beta (t)`$.
### 5.3 Simplified cases:
In the case of two PV levels $`\sigma _0`$ and $`\sigma _1`$, the transport equation (55) for the probability $`p_1`$ is equivalent to the transport equation for $`\overline{q}`$ (since $`\overline{q}=\sigma _0+p_1(\sigma _1\sigma _0)`$), which is already obtained from the curl of the shallow water equation (67). Therefore the relaxation equations reduce to the modified shallow water system
$$\frac{h}{t}+(h𝐮)=0$$
(73)
$$\frac{𝐮}{t}+\overline{q}h𝐞_z𝐮=(\frac{𝐮^2}{2}+gh)+D[𝐞_z\overline{q}\beta (t)q_2h𝐮]$$
(74)
$$\overline{q}=\frac{(𝐮).𝐞_z+2\mathrm{\Omega }}{h},q_2=(\overline{q}\sigma _0)(\sigma _1\overline{q})$$
(75)
$$\beta (t)=\frac{Dh(𝐞_z𝐮)\overline{q}d^2𝐫}{Dq_2(𝐞_z𝐮)^2h^2d^2𝐫}$$
(76)
$$\text{n}.\overline{q}=\beta (t)q_2h\text{n}\text{u}_{}(\text{on }\text{each }\text{boundary})$$
(77)
$$\text{n}.\text{u}=0(\text{on }\text{each }\text{boundary})$$
(78)
where we have omitted the over-bar on $`𝐮`$, and the expression (75) of $`q_2=\overline{q^2}\overline{q}^2`$ is easily obtained for a probability distribution with two values $`\sigma _0`$ and $`\sigma _1`$. The numerical implementation of this system will lead to the two PV levels Gibbs state.
Stating $`q_2=cte`$ instead of the expression (75) yields the Gaussian Gibbs state with linear relationship between $`\overline{q}`$ and $`\psi `$. Then the coefficient $`q_2\beta `$ used in (74) is directly obtained from (76). This is sufficient for the purpose of finding the statistical equilibrium, but more refined relaxation models can be used as discussed by Kazantsev et al. (1998) in the context of QG models.
### 5.4 The incompressible limit:
The case of ordinary 2D incompressible turbulence is recovered in the limit $`h1`$, $`q\omega `$ and $`\text{u}=\text{e}_z\psi `$. The relaxation equation (67) then becomes
$$\frac{\overline{𝐮}}{t}+(\overline{𝐮}.)\overline{𝐮}=\frac{1}{\rho }p+D(\mathrm{\Delta }\overline{𝐮}\beta (t)\omega _2\overline{𝐮})$$
(79)
where we have used the well-known identity of vector analysis $`\mathrm{\Delta }\overline{𝐮}=(\overline{𝐮})(\overline{𝐮})`$ which reduces for a two-dimensional incompressible flow to $`\mathrm{\Delta }\overline{𝐮}=𝐞_z\overline{\omega }`$. Equation (79) is valid even if $`D`$ is space dependant unlike with a usual viscosity term. In previous publications this equation was given only in its vorticity form
$$\frac{\overline{\omega }}{t}+(\overline{\omega }\overline{𝐮})=\left[D\left(\overline{\omega }+\beta (t)\omega _2\psi \right)\right]$$
(80)
and the equivalence with (79) is not obvious at first sights when $`D`$ is space dependent. At equilibrium, we have from (79) the identity
$$\mathrm{\Delta }\overline{𝐮}=\beta \omega _2\overline{𝐮}$$
(81)
which can be deduced directly from the Gibbs state. Indeed for a stationary solution $`\overline{\omega }=F(\psi )`$, the previous identity $`\mathrm{\Delta }\overline{\text{u}}=\text{e}_z\overline{\omega }`$ becomes $`\mathrm{\Delta }\overline{𝐮}=F^{}(\psi )\overline{𝐮}`$ equivalent to (81) for a Gibbs state thanks to (39).
We now account for a deformation of the fluid layer but assume that the elevation with respect to the average thickness $`H`$ is weak, so that
$$h=H(1+\eta )\text{with}\eta 1$$
(82)
To first order the flow is incompressible and equation (1) reduces to $`.𝐮=0`$, or equivalently $`𝐮=𝐞_z\psi `$ (there is a factor $`H`$ with the previous definition (9)). In the quasi-geostrophic limit of small Rossby numbers $`\omega \mathrm{\Omega }`$, the momentum equation (2) reduces at zero order to the geostrophic balance
$$𝐮=\frac{gH}{2\mathrm{\Omega }}𝐞_z\eta \text{or}\psi =\frac{gH^2}{2\mathrm{\Omega }}\eta $$
(83)
The expression for the potential vorticity then reduces to
$$\zeta Hq2\mathrm{\Omega }\omega +\frac{\psi }{L_R^2}$$
(84)
with the Rossby radius of deformation
$$L_R=\frac{\sqrt{gH}}{2\mathrm{\Omega }}$$
(85)
The term $`\frac{1}{L_R}\psi `$ in (84) creates a shielding of the interaction between vortices (similar to the Debye shielding in plasma physics) on a length scale $`L_R`$. In the limit $`1/L_R0`$, we recover the 2D incompressible equations. For finite $`L_R`$ we can extend the statistical theory of the 2D Euler equations to the QG case by simply replacing the vorticity $`\overline{\omega }`$ by the potential vorticity $`\overline{\zeta }`$ .
## 6 The case of circular domains or channel:
### 6.1 Statistical equilibrium
In a disk, the angular momentum
$$L=h(𝐫𝐮)_zd^2𝐫$$
(86)
is conserved. This additional constraint can be accounted for by adding a term $`\beta \lambda \delta L`$ in the first order variation (23). We can write $`\delta L=\delta h(𝐫𝐮)_zd^2𝐫+h(𝐞_z𝐫).\delta 𝐮d^2𝐫`$, and the second term can be expressed in terms of $`\delta \overline{\omega }`$ and $`\delta (.𝐮)`$ by a Helmholtz decomposition of $`h(𝐞_z𝐫)`$ analogous to (9), followed by an integration by part. This is analogous to the formulae (28)(29) used for expressing the energy variation. We can combine the energy and momentum variations by defining
$$h[𝐮\lambda (𝐞_z𝐫)]=𝐞_z\psi ^{}+\varphi ^{}$$
(87)
instead of (9). Adding the new terms in (31)(33) yields the Gibbs state (35)(41) for the velocity $`𝐮^{}=𝐮\lambda (𝐞_𝐳𝐫)`$ seen in the reference frame rotating at angular velocity $`\lambda `$. Accordingly, the expression of the Bernouilli function must be modified by a term of centrifugal force: we must use $`B^{}(\psi ^{})=gh+\frac{𝐮_{}^{}{}_{}{}^{2}}{2}\lambda ^2r^2`$ instead of $`B(\psi )`$. We find therefore that the Gibbs state (its locally averaged velocity field) is a solution of the shallow water equation which is steady in a reference frame rotating at a modified angular velocity $`\mathrm{\Omega }+\lambda `$. This velocity is indirectly determined by the constraint on angular momentum. Note that the result can be readily extended to the shallow water system on the sphere.
In the case of an annulus, the circulation $`C_{}=u_\theta 𝑑l`$ around the inner wall is an additional conserved quantity (the circulation $`C_+`$ around the outer wall is also conserved, but it is related to other conserved quantities, as $`C_+=\mathrm{\Gamma }C_{}`$, and the conservation of $`\mathrm{\Gamma }`$ is already included in the PV conservation). Furthermore we need in general to set $`\psi =\psi _{}0`$ at the inner wall (while we can still set $`\psi =0`$ at the outer wall). As a consequence a boundary term -$`\psi _{}\delta C_{}`$ now appears in the expression (30) for the energy variation. However we can directly set $`\delta C_{}=0`$, canceling this boundary term, without influence on the independent variables $`h\rho `$ (determining the locally averaged vorticity $`\overline{\omega }=\sigma h\rho 𝑑\sigma `$), $`.𝐮`$ and $`h`$. Therefore the only modification with respect to the disk is an additional unknown $`\psi _{}`$ in the definition (10) of $`\psi `$, determined by the corresponding additional constraint $`C_{}`$ .
The case of a straight channel can be viewed as the limit of an annulus with small width, but it can be convenient to treat it in itself. Let us consider a straight channel between two walls at coordinates $`y=\pm L_y/2`$ with periodic boundary conditions along the x-direction (with domain length $`L_x`$ ). In the absence of Coriolis force ($`\mathrm{\Omega }=0)`$ , the x-wise momentum
$$P=hu_xd^2𝐫$$
(88)
is conserved (instead of the angular momentum), as well as the circulation $`C_+=u_x𝑑x`$ along the upper wall ($`y=L_y/2`$ ). The boundary condition (10) defining $`\psi `$ is replaced by $`\psi =P/L_x`$ at the upper wall $`y=L_y/2`$ and $`\psi =0`$ (for instance) at the lower wall $`y=L_y/2`$. Unlike with angular momentum, we cannot express the variation $`\delta 𝐏`$ in terms of the variations in the independent variables $`\overline{\omega }`$ , $`.𝐮`$ , $`h`$. However we have now an additional freedom in the variational problem, as we can add a uniform x-wise velocity $`U𝐞_x`$ (use a reference frame with velocity $`U𝐞_x`$) without influence on the independent variables $`\overline{\omega }`$ , $`.𝐮`$ , $`h`$. For any choice of $`U`$, we can solve the variational problem with the velocity $`𝐮^{}=𝐮U𝐞_x`$ and corresponding energy $`E^{}=E+MU^2/2PU`$, upper wall circulation $`C_+^{_{}}=C_+UL_x`$. This yields a Gibbs state (35)(38)(41), representing a steady solution of the shallow water equation in a reference frame moving with velocity $`U𝐞_𝐱`$. Among these states, the ones with the right value of the momentum $`P=hu_xd^2𝐫`$ will be the actual solutions. Families of Gibbs states with the same structure translated in the x-direction are obtained, as discussed by Sommeria et al. (1991) in the incompresible case.
Finally in the case of an infinite domain, the two components of momentum, as well as the angular momentum are conserved. This yields to purely translating or purely rotating Gibbs states, as discussed by Chavanis and Sommeria (1997) in the incompressible case.
### 6.2 Relaxation equations
Taking the derivative of (88)(86) with respect to time and using equations (51)(52), we obtain the constraints
$$\dot{P}=hJ_{\omega y}d^2\text{r}=0$$
(89)
$$\dot{L}=h\text{J}_\omega .\text{r}d^2\text{r}=0$$
(90)
These constraints can be included in the variational principle (60) by introducing appropriate Lagrange multipliers denoted $`\beta (t)U(t)`$ and $`\beta (t)\lambda (t)`$. Then, the results of section 5 are generalized simply be replacing the velocity $`\overline{\text{u}}`$ by the relative velocity $`\overline{\text{u}}^{}=\overline{\text{u}}U(t)\text{e}_x\lambda (t)\text{e}_z\text{r}`$ where the time evolution of $`U(t)`$ and $`\lambda (t)`$ is obtained by substituting the optimal current (66), with the above transformation, in the constraints (89) and (90). In the case of a channel we have the additional conserved quantity $`C_+=u_x𝑑x`$ along the upper wall. Using (52), we readily find that $`\dot{C}_+=J_{\omega y}𝑑x=0`$ as the current $`𝐉_\omega `$is parallel to the wall, so the circulation along each wall is indeed conserved by the relaxation equations.
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# A Flat Universe from High-Resolution Maps of the Cosmic Microwave Background Radiation
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# A mean identity for longest increasing subsequence problems
## 1 Introduction
In , Prähofer and Spohn consider a certain polynuclear growth (PNG) model with stationary initial conditions, and show that it maps to the following increasing subsequence problem:
Let $`t`$ be a positive real number. Pick a random set of points in the unit square $`[0,1]\times [0,1]`$ as follows. On the left and bottom edges, take a Poisson process of mean $`t`$; inside the square, take a Poisson process of mean $`t^2`$. (Thus our total mean is $`t^2+2t`$.) A sequence of these points is “increasing” if we have $`xx^{}`$, $`yy^{}`$ whenever $`(x,y)`$ and $`(x^{},y^{})`$ are consecutive points in the sequence; the length of the sequence is defined to be the number of points. The problem is then to determine the asymptotic distribution of the length of the longest increasing subsequence. (Note that without the extra points on the left and bottom edge, this is just the standard Poisson model for increasing subsequences of random permutations .)
Prähofer and Spohn then observe that the stationarity of the initial conditions can be used to show that the length of the longest increasing subsequence has mean exactly $`2t`$. This fact is striking for two reasons. The first is that the mean in the standard model is rather complicated; it is thus surprising that a fairly simple extension gives rise to an exact formula for the extended mean. The second is that since adding points can only help the longest increasing subsequence, we conclude that $`2t`$ is an upper bound on the mean in the standard model. This bound is quite tight; indeed, in the standard model, the mean takes the form $`2t+O(t^{1/3})`$. (We could also derive this upper bound from the (strictly stronger) result of that the expected length of the longest increasing subsequence of a permutation of length $`n`$ is at most $`2\sqrt{n}`$; the present method is more generally applicable, however.)
The first object of the present paper is to generalize this fact in a number of ways. It turns out, for instance, that the standard Poisson model admits a one-parameter family of extensions with explicit means; this further extends to give explicit information about the moment generating function in a neighborhood of this family. Also, we can replace the Poisson model by the generalized model considered in (based in turn on a model of Johansson ).
Our other object is to explore the asymptotic relations between the extended models and the unextended models. It turns out that by a careful (if heuristic) analysis, we can use the moment generating function identities to determine not just the asymptotic mean of the longest increasing subsequence length, but also the asymptotic scale factor. This gives a uniform prescription for the scaling information, agreeing with the results of all of the cases that have so far been analyzed.
Section 1 defines the models of interest, as well as a certain continuous limiting case. Section 2 gives a short, algebraic proof of the moment generating function and mean identities; this is followed by a somewhat more complicated, but also more enlightening combinatorial proof in Section 3. Finally, Section 4 considers the asymptotic consequences, giving explicit conjectures for the asymptotics of the general case.
Acknowledgements The author would like to thank J. Baik, Y. Baryshnikov, and W. Whitt for helpful conversations.
## 2 Extended growth models
The model we will be considering is a generalization of a model considered by Johansson , combining the generalizations of , , and . We define a “parameter set” $`p`$ to be a triple $`(t,q,r)`$, where $`t`$ is a nonnegative number, and $`q`$ and $`r`$ are sequences of nonnegative numbers with
$$\underset{i}{}q_i+\underset{i}{}r_i<\mathrm{};$$
(2.1)
we will write such a parameter set as $`t:q/r`$, omitting $`t`$ if $`t=0`$ and omitting any trailing 0’s from $`q`$ and $`r`$. From the existence of the sum, we conclude that
$$Q(p):=\underset{i}{sup}(q_i)\text{and}R(p):=\underset{i}{sup}(r_i)$$
(2.2)
are well-defined and attained. Given a parameter set $`p`$, we define two functions
$`H(z;p)`$ $`=e^{tz}{\displaystyle \underset{i}{}}(1zq_i)^1(1+zr_i)`$ (2.3)
$`E(z;p)`$ $`=e^{tz}{\displaystyle \underset{i}{}}(1+zq_i)(1zr_i)^1;`$ (2.4)
$`H(z;p)`$ converges for $`|z|<Q(p)^1`$, while $`E(z;p)`$ converges for $`|z|<R(p)^1`$. We will say that two parameter sets $`p_+=t^+:q^+/r^+`$ and $`p_{}=t^{}:q^{}/r^{}`$ are compatible if $`Q(p_+)Q(p_{})<1`$, $`R(p_+)R(p_{})<1`$. Note that we then have
$$H(p_+;p_{}):=e^{t^+t^{}}\underset{i}{}e^{t^{}(q_i^++r_i^+)}e^{t^+(q_i^{}+r_i^{})}\underset{i,j}{}(1q_i^+q_j^{})^1(1r_i^+r_j^{})^1(1+q_i^+r_j^{})(1+r_i^+q_j^{})<\mathrm{}.$$
(2.5)
Let $`(^+)^{}`$ be a disjoint copy of $`^+`$, and consider the set $`\mathrm{\Omega }:=[0,1]^+(^+)^{}`$. Then we associate to a pair $`p_+`$, $`p_{}`$ of compatible parameter sets a random multiset $`M(p_+,p_{})\mathrm{\Omega }\times \mathrm{\Omega }`$ as follows. We let $`P(t)`$ denote a Poisson random variable of parameter $`t`$, $`g(q)`$ denote a geometric random variable with parameter $`q`$, and $`b(q)`$ denote a random variable which is 0 with probability $`\frac{q}{1+q}`$ and 1 with probability $`\frac{1}{1+q}`$.
* On $`[0,1]\times [0,1]`$, we choose $`P(t^+t^{})`$ i.i.d. uniform points,
* On $`[0,1]\times i`$, we choose $`P(t^+q_i^{})`$ i.i.d. uniform points,
* On $`[0,1]\times i^{}`$, we choose $`P(t^+r_i^{})`$ i.i.d. uniform points,
* On $`i\times j`$, we have multiplicity $`g(q_i^+q_j^{})`$,
* On $`i\times j^{}`$, we have multiplicity $`b(q_i^+r_j^{})`$,
* On $`i^{}\times j^{}`$, we have multiplicity $`g(r_i^+r_j^{})`$,
and so on, with all of the multiplicities chosen independently.
Choose a pair of total orderings (denoted $`<_+`$ and $`<_{}`$) on $`\mathrm{\Omega }`$ compatible with the usual ordering on $`[0,1]`$. A subsequence of a multiset $`M`$ in $`\mathrm{\Omega }\times \mathrm{\Omega }`$ (a sequence of points $`(x_i,y_i)`$ from $`M`$ with no point occuring more often than its multiplicity) is “increasing” if we always have $`x_i_+x_{i+1}`$, $`y_i_{}y_{i+1}`$, subject to the further conditions
$`x_i=x_{i+1}`$ $`x_i^+`$ (2.6)
$`y_i=y_{i+1}`$ $`y_i^+`$ (2.7)
In other words, the sequence must be strictly increasing along rows and columns from $`^+`$. We then define a sequence $`\lambda _i(M)`$ by setting
$$\underset{1il}{}\lambda _i(M)$$
(2.8)
equal to the size of the longest subsequence of $`M`$ which is a union of $`l`$ increasing subsequences. We then define $`\lambda (p_+,p_{}):=\lambda (M(p_+,p_{}))`$ (with the latter notation preferred when particular multisets are being compared). We recall the following results from :
###### Theorem 2.1.
For any pair $`p_+`$, $`p_{}`$ of compatible parameter sets, $`\lambda (p_+,p_{})`$ is a random partition, finite with probability 1. The distribution of $`\lambda (p_+,p_{})`$ is independent of the choice of total orderings on $`\mathrm{\Omega }`$.
###### Theorem 2.2.
For any pair $`p_+`$, $`p_{}`$ of compatible parameter sets, we have the identity
$$\mathrm{Pr}(\lambda _1(p_+,p_{})l)=H(p_+;p_{})^1𝐄_{UU(l)}det(H(U;p_+)H(U^{};p_{})).$$
(2.9)
###### Remark.
As in , this is a formal integral, defined by analytic continuation from the region $`Q(p_\pm )<1`$.
For our purposes, we need to extend the model slightly, by adding a special row and column to the random multiset. Extend $`\mathrm{\Omega }`$ to $`\mathrm{\Omega }^+`$ by adding a new element, denoted $`\mathrm{\Sigma }`$, and extend the total orderings so that $`\mathrm{\Sigma }`$ is smallest in both orderings. Then we define a new random multiset $`M(p_+,p_{};\alpha _+,\alpha _{})`$ as follows.
* On $`\mathrm{\Omega }\times \mathrm{\Omega }`$, we take $`M(p_+,p_{})`$,
* On $`[0,1]\times \{\mathrm{\Sigma }\}`$, we choose $`P(\alpha _{}t^+)`$ i.i.d. uniform points,
* On $`(i,\mathrm{\Sigma })`$, we have multiplicity $`b(\alpha _{}q_i^+)`$,
* On $`(i^{},\mathrm{\Sigma })`$, we have multiplicity $`q(\alpha _{}r_i^+)`$,
and so on, and allow increasing subsequences to be weakly increasing in the new row and column. Note that the point $`(\mathrm{\Sigma },\mathrm{\Sigma })`$ has multiplicity fixed at 0; otherwise, the new model would simply be a special case of the old model. By the argument in , we have:
###### Theorem 2.3.
For any compatible pair $`p_+`$, $`p_{}`$ of parameter sets and any $`\alpha _+`$, $`\alpha _{}`$ with $`\alpha _+R(p_{})<1`$ and $`\alpha _{}R(p_+)<1`$, we have
$$\begin{array}{c}\mathrm{Pr}(\lambda _1(p_+,p_{};\alpha _+,\alpha _{})l)=\hfill \\ \hfill E(\alpha _+;p_{})^1E(\alpha _{};p_+)^1H(p_+;p_{})^1[D_l(p_+,p_{};\alpha _+,\alpha _{})\alpha _+\alpha _{}D_{l1}(p_+,p_{};\alpha _+,\alpha _{})],\end{array}$$
(2.10)
where
$$D_l(p_+,p_{};\alpha _+,\alpha _{})=𝐄_{UU(l)}det((1+\alpha _+U)(1+\alpha _{}U^{})H(U;p_+)H(U^{};p_{})).$$
(2.11)
Aside from intrinsic interest (), this new model in principle can give us some information about the original model $`\lambda _1(p_+,p_{})`$, since we have the coupling
$$\lambda _1(M(p_+,p_{}))\lambda _1(M(p_+,p_{};\alpha _+,\alpha _{})).$$
(2.12)
Thus it is particularly interesting that, as we shall see, there is an exact formula for
$$𝐄(\lambda _1(p_+,p_{};\alpha ,\alpha ^1))$$
(2.13)
whenever $`R(p_+)<\alpha <R(p_{})^1`$.
We will also consider a continuous limiting case of the above model, combining the exponential limit of and the heavy-traffic limit of queueing theory (studied in the present context in and ).
The parameters for the continuous model consist of a pair of sequences $`\rho _i^\pm \{\mathrm{}\}`$, a nonnegative real $`u`$, and two numbers $`a_\pm \{\mathrm{}\}`$, subject to the convergence constraint
$$\underset{i}{}\frac{1}{\rho _i^\pm z}<\mathrm{},\text{for}z<inf(\rho ^\pm ),$$
(2.14)
and the compatibility constraints
$$inf(\rho ^+)+inf(\rho ^{})>0,inf(\rho ^+)+a_{}>0,a_++inf(\rho ^{})>0.$$
(2.15)
We omit any trailing $`\mathrm{}`$’s from $`\rho ^\pm `$. Define an infinite matrix $`\mathrm{\Lambda }_{ij}`$, $`0i,j`$ by
$$\mathrm{\Lambda }_{ij}=\{\begin{array}{cc}0\hfill & (i,j)=(0,0)\hfill \\ \frac{1}{\rho _i^++\rho _j^{}},\hfill & (i,j)(0,0),\hfill \end{array}$$
(2.16)
with the convention that $`\rho _0^\pm =\alpha _\pm `$, and let $`M`$ be a random matrix filled with independent exponential random variables, such that $`M_{ij}`$ has mean $`\mathrm{\Lambda }_{ij}`$. Also, for $`0i`$ such that $`\rho _i^+<\mathrm{}`$, let $`B_i`$ be a Brownian motion on $`[0,1]`$ with
$$𝐄(B_i(x))=u\rho _i^+x\mathrm{Var}(B_i(x))=ux;$$
(2.17)
by convention, $`\rho _0^\pm =a_\pm `$. Then we define a random sequence $`\chi (\rho ^+,u:\rho ^{};a_+,a_{})`$ in terms of increasing paths in $`\times ([0,1])`$, where the contribution of an interval $`i\times [x_i,y_i]`$ is $`B_i(y_i)B_i(x_i)`$ and the contribution of a point $`i\times j`$ is $`M_{ij}`$. When $`u>0`$, we require that the union of $`k`$ increasing paths used to define $`\chi _k`$ have total Lebesgue measure $`k`$ in $`\times [0,1]`$, and that the paths contain no point in $`i\times [0,1]`$ when $`\rho _i^+=\mathrm{}`$. If no such path exists, we set $`\chi _k=\mathrm{}`$.
###### Theorem 2.4.
For any valid choice of $`\rho ^+`$, $`u:\rho ^{}`$, $`a_+`$, $`a_{}`$, of compatible parameter sets, $`\chi (\rho ^+,u:\rho ^{};a_+,a_{})`$ is a random nonincreasing sequence in $`\{\mathrm{}\}`$, and is nonnegative if $`u=0`$. The distribution is invariant under reordering of the sequences $`\rho ^\pm `$.
###### Proof.
That the sequence $`\chi `$ is nonincreasing, invariant under reordering, and nonnegative when $`u=0`$ follows from the fact that it is a scaled limit of $`\lambda `$ (see below); it suffices therefore to show that $`\chi _1<\mathrm{}`$ with probability 1. We first observe that
$$\chi _1(\rho ^+,u:\rho ^{};a_+,a_{})\chi _1(\rho ^+,u:\rho ^{};a_+,\mathrm{})+\chi _1(\rho ^+,u:\rho ^{};\mathrm{},a_{}).$$
(2.18)
Since decreasing $`a_\pm `$ cannot decrease $`\chi _1`$, we conclude that it suffices to prove finiteness when $`inf(\rho ^{})<a_\pm <inf(\rho ^\pm )`$ and $`a_++a_{}>0`$. But then Corollary 3.4 below applies, expressing $`𝐄(e^{(a_++a_{})\chi _1})`$ as an infinite product. Moreover, the conditions on $`\rho ^\pm `$ suffice to force convergence of this product, and thus $`𝐄(e^{(a_++a_{})\chi _1})<\mathrm{}`$. But this immediately implies $`\chi _1<\mathrm{}`$, as desired. ∎
As we alluded to above, this is a limiting case of the discrete model:
$$\chi (\rho ^+,u:\rho ^{};a_+,a_{})=\underset{t\mathrm{}}{lim}\frac{\lambda (/e^{\rho ^+/t},ut^2:/e^{\rho ^{}/t};e^{a_+/t},e^{a_{}/t})ut^2}{t},$$
(2.19)
with the limit taken in the sense of distribution. This corresponds to the facts that if $`x`$ is an exponential random variable of mean $`1/m`$, then $`\frac{x}{l}`$ is a geometric random variable of parameter $`e^{m/l}`$, and that as the parameter tends to infinity, Poisson processes converge (with proper scaling) to Brownian motion.
The main significance of the continuous model is that it contains two classical matrix ensembles as special cases. Let $`0_n`$ denote the finite sequence consisting of $`n`$ copies of $`0`$. We obtain the Gaussian Unitary Ensemble from $`\chi (0_n,1:;\mathrm{},\mathrm{})`$ that is, $`\chi (0_n,1:;\mathrm{},\mathrm{})`$ is distributed as the (ordered) eigenvalues of an $`n\times n`$ Hermitian Gaussian matrix (extended to an infinite sequence by adding $`\mathrm{}`$). (This is essentially proved in ; to be precise, they show that one obtains traceless GUE under a constraint equivalent to $`_iB_i(1)=0`$; it follows easily that without this constraint, one obtains ordinary GUE.) Similarly, we obtain the Laguerre Unitary Ensemble from $`\chi (0_{n_+},0_n_{};\mathrm{},\mathrm{})`$; that is, the distribution of the singular values of a $`n_+\times n_{}`$ complex Gaussian matrix.
## 3 An algebraic proof
Let $`p_+`$, $`p_{}`$ be a pair of compatible parameter sets.
###### Lemma 3.1.
$`D_l(\alpha _+,\alpha _{}):=D_l(p_+,p_{};\alpha _+,\alpha _{})`$ is a polynomial, satisfying the identity
$$D_l(\alpha _+,\alpha _{})=(\alpha _+\alpha _{})^lD_l(\alpha _{}^1,\alpha _+^1).$$
(3.1)
###### Proof.
This follows immediately from the corresponding fact for $`det((1+\alpha _+U)(1+\alpha _{}U^{}))`$. ∎
For $`\alpha _+<R(p_{})^1`$, $`\alpha _{}<R(p_+)^1`$, define $`L(\alpha _+,\alpha _{})=\lambda _1(p_+,p_{};\alpha _+,\alpha _{})`$; then
###### Theorem 3.2.
For $`R(p_+)<\alpha _+<R(p_{})^1`$ and $`R(p_{})<\alpha _{}<R(p_+)^1`$,
$$𝐄((\alpha _+\alpha _{})^{L(\alpha _+,\alpha _{})})=E(\alpha _+;p_{})^1E(\alpha _{};p_+)^1E(\alpha _+^1;p_+)E(\alpha _{}^1;p_{})$$
(3.2)
###### Proof.
By Theorem 2.3 above,
$`\mathrm{Pr}(L(\alpha _+,\alpha _{})l)`$ $`=E(\alpha _+;p_{})^1E(\alpha _{};p_+)^1H(p_+;p_{})^1[D_l(\alpha _+,\alpha _{})\alpha _+\alpha _{}D_{l1}(\alpha _+,\alpha _{})]`$ (3.3)
$`=(\alpha _+\alpha _{})^lE(\alpha _+;p_{})^1E(\alpha _{};p_+)^1H(p_+;p_{})^1[D_l(\alpha _{}^1,\alpha _+^1)D_{l1}(\alpha _{}^1,\alpha _+^1)].`$ (3.4)
Then
$`{\displaystyle \underset{0lk}{}}(\alpha _+\alpha _{})^l`$ $`\mathrm{Pr}(L(\alpha _+,\alpha _{})=l)`$
$`={\displaystyle \underset{0lk}{}}(\alpha _+\alpha _{})^l[\mathrm{Pr}(L(\alpha _+,\alpha _{})l)\mathrm{Pr}(L(\alpha _+,\alpha _{})l1)]`$ (3.5)
$`=(\alpha _+\alpha _{})^{k1}\mathrm{Pr}(L(\alpha _+,\alpha _{})k)+{\displaystyle \frac{\alpha _+\alpha _{}1}{\alpha _+\alpha _{}}}{\displaystyle \underset{0lk}{}}(\alpha _+\alpha _{})^l\mathrm{Pr}(L(\alpha _+,\alpha _{})l)`$ (3.6)
$`=E(\alpha _+;p_{})^1E(\alpha _{};p_+)^1H(p_+;p_{})^1`$
$`\left[{\displaystyle \frac{1}{\alpha _+\alpha _{}}}\left(D_k(\alpha _{}^1,\alpha _+^1)D_{k1}(\alpha _{}^1,\alpha _+^1)\right)+{\displaystyle \frac{\alpha _+\alpha _{}1}{\alpha _+\alpha _{}}}D_k(\alpha _{}^1,\alpha _+^1)\right]`$ (3.7)
$`=E(\alpha _+;p_{})^1E(\alpha _{};p_+)^1H(p_+;p_{})^1[D_k(\alpha _{}^1,\alpha _+^1)(\alpha _+\alpha _{})^1D_{k1}(\alpha _{}^1,\alpha _+^1)]`$ (3.8)
$`=E(\alpha _+;p_{})^1E(\alpha _{};p_+)^1E(\alpha _+^1;p_+)E(\alpha _{}^1;p_{})\mathrm{Pr}(L(\alpha _{}^1,\alpha _+^1)k),`$ (3.9)
where the last step is valid since $`\alpha _{}^1<R(p_{})^1`$ and $`\alpha _+^1<R(p_+)^1`$.
The theorem then follows by taking the limit $`k\mathrm{}`$. ∎
Taking a limit as $`\alpha _+\alpha _{}1`$, we obtain:
###### Corollary 3.3.
Whenever $`R(p_+)<\alpha <R(p_{})^1`$,
$$𝐄(L(\alpha ,\alpha ^1))=\frac{\alpha E^{}(\alpha ;p_{})}{E(\alpha ;p_{})}+\frac{\alpha ^1E^{}(\alpha ^1;p_+)}{E(\alpha ^1;p_+)}.$$
(3.10)
Remark. If $`p=t:q/r`$, then
$$\frac{\alpha E^{}(\alpha ;p)}{E(\alpha ;p)}=\alpha t+\underset{i}{}\frac{\alpha q_i}{1+\alpha q_i}+\underset{i}{}\frac{\alpha r_i}{1\alpha r_i}$$
(3.11)
We also observe that
$$\frac{\alpha E^{}(\alpha ;p_{})}{E(\alpha ;p_{})}+\frac{\alpha ^1E^{}(\alpha ^1;p_+)}{E(\alpha ^1;p_+)}=\alpha \frac{d}{d\alpha }\mathrm{log}\frac{E(\alpha ;p_{})}{E(\alpha ^1;p_+)}$$
(3.12)
In the continuous limit, we write $`X(a_+,a_{}):=\chi _1(\rho ^+,u:\rho ^{};a_+,a_{}).`$ Taking the appropriate limit gives:
###### Corollary 3.4.
Let $`\rho ^+`$ and $`u:\rho ^{}`$ be compatible continuous parameter sets. Then whenever $`inf(\rho ^{})<a_+<inf(\rho ^+)`$ and $`inf(\rho ^+)<a_{}<inf(\rho ^{})`$,
$$𝐄(e^{(a_++a_{})X(a_+,a_{})})=e^{u(a_{}^2a_+^2)/2}\underset{i}{}\frac{\rho _i^++a_{}}{\rho _i^+a_+}\underset{i}{}\frac{\rho _i^{}+a_+}{\rho _i^{}a_{}}.$$
(3.13)
Whenever $`inf(\rho ^{})<a<inf(\rho ^+)`$,
$$𝐄(X(a,a))=ua+\underset{i}{}\frac{1}{\rho _i^{}+a}+\underset{i}{}\frac{1}{\rho _i^+a}.$$
(3.14)
## 4 A combinatorial proof
Fix parameters as in the previous section, and set
$`N_+`$ $`=|M(p_+,p_{};\alpha _+,\alpha _{})(\{\mathrm{\Sigma }\}\times \mathrm{\Omega })|`$ (4.1)
$`N_{}`$ $`=|M(p_+,p_{};\alpha _+,\alpha _{})(\mathrm{\Omega }\times \{\mathrm{\Sigma }\})|`$ (4.2)
Then we observe
$$𝐄((\alpha _+\alpha _{})^{N_+})=\frac{E(\alpha _{}^1;p_{})}{E(\alpha _+;p_{})}𝐄((\alpha _+\alpha _{})^N_{})=\frac{E(\alpha _+^1;p_+)}{E(\alpha _{};p_+)},$$
(4.3)
and for $`\alpha _+=\alpha _{}^1=\alpha `$,
$$𝐄(N_+)=\frac{E^{}(\alpha ;p_{})}{E(\alpha ;p_{})}𝐄(N_{})=\frac{E^{}(\alpha ^1;p_+)}{E(\alpha ^1;p_+)}.$$
(4.4)
So we can restate Theorem 3.2 and Corollary 3.3 as
$`𝐄((\alpha _+\alpha _{})^{L(\alpha _+,\alpha _{})})`$ $`=𝐄((\alpha _+\alpha _{})^{N_+N_{}})`$ (4.5)
$`𝐄(L(\alpha _+,\alpha _{}))`$ $`=𝐄(N_++N_{})`$ (4.6)
We give a direct proof of this fact, for $`\alpha _+\alpha _{}1`$:
###### Proof.
Let $`\alpha ^{}<\alpha _{}`$. Then $`\alpha _+\alpha ^{}<1`$, so we can extend $`M(p_+,p_{};\alpha _+,\alpha ^{})`$ by adjoining $`(\mathrm{\Sigma },\mathrm{\Sigma })`$ with multiplicity $`N_0`$ of distribution $`g(\alpha _+\alpha ^{})`$; denote the resulting random multiset by $`M^{}`$. But then by Theorem 2.1, we can change the total ordering $`<_{}`$ so that $`\mathrm{\Sigma }`$ becomes maximal instead of minimal. We then find that
$$\mathrm{\Sigma }\times \mathrm{\Omega },(\mathrm{\Sigma },\mathrm{\Sigma }),\mathrm{\Omega }\times \mathrm{\Sigma }$$
(4.7)
induces an increasing subsequence of $`M^{}`$ with respect to this new ordering; thus
$$N_0+N_++N_{}\lambda _1(M^{}).$$
(4.8)
On the other hand, this is the only maximal increasing subsequence that passes through $`(\mathrm{\Sigma },\mathrm{\Sigma })`$; any other maximal increasing subsequence can have size at most $`N_++N_{}+\lambda _1(M(p_+,p_{}))`$. We thus find
$$\lambda _1(M^{})N_++N_{}+\mathrm{max}(\lambda _1(M(p_+,p_{})),N_0).$$
(4.9)
But $`\lambda _1(M^{})`$ is distributed as $`N_0+\lambda _1(M(p_+,p_{};\alpha _+,\alpha ^{}))`$ (since before the reordering every maximal increasing subsequence passes through $`(\mathrm{\Sigma },\mathrm{\Sigma })`$, and $`N_0`$ is independent of $`\lambda _1(M(p_+,p_{};\alpha _+,\alpha ^{}))`$. So if we take the expectations and subtract/divide by the contribution of $`N_0`$, we find that we need only show
$`\underset{\alpha ^{}(1/\alpha )^{}}{lim}𝐄(\mathrm{max}(\lambda _1(p_+,p_{}),N_0))𝐄(N_0)`$ $`=0`$ (4.10)
$`\underset{\alpha ^{}(\alpha _{})^{}}{lim}𝐄((\alpha _+\alpha _{})^{N_0})^1𝐄((\alpha _+\alpha _{})^{\mathrm{max}(\lambda _1(p_+,p_{}),N_0)})`$ $`=1.`$ (4.11)
###### Lemma 4.1.
Let $`X`$ be a nonnegative-integer-valued random variable with finite first moment, and let $`Y`$ be an independent geometric random variable of parameter $`t`$. Then
$$\underset{t1^{}}{lim}𝐄(\mathrm{max}(X,Y))𝐄(Y)=0.$$
(4.12)
Similarly, for $`s<1`$, if $`𝐄(s^X)`$ is finite, then
$$\underset{ts^{}}{lim}𝐄(s^Y)^1𝐄(s^{\mathrm{max}(X,Y)})=1.$$
(4.13)
###### Proof.
$`𝐄(\mathrm{max}(X,Y)Y)`$ $`=𝐄((XY)\mathrm{Pr}(Y<X))`$ (4.14)
$`=𝐄(X)+{\displaystyle \frac{t}{1t}}(𝐄(t^X)1)`$ (4.15)
$`0.`$ (4.16)
Similarly,
$`𝐄(s^{\mathrm{max}(X,Y)})`$ $`=𝐄(s^X\mathrm{Pr}(Y<X)+\mathrm{Pr}(YX)𝐄(s^Y|YX))`$ (4.17)
$`=𝐄(s^X)+{\displaystyle \frac{t(s1)}{ts}}𝐄((t/s)^X).`$ (4.18)
and thus
$`\underset{ts^{}}{lim}𝐄(s^Y)^1𝐄(s^{\mathrm{max}(X,Y)})`$ $`=\underset{ts^{}}{lim}{\displaystyle \frac{ts}{s(t1)}}𝐄(s^X)+{\displaystyle \frac{t(s1)}{s(t1)}}𝐄((t/s)^X)`$ (4.19)
$`=1`$ (4.20)
The theorem then follows from the following lemma, since
$$\alpha _+\alpha _{}>R(p_+)R(p_{})$$
(4.21)
by assumption. ∎
###### Lemma 4.2.
For all $`z>R(p_+)R(p_{})`$, $`𝐄(z^{\lambda _1(p_+,p_{})})`$ is finite. In particular, since $`1>R(p_+)R(p_{})`$, $`\lambda _1(p_+,p_{})`$ has moments of all orders.
###### Proof.
An increasing subsequence in $`M(p_+,p_{})`$ can pass through a point on a strict row or column at most once; thus $`\lambda _1(M(p_+,p_{}))`$ is unchanged if we remove any excess multiplicity in those rows and columns. Let $`M^o`$ be the resulting multiset, then
$$\lambda _1(M(p_+,p_{}))=\lambda _1(M^o)|M^o|.$$
(4.22)
It will thus suffice to prove that $`𝐄(z^{|M^o|})<\mathrm{}`$. But the moment generating function of $`|M^o|`$ is $`f(z)/f(1)`$, where
$$f(z)=e^{t^+t^{}z}\underset{i}{}e^{t^+q_i^{}z}e^{t^+r_i^{}z}e^{q_i^+t^{}z}e^{r_i^+t^{}z}\underset{i,j}{}(1r_i^+r_j^{}z)^1(1+q_i^+r_j^{}z)(1+r_i^+q_j^{}z)((1q_i^+q_j^{})+q_i^+q_j^{}z)$$
(4.23)
This product converges to an analytic function with no pole inside the open disc $`|z|<(R(p_+)R(p_{}))^1`$, and thus the result follows. ∎
## 5 Asymptotic consequences
Since $`\lambda _1(p_+,p_{};\alpha _+,\alpha _{})`$ is nondecreasing in $`\alpha _+`$ and $`\alpha _{}`$, we obtain the following bound:
###### Theorem 5.1.
For any compatible parameters $`p_+`$, $`p_{}`$,
$$𝐄(\lambda _1(p_+,p_{}))\underset{R(p_+)<\alpha <R(p_{})^1}{inf}𝐄(\lambda _1(p_+,p_{};\alpha ,\alpha ^1)).$$
(5.1)
For instance, in the purely Poisson case, $`p_+=p_{}=t:/`$, we find
$$𝐄(\lambda _1(p_+,p_{}))\underset{\alpha >0}{inf}(\alpha +\alpha ^1)t=2t.$$
(5.2)
This bound is remarkably tight; indeed, we have
$$\lambda _1(p_+,p_{})=2tO(t^{1/3+ϵ}).$$
(5.3)
This suggests the following conjecture:
###### Conjecture 5.2.
Fix parameters $`p_+`$, $`p_{}`$, and define
$$m(\alpha ;p_+,p_{})=𝐄(\lambda _1(p_+,p_{};\alpha ,\alpha ^1)).$$
(5.4)
Then
$$\underset{n\mathrm{}}{lim}n^1\lambda _1(p_+^n,p_{}^n)=\underset{R(p_+)<\alpha <R(p_{})^1}{inf}m(\alpha ;p_+,p_{}),$$
(5.5)
with probability 1, where
$`(t:q/r)^n`$ $`:=(nt):q^n/r^n`$ (5.6)
$`q^n`$ $`:=q_1,q_1,\mathrm{}q_1,q_2,q_2,\mathrm{}q_2,\mathrm{}.`$ (5.7)
###### Remark 1.
Roughly speaking, this is an analogue of the law of large numbers. As such, it can most likely be strengthened considerably (considering different sequences of parameter sets than just $`p_+^n`$, $`p_{}^n`$). See, for instance, the result of .
###### Remark 2.
The existence of the limit (5.5) follows from superadditivity and the bound (5.1).
This has been verified in a number of special cases (see below). In each case, we in fact find that
$$\frac{\lambda _1(p_+^n,p_{}^n)\mu n}{n^{1/3}}$$
(5.8)
converges to a limit distribution.
Fix parameters $`p_+`$, $`p_{}`$. An increasing subsequence of $`M(p_+,p_{};\alpha ,\alpha ^1)`$ cannot include points from both $`\{\mathrm{\Sigma }\}\times \mathrm{\Omega }`$ and $`\mathrm{\Omega }\times \{\mathrm{\Sigma }\}`$. We would thus expect that for $`\alpha `$ “large”, the typical longest increasing subsequence will avoid $`\mathrm{\Omega }\times \{\mathrm{\Sigma }\}`$ entirely. In particular, we would expect
$$𝐄(\lambda _1(p_+,p_{};\alpha ,\alpha ^1))𝐄(\lambda _1(p_+,p_{};\alpha ,0))$$
(5.9)
whenever $`N_+N_{}`$. For asymptotic purposes, this condition is simply that $`\alpha >\stackrel{~}{\alpha }_+`$, where $`\stackrel{~}{\alpha }_+`$ minimizes $`m(\alpha ;p_+,p_{})`$. (We also define $`\stackrel{~}{\alpha }_{}=(\stackrel{~}{\alpha }_+)^1`$, which of course minimizes $`m(\alpha ;p_{},p_+)`$).
In particular, this tells us that $`\stackrel{~}{\alpha }_\pm `$ is a critical point; if $`\alpha _+<\stackrel{~}{\alpha }_+`$ and $`\alpha _{}<\stackrel{~}{\alpha }_{}`$, we have
$$𝐄(\lambda _1(p_+,p_{};\alpha _+,\alpha _{}))𝐄(\lambda _1(p_+,p_{};\stackrel{~}{\alpha }_+,\stackrel{~}{\alpha }_{})),$$
(5.10)
while if either is greater, the mean is determined by the dominant parameter.
This behaviour is, of course, confirmed by the analysis of , in which the asymptotics for general $`\alpha _\pm `$ are determined for the Poisson case $`p=p^{}=t:/`$ and the Johansson case $`p=p^{}=/\sqrt{q}^n`$ (where $`\sqrt{q}^n`$ is the finite sequence consisting of $`n`$ copies of $`\sqrt{q}`$). In both cases, we obtain the same behaviour near the critical point. This suggests that for general parameters there should exist constants $`\mu `$, $`\sigma `$, and $`\sigma _\pm `$ so that the following holds:
If we fix $`w_+`$, $`w_{}`$, and define
$$\alpha _\pm =\stackrel{~}{\alpha }_\pm \mathrm{exp}(\frac{2w_\pm }{\sigma _\pm n^{1/3}}),$$
(5.11)
then as $`n\mathrm{}`$,
$$\frac{\lambda _1(p_+^n,p_{}^n;\alpha _+,\alpha _{})\mu n}{\sigma n^{1/3}}$$
(5.12)
converges to the distribution $`H(w_+,w_{})`$ ().
We recall the following information about the distribution $`H(w_+,w_{})`$:
###### Lemma 5.3.
Let $`X`$ be distributed as $`H(w_+,w_{})`$. Then $`𝐄(\mathrm{exp}(2(w_++w_{})X))=\mathrm{exp}(\frac{8}{3}(w_+^3+w_{}^3))`$. If $`w_+=w_{}=w`$, then $`𝐄(X)=4w^2`$.
Only the latter equation was actually shown in , but essentially the same calculation gives the first equation as well. Furthermore, in the cases that have been fully analyzed, this is precisely the analogue for $`H`$ of Theorem 3.2 and Corollary 3.3 above. This suggests that we compare the results in general.
On the one hand we compute
$$\mathrm{log}(𝐄((\alpha _+\alpha _{})^{\lambda _1(p_+^n,p_{}^n;\alpha _+,\alpha _{})}))=2\mu n^{2/3}(\frac{w_+}{\sigma _+}+\frac{w_{}}{\sigma _{}})+\mathrm{log}(𝐄(\mathrm{exp}(2w_+X\sigma /\sigma _+)\mathrm{exp}(2w_{}X\sigma /\sigma _{}))),$$
(5.13)
with $`X`$ distributed according to $`H(w_+,w_{})`$ in the limit. Thus to retain the analogy, we must have $`\sigma _+=\sigma _{}=\sigma `$. On the other hand, we have:
$$\begin{array}{c}\mathrm{log}(𝐄((\alpha _+\alpha _{})^{\lambda _1(p_+^n,p_{}^n;\alpha _+,\alpha _{})}))=\hfill \\ \hfill 2n^{2/3}(\theta g)(\stackrel{~}{\alpha }_+)\frac{w_++w_{}}{\sigma }+2n^{1/3}(\theta ^2g)(\stackrel{~}{\alpha }_+)\frac{w_{}^2w_+^2}{\sigma ^2}+\frac{4}{3}(\theta ^3g)(\stackrel{~}{\alpha }_+)\frac{w_{}^3+w_+^3}{\sigma ^3}+O(n^{1/3}).\end{array}$$
(5.14)
where $`(\theta f)(z)=z\frac{d}{dz}f(z)`$ and $`g(z)=\mathrm{log}E(z;p_{})\mathrm{log}E(z^1;p_+)`$. Thus, in particular, $`(\theta g)(z)=m(z;p_+,p_{})`$. Comparing the asymptotics, we find
$$\sigma =((\theta ^3g)(\stackrel{~}{\alpha }_+)/2)^{1/3}=(\stackrel{~}{\alpha }_+^2m^{\prime \prime }(\stackrel{~}{\alpha }_+;p_+,p_{})/2)^{1/3}.$$
(5.15)
Similar considerations (based on part (iv) below) give us the scale factors for $`\alpha _+>\stackrel{~}{\alpha }_+`$, thus giving us the following conjecture (see for the definitions of the limiting distributions):
###### Conjecture 5.4.
Fix parameters $`p_+`$, $`p_{}`$, define $`\stackrel{~}{\alpha }_\pm `$ as above, and further define
$`\mu `$ $`=m(\stackrel{~}{\alpha }_+;p_+,p_{})`$ $`\sigma `$ $`=(\stackrel{~}{\alpha }_+^2m^{\prime \prime }(\stackrel{~}{\alpha }_+;p_+,p_{})/2)^{1/3}`$ (5.16)
$`\mu _+(z)`$ $`=m(z;p_+,p_{})`$ $`\sigma _+(z)`$ $`=(zm^{}(z;p_+,p_{}))^{1/2},`$ (5.17)
$`\mu _{}(z)`$ $`=m(z;p_{},p_+)`$ $`\sigma _{}(z)`$ $`=(zm^{}(z;p_{},p_+))^{1/2}.`$ (5.18)
Assume that $`\stackrel{~}{\alpha }_\pm \{0,\mathrm{}\}`$, $`\sigma >0`$, $`\sigma _+(\alpha _+)^2>0`$ for $`\stackrel{~}{\alpha }_+<\alpha _+<R(p_{})^1`$, and $`\sigma _{}(\alpha _{})^2>0`$ for $`\stackrel{~}{\alpha }_{}<\alpha _{}<R(p_+)^1`$. Then we have the following limiting distributions as $`n\mathrm{}`$:
* (i) If $`0\alpha _+<\stackrel{~}{\alpha }_+`$ and $`0\alpha _{}<\stackrel{~}{\alpha }_{}`$ are fixed, then
$$\frac{\lambda _1(p_+^n,p_{}^n;\alpha _+,\alpha _{})\mu n}{\sigma n^{1/3}}F_{\text{GUE}}$$
(5.19)
Near the critical point, set $`w_\pm `$ by $`\alpha _\pm =\stackrel{~}{\alpha }_\pm \mathrm{exp}(2w_\pm /\sigma n^{1/3})`$.
* (ii) If $`w_\pm `$ and $`0\alpha _{}<\stackrel{~}{\alpha }_{}`$ are fixed,
$$\frac{\lambda _1(p_+^n,p_{}^n;\alpha _+,\alpha _{})\mu n}{\sigma n^{1/3}}G(w_\pm ).$$
(5.20)
* (iii) If $`w_+`$ and $`w_{}`$ are fixed,
$$\frac{\lambda _1(p_+^n,p_{}^n;\alpha _+,\alpha _{})\mu n}{\sigma n^{1/3}}H(w_+,w_{}).$$
(5.21)
Finally (the Gaussian regime), let $`\alpha _+^0`$ and $`\alpha _{}^0`$ be fixed such that $`\alpha _\pm ^0>\stackrel{~}{\alpha }_\pm `$ and $`\mu _+(\alpha _+^0)=\mu _{}(\alpha _{}^0)`$, and set $`\alpha _\pm =\alpha _\pm ^0\mathrm{exp}(x_\pm /\sigma _\pm (\alpha _\pm ^0)n^{1/2})`$.
* (iv) If $`x_\pm `$ and $`0\alpha _{}<\alpha _{}^0`$ are fixed,
$$\frac{\lambda _1(p_+^n,p_{}^n;\alpha _+,\alpha _{})\mu _\pm (\alpha _\pm ^0)n}{n^{1/2}}N(x_\pm ,\sigma _\pm (\alpha _\pm ^0)^2).$$
(5.22)
* (v) If $`x_+`$ and $`x_{}`$ are fixed,
$$\frac{\lambda _1(p_+^n,p_{}^n;\alpha _+,\alpha _{})\mu _\pm (\alpha _\pm ^0)n}{n^{1/2}}\mathrm{max}(N(x_+,\sigma _+(\alpha _+^0)^2),N(x_{},\sigma _{}(\alpha _{}^0)^2)).$$
(5.23)
###### Remark 1.
This, in turn, is analogous to the central limit theorem, so again can probably be strengthened considerably (although not nearly to the same extent as Conjecture 5.2 most likely can). In particular, it is presumably sufficient for the parameters $`\alpha _\pm `$, $`w_\pm `$, $`x_\pm `$ to tend to limits as appropriate, rather than simply be fixed.
###### Remark 2.
We recall that
$$\mu _+(\alpha )=𝐄(N_+)+𝐄(N_{})$$
(5.24)
at $`\alpha _+=\alpha `$, $`\alpha _{}=\alpha ^1`$. Similarly,
$$\sigma _+(\alpha )=\mathrm{Var}(N_+)\mathrm{Var}(N_{})$$
(5.25)
###### Remark 3.
The analogous conjectures for models of the other symmetry types (, , ) are straightforward. We note in particular that when $`p_+=p_{}=p`$, we find $`\alpha m^{}(\alpha )=\alpha ^1m^{}(\alpha ^1)`$, and $`\alpha m^{}(\alpha )>0`$ whenever $`\alpha >1`$. So the hypotheses of the above conjectures hold in such cases, with $`\stackrel{~}{\alpha }_\pm =1`$.
In the continuous limit, we make a similar conjecture; the main difference is that the model is nonincreasing in the parameters, not nondecreasing, so the $`F_{\text{GUE}}`$ region is now $`a_\pm >\stackrel{~}{a}_\pm `$. The scale factors are:
$`\mu `$ $`=m_c(\stackrel{~}{a}_+;\rho ^+,u:\rho ^{})`$ $`\sigma `$ $`=(m_c^{\prime \prime }(\stackrel{~}{a}_+;\rho ^+,u:\rho ^{})/2)^{1/3}`$ (5.26)
$`\mu _+(z)`$ $`=m_c(z;\rho ^+,u:\rho ^{})`$ $`\sigma _+(z)^2`$ $`=m_c^{}(z;\rho ^+,u:\rho ^{}),`$ (5.27)
$`\mu _{}(z)`$ $`=\mu _+(z)`$ $`\sigma _{}(z)^2`$ $`=\sigma _+(z)^2`$ (5.28)
where we define
$$m_c(z;\rho ^+,u:\rho ^{})=uz+\underset{i}{}\frac{1}{\rho _i^{}+a}+\underset{i}{}\frac{1}{\rho _i^+a}$$
(5.29)
Near the critical point, we take $`a^\pm =\stackrel{~}{a}^\pm +2w^\pm /\sigma n^{1/3}`$, while in the Gaussian regime, we take $`a^\pm =a_0^\pm x^\pm /\sigma _\pm (a_0^\pm )n^{1/2}`$.
As remarked above, Conjecture 5.4 was proved in (with the exception of parts (iv) and (v), which are straightforward using the argument in section 7 of ) for the cases $`p_\pm =t:/`$ and $`p_\pm =/\sqrt{q}^n`$. The only other known results are for the case $`\alpha _\pm =0`$; the references in the following examples refer to this case alone.
###### Example 1.
If we take $`p_\pm =t:/`$, we have
$$\mu _\pm (z)=(z+z^1)t,\sigma _\pm (z)=((zz^1)t)^{1/2},\stackrel{~}{\alpha }=1,\mu =2t,\sigma =t^{1/3}.$$
(5.30)
This corresponds to the classical case of increasing subsequences of random permutations, studied in .
###### Example 2.
If we take $`p_\pm =/(\sqrt{q})^{n_\pm }`$, with $`n_+/n_{}`$ tending to a constant, we have
$$\mu _+(z)=\frac{\sqrt{q}zn_+}{1z\sqrt{q}}+\frac{\sqrt{q}n_{}}{z\sqrt{q}},\sigma _+(z)^2=\frac{\sqrt{q}zn_+}{(1z\sqrt{q})^2}\frac{\sqrt{q}zn_{}}{(z\sqrt{q})^2}$$
(5.31)
and thus
$$\stackrel{~}{\alpha }_+=\frac{\sqrt{qn_+}+\sqrt{n_{}}}{\sqrt{qn_{}}+\sqrt{n_+}},\mu =\frac{q(n_++n_{})+2\sqrt{qn_+n_{}}}{1q},\sigma =\frac{(qn_+n_{})^{1/6}(1+\sqrt{\frac{qn_+}{n_{}}})^{2/3}(1+\sqrt{\frac{qn_{}}{n_+}})^{2/3}}{1q}.$$
(5.32)
This model was analyzed in , along with the continuous (Laguerre) limit, in which case we have
$$\mu _+(z)=\frac{1}{2}[\frac{n_+}{1+2z}+\frac{n_{}}{12z}],\sigma _+(z)^2=\frac{4n_+}{(1+2z)^2}\frac{4n_{}}{(12z)^2}$$
(5.33)
and thus
$$\stackrel{~}{a}_+=\frac{\sqrt{n_+}\sqrt{n_{}}}{2(\sqrt{n_+}+\sqrt{n_{}})},\mu =(\sqrt{n_+}+\sqrt{n_{}})^2,\sigma =(n_+n_{})^{1/6}(\sqrt{n_+}+\sqrt{n_{}})^{4/3}.$$
(5.34)
###### Example 3.
If we take $`p_+=/1^n`$, $`p_{}=t:/`$, we have
$$\mu _+(z)=zt+\frac{n}{z1},\sigma _+(z)^2=zt\frac{zn}{(z1)^2},$$
(5.35)
and thus
$$\stackrel{~}{\alpha }_+=1+\sqrt{n/t},\mu =t+2\sqrt{tn},\sigma =\sqrt{t}n^{1/6}(1+\sqrt{n/t})^{2/3}.$$
(5.36)
This model, corresponding to weakly increasing subsequences of random words, was studied in and . In the continuous limit (corresponding to the $`n\times n`$ GUE ), we have:
$$\mu =2\sqrt{n},\sigma =n^{1/6}$$
(5.37)
These are precisely the scale factors required to make the largest eigenvalue of an $`n\times n`$ Gaussian Hermitian matrix tend to the limit $`F_{\text{GUE}}`$.
###### Example 4.
If we take $`p_+=(\sqrt{q})^{n_+}/`$, $`p_{}=(\sqrt{q})^n_{}/`$, with $`n_+/n_{}`$ tending to a constant, we have
$$\mu _+(z)=\frac{\sqrt{q}zn_+}{1+z\sqrt{q}}+\frac{\sqrt{q}n_{}}{z+\sqrt{q}}\sigma _+(z)^2=\frac{\sqrt{q}zn_+}{(1+z\sqrt{q})^2}\frac{\sqrt{q}zn_{}}{(z+\sqrt{q})^2}.$$
(5.38)
Here we have three cases. If $`qn_+n_{}`$, then $`\stackrel{~}{\alpha }_+=0`$, and if $`qn_{}n_+`$, then $`\stackrel{~}{\alpha }_+=\mathrm{}`$; in either case, the above conjectures do not apply (indeed, in those cases one expects the limiting distribution to be atomic, $`\lambda _1(p_+,p_{})=\mathrm{min}(n_+,n_{})`$). Otherwise,
$$\stackrel{~}{\alpha }_+=\frac{\sqrt{n_{}}\sqrt{qn_+}}{\sqrt{n_+}\sqrt{qn_{}}},\mu =\frac{2\sqrt{qn_+n_{}}q(n_++n_{})}{1q},\sigma =\frac{(qn_+n_{})^{1/6}(1\sqrt{\frac{qn_+}{n_{}}})^{2/3}(1\sqrt{\frac{qn_{}}{n_+}})^{2/3}}{1q}.$$
(5.39)
The mean in this model was derived in ; the refined asymptotics of a symmetrized version was studied in .
###### Example 5.
If we take $`p_+=(\sqrt{q})^{n_+}/`$, $`p_{}=/(\sqrt{q})^n_{}`$, with $`n_+/n_{}`$ tending to a constant, we have
$$\mu _+(z)=\frac{\sqrt{q}zn_+}{1z\sqrt{q}}+\frac{\sqrt{q}n_{}}{z+\sqrt{q}},\sigma _+(z)^2=\frac{\sqrt{q}zn_+}{(1z\sqrt{q})^2}\frac{\sqrt{q}zn_{}}{(z+\sqrt{q})^2}.$$
(5.40)
There are two cases. If $`qn_+n_{}`$, then $`\stackrel{~}{\alpha }_+=0`$, and the above remark applies. Otherwise
$$\stackrel{~}{\alpha }_+=\frac{\sqrt{n_{}}\sqrt{qn_+}}{\sqrt{n_+}+\sqrt{qn_{}}},\mu =\frac{2\sqrt{qn_+n_{}}+q(n_{}n_+)}{1+q},\sigma =\frac{(qn_+n_{})^{1/6}(1\sqrt{\frac{qn_+}{n_{}}})^{2/3}(1+\sqrt{\frac{qn_{}}{n_+}})^{2/3}}{1+q}.$$
(5.41)
The mean in this model was first derived in ; the fluctuations have been analyzed in section 5 of .
###### Example 6.
If we take $`p_+=1^n/`$, $`p_{}=t:/`$, we have
$$\mu _+(z)=zt+\frac{n}{z+1},\sigma _+(z)^2=zt\frac{zn}{(z+1)^2}$$
(5.42)
If $`nt`$, then $`\stackrel{~}{\alpha }_+=0`$; otherwise we have
$$\stackrel{~}{\alpha }_+=\sqrt{n/t}1,\mu =2\sqrt{tn}t,\sigma =(tn)^{1/6}(1\sqrt{t/n})^{2/3}.$$
(5.43)
This corresponds to strictly increasing subsequences of random words, and was studied to a small extent in .
Note that the pathological case $`\stackrel{~}{\alpha }_+=0`$ (resp. $`\stackrel{~}{\alpha }_+=\mathrm{}`$) can only occur if all the rows (resp. columns) are strict.
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# A new integrable system of symmetrically coupled derivative nonlinear Schrödinger equations via the singularity analysis
(April 13, 2000)
## Abstract
A new integrable system of two symmetrically coupled derivative nonlinear Schrödinger equations is detected by means of the singularity analysis. A nonlinear transformation is proposed which uncouples the equations of the new system.
In this paper, we study the integrability of the following system of two symmetrically coupled derivative nonlinear Schrödinger equations:
$$\begin{array}{c}q_t=\mathrm{i}q_{xx}+aq\overline{q}q_x+bq^2\overline{q}_x+cr\overline{r}q_x+dq\overline{r}r_x+eqr\overline{r}_x+\mathrm{i}fq^2\overline{q}+\mathrm{i}gqr\overline{r},\hfill \\ r_t=\mathrm{i}r_{xx}+ar\overline{r}r_x+br^2\overline{r}_x+cq\overline{q}r_x+dr\overline{q}q_x+erq\overline{q}_x+\mathrm{i}fr^2\overline{r}+\mathrm{i}grq\overline{q},\hfill \end{array}$$
(1)
where $`a,b,c,d,e,f,g`$ are real parameters, and the bar denotes the complex conjugation. By means of the singularity analysis, we detect one new integrable case of the system (1), characterized by the conditions
$$a=c=e0,b=d=g=0.$$
(2)
Then we propose a nonlinear transformation, which uncouples the equations (1) in the case (2).
We follow the Weiss-Kruskal algorithm of the singularity analysis , . With respect to $`q,\overline{q},r,\overline{r}`$, which should be considered as mutually independent, the system (1) is a normal system of four second-order equations, of total order eight. A hypersurface $`\varphi (x,t)=0`$ is non-characteristic for (1) if $`\varphi _x0`$, and we set $`\varphi _x=1`$. The substitution of the expansions
$$\begin{array}{c}q=q_0(t)\varphi ^\alpha +\mathrm{}+q_n(t)\varphi ^{n+\alpha }+\mathrm{},\hfill \\ \overline{q}=\overline{q}_0(t)\varphi ^\beta +\mathrm{}+\overline{q}_n(t)\varphi ^{n+\beta }+\mathrm{},\hfill \\ r=r_0(t)\varphi ^\gamma +\mathrm{}+r_n(t)\varphi ^{n+\gamma }+\mathrm{},\hfill \\ \overline{r}=\overline{r}_0(t)\varphi ^\delta +\mathrm{}+\overline{r}_n(t)\varphi ^{n+\delta }+\mathrm{}\hfill \end{array}$$
(3)
(the bar does not mean the complex conjugation now) into the system (1) determines the branches, i.e. admissible choices of $`\alpha ,\beta ,\gamma ,\delta `$ and $`q_0,\overline{q}_0,r_0,\overline{r}_0`$, as well as the positions $`n`$ of resonances for those branches. We require that the system (1) admits a singular generic branch, where at least one of the exponents $`\alpha ,\beta ,\gamma ,\delta `$ is negative, the number of resonances is eight, and seven of them lie in nonnegative positions. For this singular generic branch, there are three possibilities: (i) $`\alpha +\beta =\gamma +\delta =1`$, (ii) $`\alpha +\beta =1,\gamma +\delta >1`$, and (iii) $`\alpha +\beta >1,\gamma +\delta =1`$, which is related to (ii) by $`qr,\overline{q}\overline{r}`$. The choice of (i) would lead us to known integrable cases of (1), and we will give a detailed consideration of it in a separate paper. In this paper, we study the possibility (ii), which leads to a new integrable case of (1). Computations are done by using the Mathematica system .
Setting $`\beta =1\alpha ,\gamma +\delta >1`$ in (3), and assuming without loss of generality that $`q_0\overline{q}_0=\mathrm{i}`$, we obtain from (1) the following four cases:
$$\begin{array}{c}a=\frac{25\alpha 5\alpha ^2}{1+2\alpha },b=\frac{3\alpha 3\alpha ^2}{1+2\alpha },\\ c=\frac{(1+2\alpha )e+(1+\alpha )(\gamma \gamma ^2)\alpha (\delta \delta ^2)}{(1+\alpha )\gamma +\alpha \delta },d=\frac{(\alpha \gamma +(1+\alpha )\delta )e+\gamma \delta (2\gamma \delta )}{(1+\alpha )\gamma +\alpha \delta },\\ (n+1)n^3(n2)(n3)(n+\gamma +\delta +\frac{(1+2\alpha )e(2\alpha +\gamma )\delta }{(1+\alpha )\gamma +\alpha \delta })\times \\ (n+\gamma +\delta 2\frac{(1+2\alpha )e(2\alpha +\gamma )\delta }{(1+\alpha )\gamma +\alpha \delta })=0;\end{array}$$
(4)
$$\begin{array}{c}a=\frac{25\alpha 5\alpha ^2}{1+2\alpha },b=\frac{3\alpha 3\alpha ^2}{1+2\alpha },\\ c=\frac{(1+\alpha )d}{\delta }\frac{1+2\alpha +2\alpha ^2}{1+2\alpha }+\frac{(1+3\alpha +3\alpha ^2)\delta }{(1+\alpha )(1+2\alpha )},\\ e=\frac{\alpha \delta (2+2\alpha \delta )}{(1+\alpha )(1+2\alpha )},\gamma =\frac{\alpha \delta }{1+\alpha },\\ (n+1)n^3(n2)(n3)(n\frac{(1+\alpha )d}{\delta }+\frac{2\alpha ^2}{1+2\alpha }+\frac{(1+3\alpha +\alpha ^2)\delta }{(1+\alpha )(1+2\alpha )})\times \\ (n+\frac{(1+\alpha )d}{\delta }\frac{2(1+\alpha )^2}{1+2\alpha }+\frac{(1+\alpha \alpha ^2)\delta }{(1+\alpha )(1+2\alpha )})=0;\end{array}$$
(5)
$$\begin{array}{c}a=\frac{25\alpha 5\alpha ^2}{1+2\alpha },b=\frac{3\alpha 3\alpha ^2}{1+2\alpha },d=e=\gamma =\delta =0,\\ (n+1)n^3(n2)(n3)(n1+c)(n1c)=0;\end{array}$$
(6)
$$\begin{array}{c}a=2,b=e=0,c=\frac{d}{\gamma }+1\gamma ,\alpha =1,\delta =0,\\ (n+1)n^3(n2)(n3)(n+\frac{d}{\gamma }+\gamma )(n\frac{d}{\gamma }2+\gamma )=0.\end{array}$$
(7)
In the case (4), taking into account that the equation for $`n`$ should admit four positive integer solutions and that $`\gamma ,\delta `$ should be integers, we have to set
$$e=\delta \frac{\delta ^2}{1+2\alpha },\gamma =\delta .$$
(8)
Now the positions of resonances are $`n=1,0,0,0,1,1,2,3`$. We obtain from (1) and (3) the recursion relations for $`q_n,\overline{q}_n,r_n,\overline{r}_n`$, $`n=0,1,\mathrm{}`$, and then check the consistency of those relations at the resonances. At $`n=1`$, we have to set
$$g=0,\delta =12\alpha ,\alpha (\alpha +1)=0,$$
(9)
otherwise logarithmic terms should be introduced into the expansions (3). With (9), the recursion relations become consistent at $`n=2,3`$ as well. Both choices of $`\alpha `$, $`\alpha =1`$ and $`\alpha =0`$, lead through (9), (8), (4) and scaling of variables to the case (2) of the system (1).
In the cases (5) and (7), it is impossible to have four resonances in positive integer positions, if $`\gamma ,\delta `$ are integers.
In the case (6), we should set $`c=0`$ to have correct positions of resonances. Then the consistency of recursion relations at $`n=1`$ requires $`g=0`$, and the system (1) becomes a pair of non-coupled equations.
As we see, the system (1) admits a good singular generic branch in the new case (2). Setting $`a=2`$ without loss of generality and making the transformation
$$\begin{array}{c}q=q^{}(x^{},t^{})\mathrm{exp}(\mathrm{i}\omega ),\overline{q}=\overline{q}^{}(x^{},t^{})\mathrm{exp}(\mathrm{i}\omega ),\\ r=r^{}(x^{},t^{})\mathrm{exp}(\mathrm{i}\omega ),\overline{r}=\overline{r}^{}(x^{},t^{})\mathrm{exp}(\mathrm{i}\omega ),\\ \omega =\frac{f}{2}x^{}+\frac{f^2}{4}t^{},x=x^{}ft^{},t=t^{},\end{array}$$
(10)
we obtain the simplified form of the new case:
$$\begin{array}{c}q_t=\mathrm{i}q_{xx}+2(q\overline{q}q_x+r\overline{r}q_x+qr\overline{r}_x),\hfill \\ r_t=\mathrm{i}r_{xx}+2(r\overline{r}r_x+q\overline{q}r_x+rq\overline{q}_x),\hfill \end{array}$$
(11)
where the prime of the new variables is omitted. Besides the singular generic branch, already considered, and the Taylor expansions, governed by the Cauchy-Kovalevskaya theorem, we have to study the following two branches, admitted by the system (11):
$$\begin{array}{c}q_0\overline{q}_0=\mathrm{i},r_0\overline{r}_0=\mathrm{i},\alpha =\gamma =1,\beta =\delta =2,\\ (n+1)^2n^2(n2)^2(n3)^2=0;\end{array}$$
(12)
$$\begin{array}{c}q_0\overline{q}_0=\mathrm{i},r_0\overline{r}_0=\mathrm{i},\alpha =\delta =1,\beta =\gamma =0,\\ (n+1)^2n^2(n2)^2(n3)^2=0.\end{array}$$
(13)
The recursion relations turn out to be consistent at all resonances of the branches (12) and (13). Consequently, the system (11), or (1) with (2), has passed the Painlevé test for integrability well.
Using the Mathematica package condens.m , we can check that the system (11) has two conservation laws for each rank from 1 to (at least) 4. The fact, that the conservation laws appear by pairs, is highly suggestive that the equations of (11) can be uncoupled by some transformation. The first conservation laws of (11) are given by
$$\begin{array}{c}(q\overline{q})_t=\{\mathrm{i}(q_x\overline{q}q\overline{q}_x)+(q\overline{q})^2+2q\overline{q}r\overline{r}\}_x,\hfill \\ (r\overline{r})_t=\{\mathrm{i}(r_x\overline{r}r\overline{r}_x)+(r\overline{r})^2+2r\overline{r}q\overline{q}\}_x.\hfill \end{array}$$
(14)
We introduce a new set of dependent variables by
$$q=u\mathrm{exp}(\mathrm{i}_{x_0}^xv\overline{v}dx^{}),r=v\mathrm{exp}(\mathrm{i}_{x_0}^xu\overline{u}dx^{}).$$
(15)
Using the first conservation laws (14) and the definition of new variables (15), we find that the system (11) is equivalent to two independent Chen-Lee-Liu equations :
$$u_t=\mathrm{i}u_{xx}+2u\overline{u}u_x,v_t=\mathrm{i}v_{xx}+2v\overline{v}v_x.$$
(16)
Here we have assumed that the dependent variables approach zero as $`xx_0`$. The transformation (15) of (11) into (16) proves the integrability of the new case (2) of the system (1).
Acknowledgments. The work of S. Yu. S. was supported in part by the Fundamental Research Fund of Belarus, grant $`\mathrm{\Phi }`$98-044. The work of T. T. was supported by a JSPS Research Fellowship for Young Scientists.
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# The 𝛾 Velorum binary system Based on observations collected at the European Southern Observatory at La Silla, Chile. ESO proposals Nrs. 56.D-327, 57.D-517 and 56.D-0700
## 1 Introduction
The double-lined spectroscopic binary $`\gamma `$ Velorum (WC8 + O7.5) can be used to derive the mass of the Wolf-Rayet star, a quantity which cannot be directly determined in the case of single objects, but which provides an important constraint in computations of the late stages of massive star evolution. From it, and the mass-luminosity relationship for WR stars (Schaerer & Maeder 1992) a value for the WR luminosity can be derived and compared to that derived from spectroscopic analyses.
Orbital parameters together with the mass ratio were determined by Schmutz et al. (1997), while the 7.3$`\sigma `$ hipparcos distance allowed Schaerer et al. (1997) to derive the system’s total mass ($``$(O+WR)=29.5$`\pm `$15.9 M), the absolute visual magnitude ($`M_V`$(O+WR)=$`5.39`$) and to carry out a preliminary spectroscopic analysis of both components. Schaerer et al. (1997), however, modelled only the WR He i $`\lambda `$4471 / He ii $`\lambda `$4541 line strength ratio, assuming the mass-loss rate and the carbon abundance. Additionally, line blanketing was not accounted for by their model and their calculation used simple helium and carbon model atoms. Their stellar parameter determination should therefore be considered preliminary.
De Marco & Schmutz (1999; hereafter Paper I) carried out a full spectroscopic analysis of the O star through which its parameters were determined simultaneously with the light ratio between the WR and O stellar components. From a solution of the hydrodynamic wind equations, the mass of the O star was also derived and found to be consistent with the evolutionary O star mass. From it and the mass ratio of Schmutz et al. (1997), the mass of the WR star could be calculated (M=9.0$`\pm `$1.0 M). The combined mass of the system was therefore derived to be 39 M, higher than the value of Schaerer et al. (1997), but consistent within their uncertainty. From the mass of the WR star and the theoretical mass-luminosity relation of Schaerer & Maeder (1992), a value for the WR bolometric luminosity was derived (log($`L`$/L) = 5.2), significantly higher than the spectroscopic luminosity derived by Schaerer et al. (1997; log($`L`$/L) = 4.8). Finally, using the synthetic O star spectrum, the WR11 spectrum was de-convolved from the O star component allowing the WR emission lines to be measured with greater accuracy.
In this paper we determine the stellar parameters of the WR star using the clumped, line-blanketed stellar atmosphere model cmfgen (Hillier 1987, Hillier & Miller 1998). Ultimately, higher spectroscopic luminosities would lower the wind performance number, indicating that a smaller number of photon scatterings is needed to achieve the radiative driving of the wind. This would facilitate the derivation of the wind velocity law. $`\gamma `$ Velorum provides us with a new tool to test the model luminosity with an independent method in such a way as to impose a further check on the wind momentun number.
A second aim of this paper is to continue the comparison of the non-LTE stellar atmosphere code cmfgen and the independently developed Sobolev approximation code isa-wind (de Koter et al. 1993, 1997), which is used in combination with the line-blanketing Monte Carlo code of Schmutz (1991). This study follows from the comparison carried out by Crowther et al. (1999) for the WN8 star WR124.
Finally, Schmutz (1997) showed that line-line interaction between He ii Ly$`\alpha `$ and metal lines at similar wavelengths can lead to a larger spectral luminosity. The inclusion of this interaction in the models, an effect that he called “photon loss”, showed that WR6 might have a larger luminosity than previously thought, and allowed him to calculate a velocity law for its outer wind. Since line-line interaction is included in the cmfgen model code, we can test whether the approximation of Schmutz (1997) is valid and quantify the importance of the photon loss effect in the case of WR11.
In Sect. 2 we summarise the observations, while in Sect. 3 we discuss the light ratio between the O and WR stars and the reddening towards the system. In Sect. 4 we describe the model atmosphere codes used in the determination of the stellar parameters. We present our results in Sect. 5, together with the code comparison, while in Sect. 6 the effect of photon loss is studied. We finally draw our conclusions in Sect. 7.
## 2 Observations and reduction
Our optical observations of $`\gamma `$ Velorum are the same as those used by Schmutz et al. (1997), where a full account of the data reduction can be found. In summary, several spectra have been obtained at the 50 cm ESO telescope in conjunction with the Heidelberg Extended Range Optical Spectrograph (heros) in the ranges 3500-5500 Å and 5800-8600 Å covering the binary period. The resolution is R=20 000, while the S/N ratio ranges between 100 at 3600 Å and 250 at 6000 Å. The phase-averaged spectrum of the $`\gamma `$ Velorum system was rectified as described in Paper I. It is important to remind that rectification of WC stars is a very delicate operation because of blending of large numbers of broad lines, which leave few identifiable continuum points. Moreover, spectral ranges at the end of instrumental orders (5300-5500 Å and 5800-5900 Å), suffer from additional uncertainty. Unfortunately this affects the diagnostic lines He i $`\lambda `$5876, He ii $`\lambda `$5412, and C iv $`\lambda `$5471. Extreme care was taken in rectifying these ranges and consistency between these lines and those residing at other wavelengths was insured.
Observations of the 10830 Å He i line were obtained at the European Southern Observatory (ESO) New Technology Telescope on January 12 and 13, 1996, with the emmi instrument. emmi was used in the remd mode with grating No. 7, a slit-width of 1.5$`^{^{\prime \prime }}`$ and the spectrum recorded on ESO CCD No. 36. Two exposures, each lasting 10 s, were taken. Wavelength calibration was obtained with respect to a ThAr arc lamp, to achieve a resolution of 2 Å. Relative flux calibration was obtained with respect to a 180 s exposure of the standard star $`\theta `$ Crt.
Mid-IR data of WR11 were obtained as part of Guest Observer programme WRSTARS (P.I. van der Hucht), with the Short Wavelength Spectrograph (SWS; de Graauw et al. 1996) on board the ESA Infrared Space Observatory (ISO; Kessler et al. 1996). The SWS AOT6 observing mode was used to achieve full grating resolution, $`\lambda /\mathrm{\Delta }\lambda \mathrm{\hspace{0.33em}1300}2500`$, with the wavelength coverage 2.38–45.0$`\mu `$m. The observations were flux calibrated to an accuracy of 5% and $``$20% for the low and high wavelength ends, respectively. The possibility that the nearby $`\gamma ^1`$ Velorum might have contaminated the ISO observations of $`\gamma `$ Velorum was investigated, but with an aperture of 20<sup>′′</sup> centered on $`\gamma `$ Velorum it is not possible that light from $`\gamma ^1`$ Velorum, at 40<sup>′′</sup> distance, could have entered the aperture. For full details on the observations and data reduction procedures, we refer the reader to Morris et al. (2000).
Finally, we have obtained a flux-calibrated UV dataset for WR11 as follows. WR11 was too bright for absolutely flux calibrated observations with IUE/LORES. Consequently we have matched IUE observations from July 1988, to low resolution S2-68 observations from 1973 at phase 0.3 (Willis & Wilson 1976).
## 3 Light ratio correction and reddening
In this section we discuss the O-WR light ratio of $`\gamma `$ Velorum. We also re-evaluate the interstellar reddening that was adopted in Paper I, namely $`E(bv)=0.03`$ mag, obtained from assumed intrinsic colours (van der Hucht et al. 1988).
The spectrum of $`\gamma `$ Velorum used here is the same phase-averaged rectified spectrum used in Paper I. Correction for the continuum as well as the line contribution of the O star was carried out in that paper using the model O star spectrum and the deduced value of the light ratio. A light ratio of $`L`$(O) / $`L`$(O+WR) = 0.795 $`\pm `$0.020 was obtained in Paper I from an average of three lines, namely He i $`\lambda `$4471, He i $`\lambda `$4922 and He ii $`\lambda `$4541. Assuming that this light ratio is representative for this spectral range, we derive the wavelength dependence of the light ratio using synthetic spectra for the O star from Paper I, and for the WR star from the following section. The ratio of these continuum distributions, adjusted such that $`L_\lambda `$(O) / $`L_\lambda `$(WR) = 3.88$`{}_{0.4}{}^{}{}_{}{}^{+0.5}`$ at $`\lambda `$= 4700 Å, is plotted in Fig. 1, together with error bars. A calibration curve $`L_\lambda `$(O+WR) / $`L_\lambda `$(WR) is created and applied to the optical spectrum of WR11 to account for the contribution of the O star continuum to the WR equivalent widths. In Fig. 2 we show the effect that the error bars from the light ratio have on He ii $`\lambda `$5412 and C iv $`\lambda `$5471.
After the spectrum is rectified and light-ratio corrected, a model can be calculated to fit the WR lines, as discussed in the following section. First, we preempt the result of this modelling in order to verify our choice of reddening. Published broad-band Johnson V magnitudes for $`\gamma `$ Velorum span a wide range, from V = 1.70 (Cousins 1972) to V = 1.83 (Johnson et al. 1966). We suspect that the broad band observations suffer from contaminations by nearby spectral lines such as C iv $`\lambda `$5806 and that the measured brightness depends crucially on the exact sensitivity curve of the used setup. Therefore, in order to minimise the contamination of WR line emission, we prefer to adopt the narrow-band $`ubv`$ Smith (1968) photometry. Following the corrections by Schmutz & Vacca (1991) to the Smith datasets, we obtain $`v=1.70`$ mag and $`bv=0.32`$ mag. Our theoretical combined O+WR model reveals $`(bv)_0=0.33`$ mag, implying $`E(bv)=0.01`$ or $`E(BV)`$ = 0.01 mag.
We have also carried out an alternative reddening determination, by fitting the combined WR+O model to UV (IUE/LORES) and IR (ISO) spectrophotometry, de-reddened following Seaton (1979). Fig. 3 compares observed and de-reddened spectrophotometry (solid lines) with theoretical models (dotted lines), revealing $`E(bv)=0.06`$ or $`E(BV)`$ = 0.07 mag. Unfortunately, we are unable to reconcile these discrepant reddening values, and so adopt $`E(bv)=0.035`$ mag ($`E(BV)`$ = 0.04), in good agreement with van der Hucht et al. (1988). Note that for this reddening and the adopted light ratio, the combined O+WR model flux lies $``$10% higher than the UV IUE observed fluxes and a similar value below the IR ISO observations.
The distance modulus of 7.06 mag (corresponding to a distance of \[$`258_{31}^{+41}`$\] pc, hipparcos<sup>1</sup><sup>1</sup>1The hipparcos distance to $`\gamma `$Velorum has been recently put in doubt (Pozzo et al. 2000) by the finding of a new population of pre-main sequence stars which might be associated with the binary. This new value of the distance re-positions $`\gamma `$ Velorum within the Vela OB2 association at a distance of 360-490 pc, in agreement with earlier estimates. Clearly, adopting a larger distance would have a dramatic effect on distance dependent parameters, such as luminosity, radius and mass-loss (for a distance of 360-490 pc, $`\mathrm{\Delta }\mathrm{log}(L/\mathrm{L}_{})`$=0.3-0.6). We have, however, maintained the 7.3$`\sigma `$ hipparcos distance measurement, which we trust to be the best distance estimate to date. and de-reddened $`v`$ = 1.56 mag imply $`M_v`$(O+WR) = $``$5.50. The O/WR continuum light ratio is 3.61 at the $`v`$-band (Fig. 1) so that $`M_v`$(O)=–5.23 and $`M_v`$(WR) = –3.84. We may convert the Smith (1968) $`v`$ to Johnson $`V`$ via $`V=v0.365(bv)+0.007`$, adapted from Turner (1982) to take account of the $`ubv`$ re-calibration by Schmutz & Vacca (1991). From theoretical continuum O and WR distributions we obtain $`(bv)=0.34`$ and $`(bv)=0.20`$, respectively, revealing absolute line free V-magnitudes listed in Table 1. Relative to Paper I, the difference in the absolute magnitude of the O star, $`M_V`$(O), is $`0.04`$ magnitude. The errors listed in Table 1 are solely from the reddening and light ratio uncertainties. Errors in the distance, would add an additional 0.3 mag in the M<sub>V</sub> uncertainties. Propagating these errors onto the derived magnitudes and luminosities in Table 1, and subsequently onto the spectroscopic properties in Sect. 5, introduces an additional systematic error of $`\mathrm{log}(L/L_{})`$=$`\pm `$0.1 for both cases. In this way we lose sight of the actual accuracy reached by this study.
To be consistent with the slight change in the O and WR star $`V`$ magnitudes we have revised the uncertainties on $`L`$, $``$ and $`i`$, while the values of the luminosity and mass of the WR star are also slightly modified.
## 4 Stellar atmosphere codes
In this section we present a summary of the description of both codes used (a full account can be found in Hillier, 1989 and Hillier & Miller, 1998 for cmfgen and in de Koter et al., 1993, 1997, for isa-wind). Along with the description of the basic codes, we give an account of the Monte Carlo line blanketing code used in conjunction with the un-blanketed isa-wind to calculate the additional opacity due to the effect of bound-bound transitions on the continuum photons. Last, the implementation of the photon loss mechanism (Schmutz 1997) into the Sobolev code is presented here, for the first time, and a parallel is drawn with the potential effect of line-line interaction in the line-blanketed code cmfgen.
### 4.1 cmfgen
cmfgen (Hillier 1987, 1990; Hillier & Miller 1998, 1999) solves the transfer equation in the co-moving frame, subject to statistical and radiative equilibrium, assuming an expanding, spherically-symmetric, homogeneous or clumped, time-independent atmosphere. Line blanketing is treated correctly in the transfer problem, except that a simplifying ‘super level’ approach is used (Anderson 1989), by which several levels of similar energies and properties are treated as a single ‘super level’. The population of an individual atomic level in the full model atom is determined by assuming that it has the same departure coefficient as the corresponding super level to which it belongs. For the specific models dealt with in this paper the atomic model is shown in Table 2. Elemental abundances other than hydrogen, helium, carbon and oxygen, are fixed at their solar mass fraction values.
The stellar radius ($`R_{}`$) is defined as the inner boundary of the model atmosphere and is located at Rosseland optical depth of $``$20 with the stellar temperature ($`T_{}`$) defined by the usual Stefan-Boltzmann relation. Similarly, the effective temperature ($`T_{\mathrm{eff}}`$) relates to the radius ($`R_{2/3}`$) at which the Rosseland optical depth equals 2/3. There is now overwhelming evidence for the clumped nature of WR stars (e.g. Moffat 1999), so we have adopted a simple filling factor approach. Following the method proposed by Schmutz (1995), we assume that the wind is clumped with a volume filling factor, $`f`$, and that there is no inter clump material. Since radiation instabilities are not expected to be important in the inner wind we parametrise the filling factor so that it approaches unity at small velocities. Clumped and non-clumped spectra are very similar, except that line profiles are slightly narrower with weaker electron scattering wings in the former. Although non-clumped models can be easily rejected, we are unable to determine the clumping factor because of the severe line blending in WC winds. We can therefore derive only $`\dot{M}/\sqrt{f}`$ with our spectroscopic analysis.
Unfortunately, individual line blanketed co-moving frame calculations are generally computationally demanding, despite the use of super levels, so that a large parameter space cannot be quickly explored. One solution is to (i) solve the transfer problem in the classical Sobolev approximation rather than the co-moving frame, which the code isa-wind does (de Koter et al. 1993, 1997), and (ii) consider line blanketing via Monte Carlo sampling following Schmutz (1991, 1997) allowing the opacity of a huge number of lines to be considered.
### 4.2 isa-wind and the Monte Carlo line-blanketing code
The improved Sobolev approximation code (isa-wind) is described in detail by de Koter et al. (1993, 1997). The principal differences with cmfgen relate to (a) the treatment of the line radiation transfer; (b) wind electron temperature which assumes a grey atmosphere in local thermodynamic equilibrium (LTE); (c) the specific atomic model treated, as listed in Table 2.
The co-moving frame method consistently treats a possible change of properties of the medium inside the region in which line photons can be absorbed and re-emitted. In the Sobolev approximation one assumes this line interaction region is infinitely small. More exactly, the Sobolev approximation assumes that the source function and the opacities are linear functions such that they can be taken out of the transfer integral. This is generally a valid assumption for Wolf-Rayet stars, where the large velocity gradient of the flow results in a small interaction zone. The Sobolev classical approximation (Castor 1970) implies that interactions of line photons with the continuum inside of the (in reality extended) interaction region are ignored. However, the improved Sobolev approximation code isa-wind does include such continuum interactions. Note that use of the Sobolev approximation introduces great simplifications in the radiative transfer, resulting in an overall iteration process up to about 20 times faster compared to the co-moving frame approach.
Turning to line blanketing, an iterative technique including the Monte Carlo method of Schmutz (1991) is employed. The method allows the computation of intensity-weighed effective opacity factors, which account for the presence of tens of thousands of spectral lines, dominated by Fe and Ni. Based on isa-wind atmosphere calculations, the Monte Carlo code determines the line blanketing factors. An iterative procedure is used, such that blanketing factors are used by the non-LTE code to calculate a new atmosphere, which in turn is used to calculate new blanketing factors. A few iterations are generally sufficient. This is due to the fact that the scattering and absorption factors are not very sensitive to the exact model intensity.
The Monte Carlo method deals with the photon scattering in the correct manner, although no branching is taken into consideration, an effect that would tend to re-distribute UV photons to longer wavelengths. The ionization equilibrium of metal species is also approximate, in that it is determined in relation to the non-LTE structure of helium, carbon and oxygen and not by the full solution of the statistical equilibrium equations for that species.
### 4.3 Photon loss
“Photon loss” is the name given by Schmutz (1997) to the interaction between the photon field of the He ii Ly$`\alpha `$ line at 303.78 Å and nearby metal lines (e.g. the O iii Bowen lines at 305.72 and 303.65 Å). This was suggested by Schmutz (1997) to be responsible for a decrease in the ionization equilibrium of the wind of WR6 (WN4), with respect to the equilibrium obtained when no such line-line interaction was accounted for. In order to fit the observed spectrum when photon loss is included in the model, the model’s stellar temperature has to be higher, which results in a higher spectroscopic luminosity. This contributes to a lower wind performance number and allowed Schmutz (1997) to calculate a velocity structure for the outer wind of WR6.
Since line-line interaction was not accounted for in any atmospheric wind code before Hillier & Miller (1998) implemented it into the Hillier (1990) code, Schmutz (1997) calculated the photon loss factor (i.e. the amount of interaction between the He ii Ly$`\alpha `$ and nearby metal lines) by comparing the opacity of metal lines in the 303-Å region to that of the He ii Ly$`\alpha `$ line (his Appendix A). Once the amount of interaction is known for a particular model atmosphere, a corresponding fraction of photons can be taken out of the radiation field at that particular wavelength for all grid points. In this way the same fraction of photons is removed from the radiation field at all depth points, although in reality the photon loss factor depends on the relative opacity of the He ii Ly$`\alpha `$ line and that of the metal lines, which in turn depends on the depth-dependent populations. This problem however is not critical as long as the opacity of metal lines behaves similarly to that of helium. Additionally, the loss of photons is only critical for the point in the atmosphere where the wind is recombining. We implemented the photon loss approximation into the isa-wind code of de Koter et al. (1993 - see Sect. 4.3.1) and used it in the calculation of a model for WR11.
cmfgen includes line-line interaction and as such it should naturally include the photon loss effect, provided the specific metal lines in the He ii Ly$`\alpha `$ region are included in the model atom. We can therefore test whether this effect and its consequences for the stellar luminosity has been accounted for properly by our approximation.
#### 4.3.1 Implementation of Photon loss in isa-wind
Once the photon loss factor, $`f`$, is calculated for a certain model atmosphere (see Schmutz 1997, Appendix A), a fraction $`f`$ of the photons is removed from the radiation field in the wavelength range 300–304Å. To do so, one has to modify the expression where the mean intensity for this line wavelength is calculated at every code iteration. In isa-wind this was achieved via a modification of the Sobolev escape probability.
In the Sobolev approximation the radiation field $`\overline{J}`$ at radius point $`r`$ is related to the line source function $`S`$ in the following way:
$$\overline{J}(r)=(1\beta )S(r)$$
(1)
where $`\beta `$ is the line escape probability by direct flight and where
$$S(r)=\frac{N_uA_{ul}}{N_lB_{ul}N_uB_{ul}}$$
(2)
with $`N_u`$ and $`N_l`$ the upper and lower level populations, respectively, $`A_{ul}`$ and B<sub>ul</sub> the Einstein coefficients. Due to photon loss, the field at frequencies close to the He ii Ly$`\alpha `$ is modified in the following way:
$$\overline{J}(r)(1f)\overline{J}(r).$$
(3)
Substituting (1) into (3) we obtain:
$$\overline{J}(r)=(1\beta f+f\beta )S(r)$$
(4)
Comparing (5) to (1) we see that the classical Sobolev escape probability is modified in the following way:
$$\beta (\beta +ff\beta ).$$
## 5 Stellar analysis
In this section we present the stellar parameters for WR11 derived with the two codes cmfgen and isa-wind (supplemented by the Monte Carlo line-blanketing code). We will discuss the quality of the fits and the overall flux distribution of the model atmospheres.
### 5.1 Analysis technique
In our approach, diagnostic optical lines of He i ($`\lambda `$$`\lambda `$5876, 10830) and He ii ($`\lambda `$$`\lambda `$4686,5412) are chosen to derive the stellar temperature and mass-loss rate. Initially we assume the carbon and oxygen abundances to be the same as those derived for the single WC8 star WR135 by Dessart et al. (2000), because of the similarity of the two spectra (a comparison of the rectified spectra is shown in Fig.5). Once a fair agreement is reached between the modelled and the observed helium lines, carbon and oxygen abundances can be altered to improve the fits to their diagnostic lines. For carbon, lines of C iii at 6741 Å and of C iv at 5471 Å are used. For oxygen, the lack of good diagnostic lines in our spectral range and the weak and blended nature of those that can be identified, make the determination of the oxygen abundance impossible. We therefore adopt the theoretical O/C number ratio of 0.2 (Meynet et al. 1994). Changing oxygen and carbon abundances acts on the total opacity of the wind and therefore feeds back on the modelled strengths of the helium lines. Parameters have therefore to be adjusted so as to retrieve the fits of the helium lines.
We adopt the wind terminal velocity determined by St.Louis et al. (1993), $`v_{\mathrm{}}=1550`$ km s<sup>-1</sup>, from the variability of the UV spectrum. The velocity law used by the cmfgen code (Hillier & Miller 1999) is a two component “beta” law which produces a slower flow than the usual $`\beta `$ = 1 velocity law, at points in the atmosphere intermediate between the outer regions where $`vv_{\mathrm{}}`$ and the photosphere. Hillier & Miller (1999) found that it leads to better fits to the line profiles, while at the same time Schmutz (1997) determined a similar velocity law from his hydrodynamic calculation (which, however was carried out only for the outer layers of the wind). The two component velocity law used for our calculations uses $`v_0`$ = 100 km s<sup>-1</sup> (the photospheric velocity), $`v_{\mathrm{core}}`$ = 1.0 km s<sup>-1</sup> (the velocity at $`R_{}`$), $`v_{\mathrm{ext}}`$ = 1100 km s<sup>-1</sup> (the intermediate velocity) and $`v_{\mathrm{}}`$ = 1550 km s<sup>-1</sup> (the terminal velocity). The two beta exponents used are $`1.0`$ and $`50`$, for $`\beta _1`$ and $`\beta _2`$, respectively. The isa-wind calculation was performed with the same velocity law as cmfgen. A clumping filling factor of 0.1 was adopted for cmfgen. Usually this is derived together with the other parameters, by fitting line wings; for WC stars the heavy blending makes this task harder. We therefore adopted the value obtained by Hillier & Miller (1999) and confirmed by Dessart et al. (2000).
The atomic data models used for the two codes are summarised in Table 2 (see Dessart et al. (2000) for a description of the atomic data used in cmfgen, while isa-wind uses the latest Opacity Project data). Note that the number of individual and super levels ($`N_F`$ and $`N_S`$) mean different things for the two codes: cmfgen bundles $`N_F`$ levels into $`N_S`$ super-levels, while all isa-wind atomic levels, listed under the $`N_F`$ heading, are treated individually (apart for He i and C iv for which the upper 32 and 6 levels, respectively, are super-levels). The data used by cmfgen is more extensive than those used by isa-wind. Generally this can have an effect on the ionization equilibrium (and hence on the spectrum), or only have an effect on particular lines between high lying levels. Tests using cmfgen with reduced atomic data show that while individual lines can be very sensitive to the amount of levels used (e.g. C iii $`\lambda 8500`$) the overall ionization of the wind is not critically dependent on the size of the atomic model provided a minimum amount of levels is used.
### 5.2 Optical and Infra-red line fits
The overall fit quality of the helium diagnostic lines is fair and comparable for the two codes (Fig. 4). Due to the difficulties in rectifying WC spectra and the weakness of the lines in the region 4400–5450 Å region, it was decided that the He i $`\lambda `$$`\lambda `$5876,10830 would be better diagnostic lines than He i $`\lambda `$4471, although the fit to this line is not inconsistent. $`\lambda `$5876 is very well reproduced by cmfgen, while $`\lambda `$10830 is under-predicted by about 15% (one should take into account that this line was not corrected for the O star absorption line, which can be observed as a depression on its blue side). isa-wind only reproduces $`\lambda `$5876 (albeit not perfectly), due to the sensitivity of $`\lambda `$10830 to the outer wind electron temperature, which is calculated in the grey LTE approximation and fixed for outer radii to 10 000 K.
He ii $`\lambda `$5412 is well reproduced by both codes, although cmfgen slightly over-predicts it. Once again we use the 4440-4600 Å region only as a secondary diagnostic: cmfgen does reproduce the He ii $`\lambda `$4541-C iv $`\lambda `$4553 rather well, while isa-wind does not predict the C iv component in the blend with an overall poorer result. $`\lambda `$4686 is slightly under-predicted by cmfgen, with the isa-wind synthetic line fit being affected by blending with the over-predicted C iii-C iv line at 4660 Å.
The He ii lines at $`\lambda `$4859 and $`\lambda `$6560 Å are underestimated by both cmfgenand isa-wind, by about 20% and 80% respectively. In the case of $`\lambda `$6560 the presence of an under-predicted flat-topped C ii line could be responsible (see below). It is in fact not uncommon that models under-predict line emission from low ionization stages (cf. C ii$`\lambda `$4267), showing that the stratification of the wind is imperfectly reproduced. It is possible that this problem will be solved once more realistic clumping is included in the models. On the other hand there is no blend in $`\lambda `$4859 that could be blamed for the poor fit. For that line there could be a slight rectification problem, although this alone could not account for the amount by which the model under-fits the data.
In Fig. 5 we show a comparison of the observed spectra of WR11 and WR135. The two spectra are very similar with WR135 having only slightly stronger He ii $`\lambda `$$`\lambda `$4541, 4686, 5412 lines and comparable He i $`\lambda `$$`\lambda `$5876,7066 lines. On the other hand, WR11 has stronger He ii $`\lambda `$$`\lambda `$4339, 4859, 6560, indicating an opposite trend to the rest of the He ii spectrum! (He ii $`\lambda `$4100 is also shown in Fig.5 although due to the heavy blending with carbon lines, it is difficult to make a clear comparison). This points to an anomaly in the strength of these He ii Pickering lines. We should also note that the ratio of He ii5411 to H$`\beta `$ is different for the two stars, indicating something anomalous in the line formation in one or both stars.
For carbon lines, cmfgen reproduces strengths and shapes to a higher degree of accuracy. isa-wind synthetic lines underestimate C iv $`\lambda `$$`\lambda `$4440,4780 by about 20% and 40%, respectively, while C iii $`\lambda `$$`\lambda `$6741,8500 are weaker than the observations by 80% and 50%, respectively. On the other hand while cmfgen reproduces $`\lambda `$$`\lambda `$4440,4780,6741, the C iii line at 8500 Å is over-estimated by more than 100%. The $`\lambda `$4660 C iii \- C iv blend is overestimated by isa-wind ($``$60%), while it is underestimated by cmfgen ($``$10%). Some lines of C iii (e.g. $`\lambda `$$`\lambda `$8196,8500) and C iv ($`\lambda `$7063), are under-predicted or not predicted at all by isa-wind partly because they are not accounted for in the atomic model. The carbon abundance determination, however, rests critically on the ratio of He ii $`\lambda `$5412 / C iv $`\lambda `$5471 which is well matched by both codes. C ii is under-predicted by cmfgen and not predicted at all by isa-wind (which includes only 7 levels in the atomic model of C ii). However the presence of the C ii line at 4267 Å indicates that C ii lines are present in the spectrum. The spectral classification lines C iii $`\lambda `$5696 and C iv $`\lambda `$5806 are not included in the observations, so in Fig. 4 we only show the prediction for both models. These two lines are remarkably different in the two models: the spectral types deduced from their equivalent width ratio are WC6 and WC9 for the isa-wind and cmfgen codes, respectively, where the criteria of Crowther et al. (1998) have been adopted. The model for WR135 of Dessart et al. (2000), with similar parameters, has a much larger C iv $`\lambda `$5806/C iii $`\lambda `$5696 line ratio, consistent with a WC8 classification. Although these two lines form an excellent spectral type diagnostic because of their mutual proximity and their strengths, it is important to note that the correspondence of their strengths to the stellar parameters is not straight forward.
Differences in the carbon spectrum are obtained (De Marco et al. 1999b) not only between cmfgen and isa-wind, but also with the co-moving frame code of Koesterke & Hamann (1995). Such discrepancies are not due to differences in the atomic data, and it is extremely difficult to identify the exact cause, since the effect is highly dependent on the parameter space investigated. Further investigations are under-way.
A comparison of the cmfgen synthetic spectrum with the rectified ISO spectrum of $`\gamma `$ Velorum is shown in Fig. 6. Our WR+O model spectrum fits the emission lines to a good level of accuracy, demonstrating that line fits are consistent also in the IR. Spectral features in the H and K windows agree qualitatively with unpublished AAT spectroscopy of L.F. Smith (priv. com.). As mentioned in Sect.3, the combined O+WR model continuum is $``$10% lower than the flux-calibrated ISO spectrum for the adopted light ratio, reddening and observed $`v`$ magnitude.
In conclusion it was felt that given the limitations of the observed spectra, the parameters of WR11 from the optical fits have been constrained rather well. The uncertainties listed in Table 3 are determined from the fits (for the $`T_{}`$, $`\dot{M}`$, C/O and $`v_{\mathrm{}}`$ parameters) or from error propagation. In the case of the spectroscopic luminosity, $`L`$, the uncertainty derives directly from the uncertainty on $`M_V`$ (see Sect. 3) since the uncertainty on the bolometric correction BC is negligible. For the radius, $`R_{}`$, the error propagates from the luminosity and temperature uncertainties. Luminosity, mass-loss and terminal velocity determine the uncertainty on the wind efficiency $`\eta `$. Finally the uncertainty on the ionizing fluxes, $`Q_0`$ and $`Q_1`$ derive directly from $`L`$, without taking into account internal model uncertainties.
Given the different nature of the two non-LTE codes, it is remarkable how similar the deduced parameters are. This confirms that despite the approximations adopted by isa-wind, this faster code is suitable for the analysis of Wolf-Rayet stars in the WC8 parameter domain (as it was also shown to be appropriate for a WN8 star by Crowther et al. 1999).
### 5.3 UV ionizing fluxes and bolometric luminosity
In Fig. 7 we present the predicted extreme UV fluxes from cmfgen (solid line, top panel) and isa-wind (solid line, bottom panel) compared with the un-blanketed isa-wind model flux (dotted lines) on the energy vs log ($`\lambda `$F<sub>λ</sub>) plane.
As can be seen, the isa-wind flux at $`\lambda 400`$ Å is harder than the cmfgen flux, a discrepancy that was already pointed out by Crowther et al. (1999). On the other hand the Rydberg vs. $`\lambda F_\lambda `$ plot is designed to emphasise the 0-912 Å region (1 Rydberg = 13.6 eV) and we see by looking at the values of log(Q<sub>0</sub>/s<sup>-1</sup>) and log(Q<sub>1</sub>/s<sup>-1</sup>) shown in Table 3, that in fact the two fluxes are not particularly dissimilar short of 912 Å and 504 Å. On the other hand, the fluxes shortward of $``$400 Å, which influence the strengths of nebular lines such as O<sup>2+</sup>, Ar<sup>3+</sup> and Ne<sup>2+</sup>, are very different and this discrepancy could affect derived nebular properties.
The spectroscopic luminosity determined by cmfgen and isa-wind are not dissimilar, although the former is slightly higher (Table 3). Both the luminosities are lower than that determined in Paper I from alternative methods (log(L/L) = 5.2$`\pm `$0.1), even considering the formal uncertainty determination.
The luminosity of 63000 L obtained by Schaerer et al. (1997) can be compared to that obtained here by the cmfgen code of 100 000 L. We can accredit the increase to the presence of line blanketing (which includes photon loss from the He ii Ly$`\alpha `$ \- see next section), although the difference of the two codes and the neglect of carbon and oxygen might also play a role.
## 6 The effect of line-line interaction
Line-line interaction between the very strong He ii Ly$`\alpha `$ and nearby weaker lines has a large effect on the wind ionization structure. Approximations in the line-blanketing technique that affect either the position or the strength of even the weakest lines can affect greatly the derived stellar parameters. In this section we determine the size of this effect and show its consequences for the derived parameters of $`\gamma `$ Velorum.
Schmutz (1997) proposed that capturing of He ii Ly$`\alpha `$ line-photons by resonant metal transitions, would reduce the ionization balance of the wind. Including line-line interaction between He ii Ly$`\alpha `$ and lines from ions such as Ca v, Fe vi, Ni vi and O iii in his non-LTE model of the WN4 star WR6, led to a significantly lower wind ionization structure than that obtained without taking into account such interaction. As a result, he had to adopt a higher effective temperature to fit the spectrum of this star, leading to a larger luminosity and a lower wind performance number. In this way enough radiation force was obtained to drive the outer parts of the wind of WR6.
Schmutz (1997), however, carried out only an approximate calculation, which kept the fraction of photons captured by metal lines constant throughout the wind, instead of varying with radius. Additionally, the populations of the levels responsible for the metal transitions intercepting the He ii photons, were calculated in LTE, instead of being calculated consistently with the populations of all other levels via the statistical equilibrium equations. Since the extent of this effect could lead to the long-sought answer to a driving mechanism for WR stars, we thought it important to investigate it further.
The line-blanketed cmfgen code (Hillier & Miller 1999) automatically takes into account line-line interaction between all the lines included in the adopted atomic model so that photon loss should naturally take place in our calculation of Sect. 5, between He ii Ly$`\alpha `$ and the included lines of iron, oxygen, calcium and carbon. To determine the effect of line-line interaction on the wind ionization, we removed from the atomic model sets of lines in the spectral region neighboring ($`\pm 2000`$ km s<sup>-1</sup>) He ii Ly$`\alpha `$. Throughout the rest of this discussion, we will refer to the model calculated with this reduced atomic model the “no-PL model” (which stands for no photon loss), while the model calculated in Sect. 5, which fits the spectrum of $`\gamma `$ Velorum, is referred to as the “basic model”.
The elimination of 21 Fe vi lines, results in a small increase in the wind ionization balance with respect to the model of Sect.5. Next we eliminated the O iii Bowen lines at 305.72 and 303.65 Å ($`{}_{}{}^{3}P^3P^o`$ and $`{}_{}{}^{3}P^3D^o`$) and a further increase in the ionization balance was observed. Next we eliminated 1 Ca v line, with no change in the calculated spectrum and finally we removed 63 C iii lines, which produced a further small increase in the ionization balance. In Fig. 8 we compare the synthetic spectrum from the basic model (solid line) with the spectrum from the no-PL model (dashed line), while in Fig. 9 we present the corresponding ionization structures. The basic model presents a less ionized spectrum, with weaker He ii lines and stronger He i lines. The ionization structure of helium is clearly higher in the model without the metal lines, although that of carbon and oxygen appear similar, even if the lines of C iv follow the trend of the He ii lines. Eliminating some metal lines does indeed change the wind ionization balance. The ionization structure and synthetic spectrum of the basic model can be matched of we increase the effective temperature of the no-PL model by about 10 000 K.
As a further test of the effect of line-line interaction on the overall wind ionization we calculated another model, with the same atomic data as in Sect.5, but with the turbulent velocity of the lines in the He ii Ly$`\alpha `$ ($`\pm `$300 km s<sup>-1</sup>) region set at 5 km s<sup>-1</sup><sup>2</sup><sup>2</sup>2This yields a Doppler velocity of 13 km s<sup>-1</sup>. instead of 50 km s<sup>-1</sup>. In this way line-line interaction is reduced, since a smaller turbulent velocity leads to narrower lines, smaller overlap and therefore less interaction. We refer to this model as the “low-vturb model”
In Fig. 8 we show the synthetic spectrum for the low-vturb model (dotted) compared to the basic model (solid), and the no-PL model (dashed), while in Fig. 9 we compare their respective ionization structures. The ionization equilibrium of the low-vturb model is higher than for the basic model but similar to the no-PL model. By reducing the line widths we have reduced the interaction between lines with an overall shift to a higher wind ionization, similar to the shift produced by the elimination of metal lines near He ii Ly$`\alpha `$.
The interception of continuum photons by metal lines (continuum blanketing) is an important effect and the shift in ionization observed when eliminating metal lines (solid and dashed lines in Figs. 8 and 9) could also be due to the lack of continuum blanketing by the lines we eliminated. The change in ionization between the basic model and the low-vturb model (solid and dotted lines in Figs. 8 and 9), on the other hand, is due to a reduced line-line interaction. The similarity of the ionization shifts in the two trials (dashed and dotted lines in Figs. 8 and 9), indicates that line-line interaction plays an important role in the change observed when eliminating metal lines. Weak but numerous metal lines, with wavelengths near that of strong emission lines are important in the overall ionization balance of the wind and therefore in the derived parameters.
From earlier trials we can estimate that the difference between the basic model (solid) and the no-PL or the low-vturb models (dashed and dotted, respectively) in Fig. 8 is equivalent to about 10 kK in effective temperature, but only about 20% in bolometric luminosity (because the radius to which the stellar temperature refers also changes). We can therefore conclude that for this parameter space photon-loss from the He ii Ly$`\alpha `$, does lead to a higher luminosity, although not as high as that determined via the evolutionary mass-luminosity relation.
## 7 Discussion
In the following sections we present a critical summary of the work carried out in this paper.
### 7.1 Stellar parameters and bolometric luminosity
We have presented a detailed spectroscopic analysis of the WR star in the $`\gamma `$ Velorum WR+O binary system. Using the wavelength-dependent light ratio and the synthetic O star spectrum from Paper I we have generated a corrected WR optical and IR spectrum which we have fitted with a synthetic spectrum using the cmfgen non-LTE model atmosphere code. Results are summarised in Table 4. Our derived mass-loss is in excellent agreement with the mass-loss derived by Stevens et al. (1996) from X-ray asca observations, once their value is scaled to the hipparcos distance. Although the parameters of temperature, mass-loss and carbon abundance are well constrained, the oxygen abundance is assumed. The discrepancy between the He ii lines at 4340, 4859 and 6560 Å and other He ii lines remains to be explained.
Previous analyses of WR11 were carried out by Schaerer et al. (1997) and by De Marco et al. (1999a - using the stellar atmosphere code of Koesterke & Hamann, 1995). Schaerer et al. determined the stellar effective temperature by fitting the He i $`\lambda `$4471 / He ii $`\lambda `$4541 ratio, assuming a (higher) mass-loss and carbon abundance. De Marco et al. (1999a) derived larger mass-loss and temperature, and a lower wind terminal velocity. This can be attributed to the fact that in the parameter space of $`\gamma `$ Velorum the helium spectrum can be reproduced approximately with a range of temperature and mass-losses. On the other hand, when the lines are analysed in detail, it is clear that for temperatures larger than about 57 kK and log($`\dot{M}`$/M yr<sup>-1</sup>)$`>`$–4.45, predicted line shapes deteriorate, becoming wider. This effect was partly compensated for by De Marco et al. (1999a) by adopting a lower wind velocity, although it is clear from their fits that most lines are still too wide. Further, their neglect of line-blanketing and therefore of photon loss (no line-line interaction between the key lines), might have contributed to the different parameters.
With respect to the analysis of Schaerer et al. (1997), the star becomes twice as luminous, 10% hotter and shows a mass-loss lower by a factor of two. The carbon abundance is lower by 40% (although we remind the reader that Schaerer et al. assumed both the mass-loss and the carbon abundance). This carries implications for studies of wind-wind interaction. From this study and the results presented in Paper I for the O star, we determine:
$$\frac{(\dot{M}v_{\mathrm{}})_{WR}}{(\dot{M}v_{\mathrm{}})_O}33$$
compared to a value of 208 derived by Schaerer et al. (1997). Adopting the relationships of Eichler & Usov (1993) we determine that the region of stellar wind collision is approximately 9-30 R from the O star (0.15 of the distance that separates the two stars (63-200 R \- Schmutz et al. 1997) vs. 0.06 for the study of Schaerer et al. (1997)). Overall the stellar parameters derived are not dissimilar from those derived for the single WC8 star WR135 by Dessart et al. (2000; see Table 4). This indicates that it is unlikely that substantial levels of emission are produced by the collision region or that an anomalous ionization is induced by the O star ionizing flux.
The derived spectroscopic luminosity ($`\mathrm{log}L/L_{}=5.0\pm 0.1`$) is increased with respect to that found in previous studies, partly due to line-line interaction, and is close to the luminosity derived in Paper I via the mass-luminosity relationship for WR stars ($`\mathrm{log}L/L_{}=5.2\pm 0.1`$). An important implication of the higher luminosity and of a 10% clumping factor, is that it helps to reduce the performance number by a factor of $``$20, from 144 to 7, assuming a 10% clumping filling factor. Lucy & Abbott (1993) showed that radiation pressure on spectral lines should in principle be able to drive winds with performance numbers $``$ 10. So, the strongly reduced $`\eta `$ suggests that line driving is the mechanism responsible for the dense wind of Wolf-Rayet component in $`\gamma `$ Velorum.
### 7.2 cmfgen vs. isa-wind
Parameters derived with the fast Sobolev approximation code isa-wind are similar to those derived with cmfgen, apart from the derived carbon mass fraction, with cmfgen resulting in a carbon abundance which is a factor of two larger than isa-wind. This discrepancy is not totally understood although we should note that the co-moving frame code of Koesterke & Hamann (1995) compares favourably with cmfgen for cases in common (e.g. for WR135 - Dessart et al. 2000).
A secondary difference is in the fluxes at $`\lambda `$$`<`$400 Å. Crowther et al. (1999) tested the far UV model fluxes of cmfgen and isa-wind by modelling the H ii region associated with the WN8 star WR124, showing that, for that system, the isa-wind flux was too hard. Also in the current analysis we find that the isa-wind flux is harder for $`\lambda `$$`<`$400 Å: we suspect that this discrepancy is due to the lack of photon branching in the Monte Carlo code, which would lead to distributing energetic photons to longer wavelengths. Although the ionising properties of the two synthetic atmospheres (the values of $`Q_0`$ and $`Q_1`$) are similar, the difference in the extreme UV can carry implications when model WR fluxes are used to interpret nebular line emission from extra-galactic stellar populations (Leitherer et al. 1999).
We stress that the comparison between the two codes is reasonable only if photon loss is included in isa-wind via the approximation detailed in Sect.4.3. If photon loss were not accounted for, the determined isa-wind effective temperature would be about 8000 K lower, comparable to the difference found in Sect.6 when reducing the interaction between metal lines and He iiLy$`\alpha `$.
### 7.3 The photon loss effect
The photon loss effect demonstrated by Schmutz (1997) as the cause of a lower wind ionization, was tested here using the cmfgen code which treats line-line interaction as part of the blanketing. For this parameter space the interaction between the He ii Ly$`\alpha `$ and nearby O iii and C iii lines causes a shift to lower ionization equilibrium, equivalent to $``$10 000 K or a $``$20% increase in the luminosity<sup>3</sup><sup>3</sup>3The quantification of the luminosity shift resulting from a different temperature depends on the location of the $`R_{}`$, since this can change from model to model, it is not possible to simply derive the change in luminosity from the change in temperature and the Boltzman relationship.. By reducing the turbulent velocity from 50 to 5 km s<sup>-1</sup> for lines residing near 303 Å, we have shown that the wind ionization is enhanced in a similar way, demonstrating that line-line interaction, and not continuum blanketing, is at the origin of the ionization shift observed when eliminating metal lines near the He ii Ly$`\alpha `$. If line-line interaction between strong resonance lines and weaker but more numerous metal lines leads to a shift in the wind ionization, a detailed treatment of even the weakest lines is important. We remind the reader that our test remains incomplete because, although cmfgen treats line-line interaction not all lines are accounted for. Amongst the ones that are not treated there might be some that could play an important role in this mechanism. Additionally, our conclusion is only valid within the parameter space appropriate for WR11, a WC8 star, since for different temperature and mass-loss the strength of the He ii and other strong lines will be different (cf. Crowther et al. (1999) who found that photon-loss did not play a role in the atmosphere of their WN8 star).
One should conclude from the presence of the photon loss mechanism that a proper treatment of the metal lines not only requires one to include very many lines (to realistically model the blanketing), but also the very critical ones, such as those near the strong He ii Ly$`\alpha `$ and possibly at other resonance lines such as C iv at 1550 Å. Furthermore, the presence of photon loss shows that applying mean blanketing factors (Pauldrach et al. 1996) - averaged over 20–50 Å intervals - may miss out some interaction with an overall effect on the derived stellar parameters for WR stars. Finally, it is clearly pointed out that macro-turbulence is connected with photon loss, as the amount of turbulence essentially dictates the size and strength of the set of lines involved in the line-line interactions.
###### Acknowledgements.
OD acknowledges support from PPARC grant PPA/G/S/1997/00780. PAC is funded by a Royal Society University Research Fellowship. DJH acknowledges support from NASA grants NAG5-8211 and NAGW-3828. AdK acknowledges support from NWO Pionier grant 600-78-333 to L.B.F.M. Waters and from NWO Spinoza grant 08-0 to E.P.J. van den Heuvel. JS acknowledges the grant Wo 296/22-1,3 by the Deutsche Forschungsgemeinschaft. Thanks to Patrick Morris and Karel van der Hucht for providing the ISO data, which is an ESA project with instruments funded by ESA member states (especially the PI countries: France, Germany, The Netherlands and the United Kingdom) with the partecipation of ISAS and NASA.
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# The von Neumann algebra of the non-residually finite Baumslag group ⟨𝑎,𝑏∣𝑎𝑏³𝑎⁻¹=𝑏²⟩ embeds into 𝑅^𝜔
## 1. Introduction
In this paper we analyze the structure of sets of (non-commutative) moments $`\tau \left(x_1\mathrm{}x_n\right)`$ of variables $`x_1,x_2,\mathrm{},x_n`$ in a type $`\mathrm{II}_1`$ factor $`M`$. We analyze the structure of these sets, for the case of projections and unitaries (see also \[Ra\] for odd moments of selfadjoint elements). While the understanding of these structures is far from being complete, we prove that any discrete (i.c.c.) group $`\mathrm{\Gamma }`$ that can be faithfully embedded into the unitary group of $`R^\omega `$ has the property that $`(\mathrm{\Gamma })R^\omega `$.
By using techniques pertaining to free probability we prove that the Baumslag group $`\mathrm{\Gamma }=a,bab^3a^1=b^2`$, which is non-residually finite, embeds (faithfully) into $`𝒰(R^\omega )`$. Note also that by \[CeGr\], the algebra $`(\mathrm{\Gamma })`$ does not have property $`\mathrm{\Gamma }`$. In particular, it is non-hyperfinite. (We are indebted to P. de la Harpe for bringing this to our attention.)
This is positive evidence towards Connes’s conjecture that any separable $`\mathrm{II}_1`$ factor is embedded into $`R^\omega `$.
In the second part of the paper we analyze the structure of the sets of first and second order of moments $`\tau \left(e_i\right)`$, $`\tau \left(e_ie_j\right)`$ of finite sets of projections $`e_1,\mathrm{},e_n`$ in a type $`\mathrm{II}_1`$ factor. By work of Kirchberg \[Ki\], if the closure of the sets of first and second order of moments of unitaries is independent of the type $`\mathrm{II}_1`$ factor considered, then Connes’s conjecture should be true.
It is obvious how to translate this statement in terms of projections. By using methods from \[Ra\], it follows that the corresponding set of moments, i.e., the set of moments (of order $`1`$ and $`2`$) of projections, is convex and multiplicative for a $`\mathrm{II}_1`$ factor $`M`$ such that $`MMM`$ and $`(M)=_+\backslash \{0\}`$. We will also analyze the structure of faces of these sets, which gives some additional data on the geometric structure of these sets.
This work was supported by NSF grant DMS99-70486.
*Definitions and Notations*: We recall that for a von Neumann algebra $`M`$, the unitary group is denoted by $`𝒰(M)`$, while $`𝒫(M)`$ stands for the set of selfadjoint projections.
For a type $`\mathrm{II}_1`$ factor $`M`$, the fundamental group $`(M)`$ of $`M`$ \[MvN\] \[vN\] is defined as the multiplicative group of all $`t>0`$, such that $`M_tM`$.
For $`\mathrm{\Gamma }`$ a countable discrete group, with infinite conjugacy classes (briefly i.c.c.), the algebra $`(\mathrm{\Gamma })`$ is the weak closure of the group algebra $`(\mathrm{\Gamma })`$ embedded (via left regular representation) in $`B(\mathrm{}^2(\mathrm{\Gamma }))`$.
If $`\omega `$ is an ultrafilter on $``$, then following \[McD\] and \[Co\], one defines for any $`\mathrm{II}_1`$ factor $`M`$ the ultrafilter product $`M^\omega `$, obtained via G.N.S. construction, by defining the trace of an element $`(x_n)`$ in the infinite product of copies of $`M`$ to be $`lim_{n\omega }`$ $`\tau (x_n)`$ (in the hypothesis that $`supx_n<\mathrm{}`$). We refer to Connes’s \[Co\] paper on injectivity for full details on this construction.
Finally, we recall a construction from \[Po\] (see also \[Vo\]). Consider two von Neumann algebras $`N_1`$, $`N_2`$ which have a common subalgebra $`B`$, containing the unit. Also assume that the algebras $`N_i`$ have faithful traces whose restriction to $`B`$ coincides. We assume that we are given conditional expectations $`E_i`$ from $`N_i`$ onto $`B`$ that are trace preserving.
The trace on the reduced amalgamated free product von Neumann algebra $`C_1_BC_2`$ is defined by the requirement that a product $`c_{1,1}c_{2,1}c_{1,2}c_{2,2}c_{1,3}c_{2,3}\mathrm{}`$, where $`c_{1,i}`$ belongs to $`C_1`$ and $`c_{2,i}`$ belongs to $`C_2`$, has zero trace if $`IdE_B(c_{ij})0`$ for all $`i=1,2`$, $`j=1,2,\mathrm{}`$.
## 2. Moments of unitaries
In this section we define some sets of non-commutative moments of unitaries $`\tau \left(u_1u_2\mathrm{}u_p\right)`$ associated with a type $`\mathrm{II}_1`$ factor $`M`$. We will use these sets to check that for any discrete i.c.c. group $`\mathrm{\Gamma }`$ that can be embedded into the unitary group of $`R^\omega `$, then also $`(\mathrm{\Gamma })`$ can be embedded into $`R^\omega `$ (as a unital subfactor).
First we consider the set of all possible embeddings (up to order $`N`$) of a group-like algebra.
By $`𝒱_{n,p}`$ we denote the set of all indices $`(i_1,i_2,\mathrm{},i_k)`$ with $`1kp`$ and $`i_1,i_2,\mathrm{},i_k`$ in $`\{1,2,\mathrm{},n\}`$. By $`u_I`$ we denote the product $`u_{i_1}u_{i_2}\mathrm{}u_{i_k}`$ if $`I=(i_1,i_2,\mathrm{},i_k)`$.
###### Definition 2.1.
Let $`M`$ be a separable type $`\mathrm{II}_1`$ factor and consider the following subset $`(`$of $`\{0,1\}^{2^{n^p}})`$, denoted by $`K_M^{n,p}`$. We define $`K_M^{n,p}`$ by requiring that $`(\epsilon _I)_{\left|I\right|p}`$ belongs to $`K_M^{n,p}`$ if and only if there exist unitaries $`u_1,u_2,\mathrm{},u_n`$ in $`𝒰(M)`$ such that $`\tau (u_I)=1`$ or $`0`$, and $`\epsilon _I=\tau (u_I)`$, for $`\left|I\right|p`$. Here $`I`$ is an index set $`(i_1,i_2,\mathrm{},i_k)`$, with $`ij\{1,2,\mathrm{},n,\}`$, $`1kp`$ and $`\left|I\right|=k`$.
###### Remark 2.2.
If $`M=\left(\mathrm{\Gamma }\right)`$ and $`\mathrm{\Gamma }`$ has non-solvable world problem, then there are $`2^{n^p}`$-uples of $`0`$’s and $`1`$’s about which it might be undecidable whether they belong to $`K_M^{n,p}`$.
The above definition might be restrictive for some purposes because it requires that $`\tau (u_I)`$ is either $`0`$ or $`1`$. In fact $`\tau (u_I)=1`$ is equivalent to the fact that $`u_I=1`$. (Here by $`u_I`$ we mean the product $`u_{i_1}u_{i_2}\mathrm{}u_{i_k}`$ if $`I=(i_1,i_2,\mathrm{},i_k)`$.)
###### Definition 2.3.
Let $`M`$ be a finite von Neumann algebra and let for fixed integers $`n,p`$,
$$\begin{array}{c}L_M^{n,p}=\{(\tau (u_I)=\tau (u_{i_1}u_{i_2}\mathrm{}u_k))_{I𝒱_{n,p}}\text{for all}\hfill \\ \hfill \text{ unitaries }u_1,u_2,\mathrm{},u_n\text{ in }𝒰(M)\}.\end{array}$$
The following properties are easy to observe (see also \[Ra\]).
###### Proposition 2.4.
1. Let $`M_1`$, $`M_2`$ be finite von Neumann algebras. Denote by $``$, the pointwise product on $`^s`$, for all $`s`$. Then for all positive integers $`n,p`$,
$$L_{M_1}^{n,p}L_{M_2}^{n,p}L_{M_1M_2}^{n,p},K_{M_1}^{n,p}K_{M_2}^{n,p}K_{M_1M_2}^{n,p}.$$
2. In particular if $`M`$ is such that $`MMM`$, then $`K_M^{n,p}`$ and $`L_M^{n,p}`$ are closed under pointwise multiplication.
3. If $`\lambda (M)`$, then $`\lambda L_M^{n,p}+(1\lambda )L_M^{n,p}L_M^{n,p}`$ for all integers $`n,p1`$. In particular if $`(M)=_+\backslash \{0\}`$, then $`L_M^{n,p}`$ is convex.
4. $`\overline{L_M^{n,p}}L_{M^\omega }^{n,p}`$ and $`L_{M^\omega }^{n,p}`$ is closed in the product topology of $`^{\left|𝒱_{n,p}\right|}`$.
5. In particular if $`MMM`$ and $`(\lambda _I)_{I𝒱_{n,p}}`$ is an element in $`L_M^{n,p}`$ such that either $`\left|\lambda _I\right|<1`$ or $`\lambda _I=1`$, then by replacing the components in $`L_M^{n,p}`$ which are not $`1`$, by zero, we obtain an element in $`K_{M^\omega }^{n,p}`$.
6. Let $`\mathrm{\Phi }_{u_1}`$ be the operation on $`L_M^{n,p}`$ which replaces in $`(\lambda _I)_{IL_M^{n,p}}`$ any monomial $`\lambda _I=\tau (u_I)`$, corresponding to a nonzero total power of $`u_1`$, by zero. Assume $`(M)=_+\backslash \{0\}`$. Then
$$\mathrm{\Phi }_{u_1}\left(L_M^{n,p}\right)L_M^{n,p}.$$
7. For all separable $`\mathrm{II}_1`$ factors $`M`$ we have that $`L_M^{n,p}L_R^{n,p}`$. Moreover if $`MR^\omega `$, then $`\overline{L_M^{n,p}}=\overline{L_R^{n,p}}`$ (and similarly for $`K`$).
Open Question: Does $`K_{(\mathrm{SL}_3())}^{n,p}K_{(F_{\mathrm{}})}^{n,p}`$ for all $`n,p`$?
The proof of properties $`(\mathrm{a})`$$`(\mathrm{c})`$ is obvious and is basically contained in \[Ra\]. To check property $`(\mathrm{d})`$ we have only to verify that $`L_{M^\omega }^{n,p}`$ is closed. But if $`(\lambda _I)_{I𝒱_{n,p}}`$ is an accumulation point for $`L_M^{n,p}`$, then take unitaries $`(u_i^s)_{i=1}^n`$, for all $`s`$, such that $`lim_s\mathrm{}\tau (u_I^s)=\lambda _I`$, $`I𝒱_{n,p}`$. Then $`u_i=(u_i^s)_s`$ are unitaries in $`R^\omega `$, whose non-commutative moments give $`(\lambda _I)_{I𝒱_{n,p}}`$.
Property $`(\mathrm{e})`$ follows from properties $`(\mathrm{a})`$ and $`(\mathrm{c})`$. Property $`(\mathrm{f})`$ follows by convexity, and integration over $`\theta `$, where gauging $`u_1`$ by $`e^{2\pi i\theta }`$. Property $`(\mathrm{g})`$ is obvious.
###### Proposition 2.5.
Let $`\mathrm{\Gamma }`$ be a discrete i.c.c. group that embeds faithfully into the unitary group of $`R^\omega `$. Then $`(\mathrm{\Gamma })R^\omega `$.
###### Proof.
Fix $`n,p`$ and let $`u_1,u_2,\mathrm{},u_n,\mathrm{}`$ be a system of generators of $`\mathrm{\Gamma }`$. Let $`\epsilon _I`$ be the traces of $`(u_I)_{I𝒱_{n,p}}`$ in the left regular representation of $`(\mathrm{\Gamma })`$. By hypothesis there exist unitaries $`v_1,v_2,\mathrm{},v_n`$ in $`R^\omega `$ such that $`\left|\tau (v_I)\right|<1`$ if $`u_I1`$ in $`\mathrm{\Gamma }`$ and $`v_I=1`$ if $`u_I=1`$ in $`\mathrm{\Gamma }`$. Let $`(\alpha _I)^s=\tau (v_I^s)=\tau (v_I)^s`$. Then $`(\alpha _I)_{I𝒱_{n,p}}^s`$ belongs to $`\overline{L_R^{n,p}}`$ and hence so does the limit
$$\epsilon _I=\underset{s\mathrm{}}{lim}(\alpha _I)^s\text{ for }I𝒱_{n,p}.$$
Thus $`(\epsilon _I)_{I𝒱_{n,p}}`$ $`\overline{K_R^{n,p}}`$ for all $`n,p`$. Hence $`(\mathrm{\Gamma })R^\omega `$. ∎
###### Definition 2.6.
We call an i.c.c. group $`\mathrm{\Gamma }`$ hyperlinear if $`\mathrm{\Gamma }`$ embeds faithfully into $`𝒰(R^\omega )`$. Clearly any residually finite group is hyperlinear. The class of hyperlinear groups is obviously closed under free products.
###### Theorem 2.7.
The class of hyperlinear groups is strictly larger than the class of residually finite groups. More precisely, the Baumslag group $`a,bab^3a^1=b^2`$ \[Ba\] \[Ma\] is hyperlinear and non-residually finite. (Note that by \[CeGr\], $`(\mathrm{\Gamma })`$ does not have property $`\mathrm{\Gamma }`$.)
###### Proof.
We divide the proof into several steps. We construct first an approximate embedding of the relation $`ab^3a^1=b^2`$ into $`M_n()`$. We then take the free amalgamated product (over $`b^2`$) by a unitary that commutes with $`b^2`$ and perturb $`a`$ with this unitary. This gives an approximate embedding of $`\mathrm{\Gamma }`$ into the nonscalar unitaries in some free product algebras. Since these algebras are embeddable into $`R^\omega `$ \[Vo\] \[Wa\], the result will follow, by Proposition 2.5.
*Step* I. Construction of an approximate embedding.
There exist unitaries $`v_n`$, $`b_n`$ (of zero trace) in $`M_{6n}()`$ with the following properties:
Property 1)
$$v_nb_n^3v_nb_n^2_{\mathrm{}}\frac{K_1}{n}\text{for a universal constant }K_1.$$
Property 2)
> Denote by $`B_0^n`$ the abelian algebra generated by $`b_n^2`$ and let $`E_{B_0^n}`$ be the corresponding conditional expectation. Let $`\mathrm{\Phi }^n=IdE_{B_0^n}`$. Then for all $`\alpha \{\pm 1,\pm 2\}`$,
>
> $$\mathrm{\Phi }(v_nb_n^\alpha v_n^{})_2K_2$$
>
> for a universal constant $`K_2`$. Here $`_2`$ is the normalized Hilbert-Schmidt trace on matrices.
Property 3)
$$E_{B_0^n}(b_n^{\pm 1})=0,E_{B_0^n}\left(vb^\alpha \right)=0,E_{B_0^n}\left(b^\alpha v\right)=0,\alpha ,\alpha 0.$$
We describe first the construction of the unitaries $`v_n`$, $`b_n`$.
Let $`e_0,e_1,\mathrm{},e_{6n1}`$ be the diagonal algebra of $`M_{6n1}()`$, and for convenience we think of $`e_k`$ as being identified with $`\chi _{[\frac{k}{6n},\frac{k+1}{6n})}`$, $`k=0,1,\mathrm{},6n1`$.
Let $`f_k`$, for $`k=0,1,2,\mathrm{},n1`$, be the projection $`\chi _{[\frac{3k}{6n},\frac{3k+3}{6n})+\{0,\frac{1}{2}\}}`$ and let
$$g_k=\chi _{[\frac{2k}{6n},\frac{2k+2}{6n})+\{0,\frac{1}{3},\frac{2}{3}\}}.$$
Let $`b_n`$ be the unitary defined by
$$b_n=\underset{k=0}{\overset{6n1}{}}e^{\frac{2\pi ik}{6n}}e_k.$$
Let $`v_n`$ be a unitary such that $`v_n^{}f_k=g_kv_n^{}`$, $`k=0,1,2,\mathrm{},n1`$, and such that
$`(Adv_n^{})(e_{3k+\epsilon })`$ $`=e_{2k+\epsilon },`$
$`(Adv_n^{})(e_{3k+3n+\epsilon })`$ $`=(e_{2k+4n+\epsilon })`$
for all $`k=0,1,2,\mathrm{},n1`$, $`\epsilon =0,1`$. (See Fig. 1.)
For $`\epsilon =2`$, we define $`v_n^{}`$ by the requirement
$`Adv_n^{}(e_{3k+2})`$ $`=e_{2k+2n},`$
$`Adv_n^{}(e_{3k+3n+2})`$ $`=e_{2k+2n+1}.`$
Observe that the definition of $`b_n`$, $`v_n`$ implies, for a universal constant $`K_1`$, that
(1)
$$b_n^2f_ke^{\frac{2\pi ik}{n}}f_k_{\mathrm{}}\frac{K_1}{n},$$
(2)
$$b_n^3g_ke^{\frac{2\pi ik}{n}}g_k_{\mathrm{}}\frac{K_1}{n},k=0,1,\mathrm{},n1.$$
Moreover, $`v_n^{}f_k=g_kv_n^{}`$, and
(3)
$$(v_nb_nv_n^{})e_{3k+\epsilon +\alpha 3n}=e^{2\pi i(\frac{2k+\epsilon }{6n}+\alpha \frac{2}{3})}e_{3k+\epsilon +\alpha 3n}\text{ for }\epsilon =0,1,\alpha =0,1.$$
In the remaining case, we have that
(4)
$$(v_nb_nv_n^{})e_{3k+2+\alpha 3n}=e^{2\pi i(\frac{2k+\alpha }{6n}+\frac{1}{3})}e_{3k+2+\alpha 3n}\text{ for }\alpha =0,1.$$
We now proceed to the proof of the properties 1), 2), 3).
Since $`v_k^{}f_k=g_kv_k^{}`$, $`v_kg_k=f_kv_k`$, $`v_ng_kv_n^{}=f_k`$, and $`f_k,g_k`$ commute with $`b_n`$, it follows that
$`v_nb_n^3v_n^{}b_n^2_{\mathrm{}}`$ $`=\underset{k=0,1\mathrm{},n1}{\mathrm{max}}(v_nb_n^3v_n^{}b_n^2)f_k`$
$`=\underset{k=0,1\mathrm{},n1}{\mathrm{max}}v_n(b_n^3g_ke^{\frac{2\pi ik}{n}}g_k)v_n^{}(b_n^2f_ke^{\frac{2\pi ik}{n}}f_k)_{\mathrm{}}.`$
But this quantity is less than $`\frac{K_1}{n}`$ by (1), (2). This completes the proof of Property 1.
To prove Property 2 we need to describe first $`E_{B_0^n}`$. But it is obvious that
$$E_{B_0^n}\left(\underset{k=0}{\overset{6n1}{}}\lambda _ke_k\right)=\underset{k=0}{\overset{3n1}{}}\frac{1}{2}\left(\lambda _k+\lambda _{k+3n}\right)(e_k+e_{k+3n}).$$
Consequently
$$\mathrm{\Phi }_n\left(\underset{k=0}{\overset{6n1}{}}\lambda _ke_k\right)=\underset{k=0}{\overset{3n1}{}}\frac{\lambda _k\lambda _{k+3n}}{2}e_k+\underset{k=0}{\overset{3n1}{}}\frac{\lambda _{k+3n}\lambda _k}{2}e_{k+3n}.$$
We use the above formula for $`v_nb_nv_n^{}`$ and use (3) and (4). We take $`P_n=_{k=0}^{n1}\chi _{[\frac{3k}{6n},\frac{3k+1}{6n})}=_{k=0}^{n1}e_{3k}`$. Then $`P_n`$ has trace $`\frac{1}{6}`$. Since $`\left(v_nb_nv_n^{}\right)e_{3k}=e^{2\pi i\frac{k}{3n}}e_{3k}`$ and $`\left(v_nb_nv_n^{}\right)e_{3k+3n}=e^{2\pi i\left(\frac{k}{3n}+\frac{2}{3}\right)}e_{3k+3n}`$ the above formula shows that
$$P_n\mathrm{\Phi }_n(vbv^{})=P_n(vbv^{}E_{B_0^n}(vbv^{}))=\underset{k=0}{\overset{n1}{}}e^{2\pi i\frac{k}{3n}}\left(1e^{2\pi i\frac{2}{3}}\right)e_{3k}.$$
Hence
$$P_n\mathrm{\Phi }_n(vbv^{})_2^2>\left|1e^{2\pi i\frac{2}{3}}\right|P_n_2^2=\frac{1}{6}\left|1e^{2\pi i\frac{2}{3}}\right|.$$
The computations for $`vb_n^{\pm 2}v^{}`$, $`vb_n^1v^{}`$ are similar, eventually the factor $`\frac{2}{3}`$ being replaced by $`\frac{4}{3}`$ or $`\frac{2}{3}`$. This completes the proof of Property 2.
It is obvious that $`E_{B_0^n}(vb_n^\alpha )`$, $`E_{B_0^n}(b_n^\alpha v)`$, $`\alpha 0`$, and $`E_{B_0^n}(b_n^{\pm 1})`$ are vanishing.
*Step* II. In this step we consider the amalgamated free product of the algebra $`\{v_n,b_n\}^{\prime \prime }`$ described above and $`()B_0^n`$. The amalgamated free product is considered over $`B_0^n`$ (the von Neumann algebra generated by $`b_n^2`$).
Let $`()`$ have the canonical generator $`a_1`$, a Haar unitary. $`()`$ is endowed with the standard trace. Consider the algebra
$$𝒜_n=(()B_0^n)_{B_0^n}\{v_n,b_n\}^{\prime \prime }$$
with the canonical amalgamated free product trace (see section on definitions, \[Po\], and \[Vo\]).
By \[Ra\] (see also \[Dy\], \[Shly\]), we have that $`𝒜_n`$ is a free group factor. Using the ultrafilter construction (\[Co\], \[McD\]), we construct algebras $`𝒜^\omega `$, $`B_0^\omega `$ and $`B^\omega `$ consisting of bounded sequences of elements in the algebras $`𝒜_n`$, $`B_0^n`$ and $`\left\{b_n\right\}^{\prime \prime }`$ respectively. It is obvious that for $`x=\left(x_n\right)_n`$ in $`𝒜^\omega `$ we have
$$E_{B_0^\omega }(\left(x_n\right)_n)=\left(E_{B_0^n}(x_n)\right)_n.$$
Let $`b`$ be the unitary element $`\left(b_n\right)_nB^\omega `$ and let $`B_0`$ be the (abelian) von Neumann algebra generated by $`b`$. Let $`v`$ be the unitary $`v=\left(v_n\right)_n𝒜^\omega `$. We identify $`\left(\right)𝒜^\omega `$ with constant sequences with elements in $`\left(\right)`$.
Since, by \[Wa\], \[Vo\], any type $`\mathrm{II}_1`$ free group factor embeds into $`R^\omega `$, we obtain that the algebra $`𝒜^\omega `$ is embedded into $`R^\omega `$ and hence
$$𝒜=()B_0_{B_0}\{v,b\}^{\prime \prime }R^\omega .$$
The trace on $`𝒜`$ is the amalgamated free product trace and coincides with the restriction of the trace on $`R^\omega `$. Let $`\mathrm{\Phi }`$ be the identity minus the conditional expectation $`E_{B_0}`$ from $`𝒜`$ (or $`𝒜^\omega `$) onto $`B_0`$. The following properties hold true:
1. $`v,b`$ are Haar unitaries, $`vb^3v^{}=b^2`$.
2. $`\mathrm{\Phi }\left(vb^\alpha v^{}\right)_2\frac{1}{6}`$, $`\alpha \{\pm 1,\pm 2\}`$; $`E_{B_0}(vb^k)=0`$, $`k`$, $`E_{B_0}(b^{\pm 1})=0`$.
To prove Property 2, note that $`E_{B_0^\omega }\left(vb^\alpha v^{}\right)_2\frac{1}{6}`$ because of the corresponding property for $`v_nb_n^\alpha v_n^{}`$. Since $`B_0B_0^\omega `$ it follows also that $`E_{B_0}\left(vb^\alpha v^{}\right)_2\frac{1}{6}`$.
Let $`a_1`$ be the standard generator of $`()`$ and let $`A=a_1v`$, $`B=b`$.
*Step* III. Let $`A,B`$ be the unitaries defined in Step II. Clearly $`AB^3A^1=B^2`$, as $`a_1`$ commutes with $`b^2`$. Let
$$W=A^{\alpha _1}B^{\beta _1}A^{\alpha _2}\mathrm{}A^{\alpha _n}B^{\beta _{n+1}}$$
be a word in $`A,B`$ such that $`\beta _10,\mathrm{},\beta _n0`$, $`\alpha _20,\mathrm{},\alpha _n0`$. Consider the following assumption on the sequence of the indices $`\alpha _i`$.
*Assumption* on $`W=A^{\alpha _1}B^{\beta _1}A^{\alpha _2}\mathrm{}A^{\alpha _n}B^{\beta _{n+1}}`$. One of the following possibilities occurs (about consecutive indices):
1. Either $`\alpha _i`$, $`\alpha _{i+1}`$ are both positive or negative (except for the case when $`\alpha _1=0)`$.
2. If $`\alpha _i<0`$, $`\alpha _{i+1}>0`$, then $`\beta _i\{\pm 1\}`$.
3. If $`\alpha _i>0`$, $`\alpha _{i+1}<0`$, then $`\beta _i\{\pm 1,\pm 2\}`$.
###### Claim.
If the word $`W`$ is subject to the conditions A1, A2, A3, then $`W`$ is not a multiple of a scalar (and hence $`\left|\tau (W)\right|<1`$).
###### Proof of the claim in Step III.
We use the following property of an amalgamated free product $`=E_CF`$ where $`E`$, $`F`$ are finite algebras with faithful traces $`\tau _1`$, $`\tau _2`$ whose restrictions coincide on the common unital subalgebra $`C`$.
Assume $`w=e_1f_1e_2f_2\mathrm{}e_nf_{n+1}`$ is a word in $`E_CF`$, $`e_iE`$, $`f_iF`$, such that $`IdE_C(f_1)`$, $`IdE_C(e_2)0`$, $`\mathrm{}`$, $`IdE_C(e_n)0`$. Then $`w`$ is not a scalar multiple of the identity. This follows for example from the construction in \[Po\].
Then for the word $`W=A^{\alpha _1}B^{\beta _1}A^{\alpha _2}\mathrm{}A_g^{\alpha _{n+1}}`$, we use the fact that $`A=a_1v`$, $`B=b`$. Since
$$E_{B_0}(b^\theta v^{}),\text{ }E_{B_0}(vb^\theta )$$
are always zero for all $`\theta `$, the only instances in the product in $`W`$ where we could have elements with $`IdE_{B_0}`$ nonzero are in subsequences of the form
$$\mathrm{}avb^{\pm \alpha }v^{}a\mathrm{},\alpha \{1,2\}\text{ (in }\mathrm{}AB^{\pm \alpha }A^1\mathrm{}\text{),}$$
or
$$\mathrm{}v^{}a^{}b^{\pm 1}av\mathrm{}\text{ (in }\mathrm{}A^1B^{\pm 1}A\mathrm{}\text{).}$$
But in these cases $`\mathrm{\Phi }=IdE_{B_0}`$ applied to the elements $`vb^{\pm 1}v^{}`$, $`vb^{\pm 2}v^{}`$, and $`b^{\pm 1}`$ is nonzero. The remaining two cases correspond to subsequences involving $`A^nB^\theta A^m`$ with $`\theta 0`$ and $`n`$, $`m`$ both strictly positive or both strictly negative. The case $`n,m>0`$ corresponds to a subsequence of the form $`\mathrm{}a_1vb^\alpha a_1v\mathrm{}`$ or $`\mathrm{}a_1b^\alpha v^{}a_1v\mathrm{}`$. In either case we use the fact that $`E_{B_0}\left(b^\alpha v\right)=0`$, $`E_{B_0}\left(v^{}b^\alpha \right)=0`$, for $`\alpha 0`$. The case $`n,m<0`$ is similar.
Hence the property of the amalgamated free product applies, and $`W`$ is non-scalar.
*Step* IV. Any word (except the identity) in the Baumslag group $`a,bab^3a^1=b^2`$, of total degree zero in $`a`$, is equal to one of the words
$$a^{\alpha _1}b^{\beta _1}a^{\alpha _2}b^{\beta _2}\mathrm{}b^{\beta _n}a^{\alpha _{n+1}}$$
for which all the Assumptions A1–A3 on consecutive indices, described in Step III, hold. Note that by Proposition 2.4$`(\mathrm{f})`$ we can reduce the proof of the theorem to words of total degree $`0`$ in $`a`$.
To prove the claim of Step IV, the following two lemmas, dealing with easier situations, will be used.
###### Lemma 2.8.
Let $`n1`$ and $`k2`$. Then $`a^nb^ka^n`$ is equal to a product of the form
$$b^{\theta _0}a^{i_1}b^{\theta _1}a^{i_2}\mathrm{}a^{i_p}b^{\theta _p}a^{i_{p+1}}$$
for some strictly positive numbers $`i_1,\mathrm{},i_{p+1}\{1,2,\mathrm{},n\}`$, $`i_{p+1}n`$, and $`\theta _0`$ is nonzero, $`\left|\theta _0\right|3`$, $`\theta _1,\mathrm{},\theta _p=1`$.
For example $`a^nb^2a^n=b^3a^1ba^1b\mathrm{}ba^1ba^{n1}`$, where the product involves $`n`$ occurrences of the letter $`b`$ (not counting powers).
###### Proof.
We start with $`k=2^lq`$, $`q`$ odd. Then $`a^nb^ka^n=a^{(nl)}b^{3^lq}a^{nl}`$. We then split this product as:
$$\left(a^{(nl)}b^{3^lq1}a^{nl}\right)\left(a^{(nl)}ba^{(nl)}\right).$$
We repeat this procedure with
$$3^lq1=2^{l_1}q_1,\text{ }q_1\text{ odd }$$
and will obtain
$$a^nb^ka^n=a^{(nll_1)}b^{3^{l_1}q_1}a^{(nll_1)}a^{(nl)}ba^{(nl)}.$$
By repeating this procedure and stopping when running out of powers of $`a`$, we get the required result. ∎
Similarly one proves:
###### Lemma 2.9.
Let $`n1`$, $`k3`$. Then $`a^nb^ka^n`$ is equal to a product of the form
$$b^{\theta _0}a^{\epsilon _1}b^{\theta _1}a^{\epsilon _2}b^{\theta _2}\mathrm{}a^{\epsilon _s}b^{\theta _s}a^{\epsilon _{s+1}},$$
where $`0<\epsilon _{s+1}n`$, $`\epsilon _1,\epsilon _2,\mathrm{},\epsilon _s>0`$. Moreover, $`\theta _02`$ and $`\theta _s\{1,2\}`$ for $`s1`$.
###### Proof of the claim in Step IV.
We start with one arbitrary word
$$W=b^{k_0}a^{n_1}b^{k_1}a^{n_2}b^{k_2}\mathrm{},n_i0,k_i0,i1,$$
where no obvious cancellations are possible. By moving from the left to the right we look at the first change in sign in the sequence $`n_1,n_2,\mathrm{}`$. Say this occurs when $`i=i_0`$. At that point, if $`\left|k_{i_0}\right|2`$ when $`n_{i_0}<0`$, $`n_{i_0+1}>0`$, or if $`\left|k_{i_0}\right|3`$ when $`n_{i_0}>0`$, $`n_{i_0+1}<0`$, we apply one of the two preceding lemmas to replace $`a^{n_{i_0}}b^{k_{i_0}}a^{n_{i_0+1}}`$ by one of the sequences described in the lemmas. More precisely, if, for example, $`n_{i_0}<0`$, $`n_{i_0+1}>0`$, we apply Lemma 2.8 for
$$x=a^{\mathrm{min}(n_{i_0},n_{i_0+1})}b^{k_{i_0}}a^{\mathrm{min}(n_{i_0},n_{i_0+1})}.$$
By replacing $`x`$ in the word by the form given in Lemma 2.8, the structure of the word up to the next power of $`a`$ following $`b^{k_{i_0}}`$ would fulfill the requirements of the claim. The only case, when in doing this replacement, a change of structure could occur in the structure of the word, before $`a^{n_{i_0}}`$, is when $`n_{i_01}>0`$. But in this case $`\left|k_{i_01}\right|2`$ so $`b^{k_{i_01}}`$ won’t cancel the $`b^3`$ appearing at the beginning of the word from Lemma 2.8. Here we reiterate the procedure. A similar argument works for $`n_{i_0}>0`$, $`n_{i_0+1}<0`$.
By induction, this completes the proof of Step IV. By Steps III and IV we conclude the proof of our theorem. ∎
## 3. Extremal finite von Neumann algebras
In this section we consider the structure of the set of moments of families of projections in a finite von Neumann algebra. Note that by Kirchberg’s technique \[Ki\], for Connes’s conjecture to be true, one should prove that the closure of this set is independent of the finite von Neumann algebra for which we consider the set of moments.
###### Definition 3.1.
Let $`M`$ be a finite separable von Neumann algebra and let $`\tau `$ be a faithful, normalized trace on $`M`$. For any integer $`n1`$, let $`K_M^n`$ be the subset of $`[0,1]^{\frac{n(n+1)}{2}}`$ consisting of the following ordered pairs:
$$K_M^n=\{(\tau (e_ie_j))_{1ijn}(e_1,e_2,\mathrm{},e_n)(𝒫(M))^n\}.$$
###### Proposition 3.2.
Let $`M`$ be a type $`\mathrm{II}_1`$ factor with trace $`\tau `$. Then for all integers $`n1`$,
1. $`K_M^n`$ is convex if $`(M)=_+\backslash \{0\}`$,
2. $`K_M^n`$ is closed under pointwise multiplication, if $`MMM`$,
3. $`K_{M^\omega }^n`$ is closed in the standard topology of $`[0,1]^{\frac{n(n+1)}{2}}`$,
4. $`K_M^n`$ $`K_R^n`$, where $`R`$ is the hyperfinite $`\mathrm{II}_1`$ factor.
The proof of this proposition is identical to the proof of the properties for the set of moments associated with unitaries in a $`\mathrm{II}_1`$ factor. Note that by Kirchberg’s results \[Ki\], $`\overline{K_M^n}=\overline{K_R^n}`$ for all $`n`$, if and only if $`MR^\omega `$.
It is very easy to describe the geometry of a diffuse abelian von Neumann algebra. Indeed,
###### Proposition 3.3.
Let $`Y_n[0,1]^{\frac{n(n+1)}{2}}`$ consist of all $`(\epsilon _{ij})_{1ijn}`$ such that there are sets $`A_1,\mathrm{},A_nX`$, $`X`$ nonvoid, $`A_i=\mathrm{}`$ or $`A_i=X`$ such that $`\epsilon _{ij}=1`$ if $`A_iA_j=X`$ and $`\epsilon _{ij}=0`$ if $`A_iA_j=\mathrm{}`$. Then
$$K_{L^{\mathrm{}}([0,1])}^n=coY_n.$$
In this section we analyze the structure of the closed convex subsets $`K_{M^\omega }^n[0,1]^{\frac{n(n+1\}}{2}}`$. To determine completely this set it would be sufficient to know, for all choices of real numbers $`(a_{ij})_{1ijn}`$ of the value of
$$\underset{1ijn}{\mathrm{max}}\{a_{ij}\lambda _{ij}(\lambda _{ij})K_M^n\}.$$
This is difficult to handle, but we are able to prove at least one geometrical property related to this maximum value: a type of separation of variables at maximum points in $`K_M^n`$.
The following lemma is an easy consequence of the fact that whenever a maximum point is attained at $`(e_1^0,\mathrm{},e_n^0)`$, then for any other projection $`e_1e_1^0`$ or $`e_1e_1^0`$ we get a lower value.
###### Lemma 3.4.
Fix
$$(\tau (e_i^0e_j^0))_{1ijn}$$
a maximum point for the fixed functional
$$L(\lambda _{ij})=\underset{1ijn}{}a_{ij}\lambda _{ij}\text{ on }K_M^n.$$
For all $`i=1,2,\mathrm{},n`$, let
$$\mathrm{\Omega }_i=\mathrm{\Omega }_i(e_1^0,\mathrm{},e_0^n)=\underset{ji}{}a_{ij}e_j^0+a_{ii}e_i^0.$$
Then $`e_i^0\mathrm{\Omega }_ie_i^00`$ and $`(1e_i^0)\mathrm{\Omega }_i(1e_i^0)0`$ for all $`i=1,2,\mathrm{},n`$.
###### Proof.
Fix $`i`$ in $`\{1,2,\mathrm{},n\}`$ and let $`e_i`$ be any projection less than $`e_i^0`$. The fact that
$$\underset{1ijn}{}a_{ij}\tau \left(e_i^0e_j^0\right)$$
is a maximum value for $`L`$ on $`K_M^n`$, implies that
$$\underset{ji}{}a_{ij}\tau \left((e_i^0e_i)e_j^0\right)+a_{ii}\tau (e_i^0e_i)0.$$
Thus for any projection $`e`$ less than $`e_i^0`$ we have that
$$\tau \left(e\left(\underset{ji}{}a_{ij}e_j+a_{ii}Id\right)\right)0.$$
But this gives exactly that
$$e_i^0\mathrm{\Omega }_ie_i^00.$$
Similarly for $`1e_i^0`$. ∎
###### Corollary 3.5.
If $`\left(\lambda _{ij}^0\right)_{1ijn}`$ in $`K_M^n`$ is a maximum point for
$$\left(\lambda _{ij}\right)_{1ijn}K_M^n\underset{1ijn}{}a_{ij}\lambda _{ij},$$
then for all $`i=1,2,\mathrm{},n`$ we have that
$$0\underset{ji}{}a_{ij}\lambda _{ij}^0+a_{ii}\lambda _{ii}^0\underset{ji}{}a_{ij}\lambda _{jj}^0+a_{ii}.$$
###### Proof.
This follows by writing down explicitly that
$$\tau (e_i^0\mathrm{\Omega }_i)0\text{}\tau \left(\left(1e_i^0\right)\mathrm{\Omega }_i\right)0.$$
We will use a method similar to the method of Lagrange multipliers to determine the finer structure of a set of projections $`e_i^0,\mathrm{},e_n^0`$ at which a maximum point is attained.
To do this we need to show that the grassmanian manifold associated with a type $`\mathrm{II}_1`$ factor is large enough.
###### Lemma 3.6.
Let $`M`$ be a $`\mathrm{II}_1`$ factor, $`e`$ be a non-trivial projection in $`M`$ and $`𝒯_e`$ be the linear space consisting of all $`Z`$ in $`M`$, such that $`Z=Z^{}`$ and $`eZe=0`$, $`(1e)Z(1e)`$. Let $`\stackrel{}{𝒯}_e`$ be the set of all $`Z`$ in $`𝒯_e`$ such that there exists a one-parameter family $`e(t)`$ of projections in $`M`$, weakly differentiable at $`0`$, such that
$$e(0)=e,\dot{e}(0)=Z.$$
Then the space of $`\stackrel{}{𝒯}_e`$ is weakly dense in $`𝒯_e`$.
###### Proof.
Assume first that $`\tau (e)=\frac{1}{2}`$, and let $`v`$ be any partial isometry mapping $`e`$ onto $`1e`$. We will show that $`Z=v+v^{}`$ belongs to $`\stackrel{}{𝒯}_e`$.
Indeed $`\{e,v\}^{\prime \prime }`$ can be identified with $`M_2()`$ in such a way that
$$v+v^{}=\left(\begin{array}{cc}0\hfill & 1\hfill \\ 1\hfill & 0\hfill \end{array}\right),\text{ }e=\left(\begin{array}{cc}0\hfill & 0\hfill \\ 0\hfill & 1\hfill \end{array}\right).$$
But then we take
$$e(\theta )=\left(\begin{array}{cc}\mathrm{sin}^2\theta & \mathrm{sin}\theta \mathrm{cos}\theta \\ \mathrm{sin}\theta \mathrm{cos}\theta & \mathrm{cos}^2\theta \end{array}\right).$$
If the trace of $`e`$ is different from $`\frac{1}{2}`$, we may then assume that $`\tau (e)<\frac{1}{2}`$. By the above argument, any partial isometry $`v`$ mapping $`e`$ into a projection under $`1e`$, determines an element $`v+v^{}`$ in $`\stackrel{}{𝒯}_e`$.
Thus $`\stackrel{}{𝒯}_e`$ contains $`v+v^{}`$ for any partial isometry $`v`$, such that $`v^{}v=e`$, $`vv^{}1e`$. Let $`u`$ be any unitary in $`eMe`$ and let $`w`$ be any unitary in $`\left(1e\right)M\left(1e\right)`$. The same argument shows that $`wvu+u^{}v^{}w^{}`$ belongs to $`\stackrel{}{𝒯}_e`$. Since any element in $`eMe`$ and $`\left(1e\right)M\left(1e\right)`$ is a linear combination of unitaries, this shows that $`yvx+x^{}v^{}y^{}`$ is always in the linear span of $`\stackrel{}{𝒯}_e`$ for all $`x`$ in $`eMe`$ and $`y`$ in $`\left(1e\right)M\left(1e\right)`$. This set is obviously weakly dense in $`𝒯_e`$. ∎
###### Corollary 3.7.
Fix $`n1`$ and real numbers $`(a_{ij})_{1ijn}`$. Let
$$(e_1^0,e_2^0,\mathrm{},e_n^0)$$
be a family of projections in $`𝒫(M)`$ such that the maximum of
$$L\left(\left(\lambda _{ij}\right)_{1ijn}\right)=\underset{1ijn}{}a_{ij}\lambda _{ij}$$
for $`\left(\lambda _{ij}\right)`$ in $`K_M^n`$ is attained at $`(\tau (e_i^{}e_j^{}))`$. Let
$$\mathrm{\Omega }_i^0=\mathrm{\Omega }_i^0(e_i^0,\mathrm{},e_n^0,a_{ij})=a_{ii}Id+\underset{ji}{}a_{ij}e_j^0.$$
Then $`[e_i^0,\mathrm{\Omega }_i^0]=0`$. By using Lemma 3.4 it follows that
$`e_i^0`$ $`supp(\mathrm{\Omega }_i^0)_+,`$
$`1e_i^0`$ $`supp\left(\mathrm{\Omega }_i^0\right)_{}.`$
###### Proof.
Indeed if
$$(\tau (e_i^0e_j^0))_{1ijn}$$
is such a maximum point for the functional $`L`$ on $`K_M^n`$, then for all $`Z_i`$ in $`\stackrel{}{𝒯}_{e_i^0}`$ we have that
$$\tau (\mathrm{\Omega }_i^0Z_i)=0.$$
But then this will give that
$$\tau (\mathrm{\Omega }_i^0Z_i)=0$$
for all $`Z_i=(1e_i^0)Y_1e_i^0+e_i^0Y_1(1e_i^0)`$, $`Y_1M_{sa}`$. Thus for all $`Y=Y^{}`$ in $`M`$ we have
$$\tau \left(\mathrm{\Omega }_i^0e_i^0Y(1e_i^0)+\mathrm{\Omega }_i^0(1e_i^0)Ye_i^0\right)=0$$
and hence
$$\tau \left([(1e_i^0)\mathrm{\Omega }_i^0e_i^0+e_i^0\mathrm{\Omega }_i^0(1e_i^0)]Y\right)=0$$
for all $`Y`$ selfadjoint in $`M`$. Since $`\mathrm{\Omega }_i^0`$ is also selfadjoint, it follows that
$$(1e_i^0)\mathrm{\Omega }_i^0e_i^0+(1e_i^0)\mathrm{\Omega }_i^0e_i^0=0$$
and hence that $`\mathrm{\Omega }_i^0`$ commutes with $`e_i^0`$. ∎
###### Remark 3.8.
The above proposition suggests that for Connes’s embedding problem, it is sufficient to consider finite von Neumann algebras (which we call *extremal finite von Neumann algebras*) that are generated by families of projections $`e_1,e_2,\mathrm{},e_n`$ such that there exists a matrix of real numbers $`(a_{ij})_{1ijn}`$ with the following property.
For each $`i`$, let
$$\mathrm{\Omega }_i=\underset{ji}{}a_{ij}e_j+a_{ii}Id.$$
Let $`s_+^i`$ be the init of the positive part of $`(\mathrm{\Omega }_i)_+`$ and $`s_{}^i`$ be the projection onto the init space of $`(\mathrm{\Omega }_i)_{}`$.
Then $`1s_{}^ie_is_+^i`$; in particular, $`e_i`$ commutes with $`\mathrm{\Omega }_i`$, for all $`i`$.
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# Twisted 𝐺𝐿_𝑛 Loop Group Orbit and Solutions of the WDVV Equations
## 1 Introduction
In the early 90’s B. Dubrovin noticed that the local classification of massive topological field theories can be solved by classifying certain flat diagonal metrics
$$ds^2=\underset{i=1}{\overset{n}{}}h_i^2(u)(du_i)^2,u=(u_1,\mathrm{},u_n),$$
(1.1)
with
$$_jh_i^2(u)=_ih_j^2(u),\text{and }\underset{k=1}{\overset{n}{}}_kh_i(u)=0,_i=\frac{}{u_i}.$$
(1.2)
For the flat coordinates $`t^i`$, $`1in`$, of the metric (1.1), the functions
$$c_k\mathrm{}^m(t)=\underset{i=1}{\overset{n}{}}\frac{t^m}{u_i}\frac{u_i}{t^k}\frac{u_i}{t^{\mathrm{}}},$$
(1.3)
satisfy
$$\underset{k=1}{\overset{n}{}}c_{ij}^k(t)c_{km}^{\mathrm{}}(t)=\underset{k=1}{\overset{n}{}}c_{jm}^k(t)c_{ik}^{\mathrm{}}(t).$$
(1.4)
If one writes down these equations for the function $`F(t)`$ for which
$$\frac{^3F(t)}{t^kt^{\mathrm{}}t^m}=c_{k\mathrm{}m}(t)=\underset{i=1}{\overset{n}{}}\eta _{mi}c_k\mathrm{}^i(t),\text{where}\eta _{pq}=\underset{i=1}{\overset{n}{}}h_i^2(u)\frac{u_i}{t^p}\frac{u_i}{t^q},$$
(1.5)
with the constraint
$$\frac{^3F(t)}{t^1t^{\mathrm{}}t^m}=\eta _\mathrm{}m,$$
one obtains the well-known the Witten-Dijkgraaf-E. Verlinde-H. Verlinde (WDVV)-equations , .
Vanishing of the curvature of these metrics (1.1) can be written in the form of a system of partial differential equations in the cannonical coordinates $`u_i`$ for the rotation coefficients
$$\gamma _{ij}=\frac{_jh_i(u)}{h_j(u)},ij,$$
(1.6)
which is known under the name the Darboux-Egoroff system:
$`\gamma _{ij}(u)`$ $`=\gamma _{ji}(u),`$ (1.7)
$`_k\gamma _{ij}(u)`$ $`=\gamma _{ik}(u)\gamma _{kj}(u),`$
$`{\displaystyle \underset{k=1}{\overset{n}{}}}_k\gamma _{ij}(u)`$ $`=0,`$
From (1.2) and (1.6), we see that the Lamé coefficients $`h_i`$ satisfy:
$$_jh_i(u)=\gamma _{ij}(u)h_j(u),ij,_ih_i(u)=\underset{ji}{}\gamma _{ij}(u)h_j(u).$$
(1.8)
The flat coordinates $`t^1,\mathrm{},t^n`$ of this metric can be found from the linear system
$$_i_jt^k=\mathrm{\Gamma }_{ij}^i_it^k+\mathrm{\Gamma }_{ji}^j_jt^k,ij;_i_it^k=\underset{j=1}{\overset{n}{}}\mathrm{\Gamma }_{ii}^j_jt^k,$$
(1.9)
where $`\mathrm{\Gamma }_{ij}^k`$ are Christoffel symbols:
$$\mathrm{\Gamma }_{ij}^i=\frac{_jh_i}{h_i},\mathrm{\Gamma }_{ii}^j=(2\delta _{ij}1)\frac{h_i_jh_i}{h_j^2}.$$
(1.10)
R. Martini and the author constructed in solutions of the Darboux-Egoroff system (1.7). These solutions were related to certain points in the $`SL_n`$ Loop group orbit of the highest weight vector of the homogeneous realization of the basic representation and were related to a reduction of the $`n`$-component KP hierarchy. They had no real representation theoretical explanation for these particular solution. In this paper we show that all elements in the twisted Loop group orbit of $`GL_n`$ lead to solutions of the Darboux-Egoroff system. The solutions of are certain homogeneous solutions in this orbit. This makes it possible to construct non-homogeneous WDVV prepotentials $`F`$. Inspired by the papers and , we construct in section 5 for all elements in the twisted Loop group orbit besides solutions of the Darboux-Egoroff system, also the related (deformed) flat coordinates and the WDVV prepotential $`F`$.
## 2 Groups and Grassmannians
Consider the space $`H_n`$ of Laurent series in $`t`$ (this $`t`$ has nothing to do with the flat coordinates appearing in the previous section) with coefficients in $`^n`$
$$H_n=\{\underset{j}{}c_jt^j|c_j^n,c_j=0\text{for }j<<0\}.$$
Let $`e_j`$, $`1jn`$, be a basis of $`^n`$, Define
$$v_{n(k+1)+j\frac{1}{2}}=v_{k\frac{1}{2}}^{(j)}:=e_jt^k,1jn,k,$$
(2.1)
then an element of $`H_n`$ is a unique linear combination of, possibly infintely many, $`v_{\mathrm{}}`$’s or equivalently $`v_{\mathrm{}}^{(j)}`$’s. The space $`H_n`$ has a natural filtration ($`j)`$
$$\mathrm{}H_n^{(j1)}H_n^{(j)}H_n^{(j+1)}\mathrm{},$$
where
$$H_n^{(j)}=\{\underset{k}{}b_kv_k|b_k,b_k=0\text{for }k>j\}.$$
(2.2)
Let $`H_n^{}`$ denote the space of linear functions $`f`$ on $`H_n`$ such that $`f(H_n^{(j)})=0`$ for $`j<<0`$. On the direct sum $`\overline{H}_n=H_nH_n^{}`$ of these two spaces, one has a natural symmetric nondegenerate bilinear form $`(,)`$ for which the spaces $`H_n`$ and $`H_n^{}`$ are isotropic, it is defined by $`(f,v)=f(v)`$ for $`vH_n`$ and $`fH_n^{}`$. Define ”dual basis” elements to the elements defined in (2.1) as follows
$$v_{nkj+\frac{1}{2}}^{}=v_{k\frac{1}{2}}^{(j)}=e_j^{}t^k,$$
(2.3)
then the bilinear form is given by the following formula’s
$$(e_i^{}t^{\mathrm{}},e_jt^k)=\delta _{k+\mathrm{},1}\delta _{ij},(v_r^{(i)},v_s^{(j)})=\delta _{r,s}\delta _{ij},(v_r^{},v_s)=\delta _{r,s}.$$
(2.4)
Let $`^n=_{j=1}^ne_j^{}`$ then clearly
$$H_n^{}=\{\underset{j}{}c_j^{}t^j|c_j^{}^n,c_j^{}=0\text{for }j<<0\}.$$
In analogy with (2.2), one also has the subspaces
$$H_n^{(j)^{}}=\{\underset{k}{}b_kv_k^{}|b_k,b_k=0\text{for }k>j\}.$$
(2.5)
Let $`(|)`$ be the symmetric bilinear form on $`^n^n`$ for which $`^n`$ and $`^n`$ are isotropic and $`(e_i^{}|e_j)=\delta _{ij}`$. Let $`v(t)H_n`$ and $`f(t)H_n^{}`$, then (2.4) is equivalent to
$$(f(t),v(t))=\text{Res}_t(f(t)|v(t)),$$
(2.6)
where $`\text{Res}_ta_it^i=a_1`$. Notice that the left hand side of (2.6) makes sense.
Now, notice that the space $`\overline{H}_n^{(j)}:=H_n^{(j)}H_n^{(j)^{}}`$ is a maximal isotropic subspace of $`\overline{H}_n`$ with respect to our bilinear form $`(,)`$. Having these specific isotropic spaces in mind, we want to define more general maximal isotropic subspaces. Let $`\overline{V}\overline{H}_n`$ be a maximal isotropic linear subspace, which satisfies the following conditions:
$$\overline{V}=VV^{}\text{with }VH_n,VH_n^{}$$
(2.7)
and
$$H_n^{(j)}V,H_n^{(j)^{}}V^{}\text{for }j<<0.$$
(2.8)
All such subspaces form the points of an infinite (isotropic) Grassmannian $`\overline{Gr}`$. Clearly, since (2.7) holds, a $`\overline{V}\overline{Gr}`$ induces two unique spaces $`V`$ and $`V^{}`$, which can be separately regarded as points of two Grassmannians
$$Gr=\{VH_n|H_n^{(j)}V,\text{for }j<<0\}\text{and }Gr^{}=\{V^{}H_n^{}|H_n^{(j)}V,\text{for }j<<0\}.$$
It is obvious that the converse also holds, i.e., if $`VGr`$ (or $`V^{}Gr^{}`$), then there exists a unique maximal space $`V^{}Gr^{}`$ (resp. $`VGr`$) such that $`(V^{},V)=0`$ and hence $`\overline{V}=VV^{}`$ is a unique point of $`\overline{Gr}`$.
The space $`Gr`$ is Sato’s polynomial Grassmannian . Here it is coupled to $`Gr^{}`$ and $`\overline{Gr}`$. The reason why we introduce the latter two spaces will become clear later on. It is well-known that the space $`Gr=_mGr_m`$, disjoint union of the spaces
$$Gr_m=\{VGr|H_n^{(j)}V\text{and dim }V/H_n^{(j)}=mj\text{for }j<<0.\}$$
If $`VGr_m`$, then $`V^{}`$ belongs to
$$Gr_m^{}=\{V^{}Gr|H_n^{(j)}V^{}\text{and dim }V^{}/H_n^{(j)}=mj\text{for }j<<0.\}$$
In this situation, $`\overline{V}=VV^{}`$ is an element of
$$Gr_m=\{\overline{V}=VV^{}\overline{Gr}|VGr_m\text{and }V^{}Gr_m^{}\}.$$
Suppose that the space $`V`$ is invariant under multiplication with $`t`$, i.e., $`V`$ is an element of the restricted Grassmannian
$$gr:=\{VGr|tVV\},$$
then it is clear from the above construction that its ”dual space” $`V^{}`$ is unique. Now let $`f(t)V^{}`$, then for all $`v(t)V)`$ we have $`\text{Res}_t(f(t)|v(t))=0`$. Since $`tv(t)V`$, also
$$\text{Res}_t(tf(t)|v(t))=\text{Res}_t(f(t)|tv(t))=0\text{for all }v(t)V,$$
hence $`tf(t)V^{}`$. This means that also $`V^{}`$ and $`\overline{V}=VV^{}`$ satisfy $`tV^{}V^{}`$, $`t\overline{V}\overline{V}`$ and that it makes sense to define
$$gr^{}:=\{V^{}Gr^{}|tV^{}V^{}\}\text{and }\overline{gr}:=\{\overline{V}\overline{Gr}|t\overline{V}\overline{V}\}$$
and
$$gr_m:=grGr_m,gr_m^{}:=gr^{}Gr_m^{}\text{and }\overline{gr}_m:=\overline{gr}\overline{Gr}_m.$$
Consider $`H_n^{(j)}`$, $`j`$, to be the a fundamental system of neighborhoods of zero, then $`H_n`$ becomes a topological vector space. let $`\overline{a}_{\mathrm{}}`$ be the algebra of all continuous endomorphisms of $`H_n`$. If one considers an elmenent of $`H_n`$ as an infinite vector with respect to the basis $`v_i`$, $`i\frac{1}{2}+`$, then (see )
$$\overline{a}_{\mathrm{}}=\{(a_ij)_{i,j\frac{1}{2}+}|\text{for each }k\text{the number of non-zero }a_{ij}\text{with }ik\text{and }jk\text{is finite}\}.$$
(2.9)
Denote by $`\overline{A}_{\mathrm{}}`$ the group of invertible elements of the associative algebra $`\overline{a}_{\mathrm{}}`$. Then $`\overline{A}_{\mathrm{}}`$ acts trasitively on $`Gr`$. We can extend an element $`g\overline{A}_{\mathrm{}}`$ to an element, which by abuse of notation we denote by the same letter, g of $`O(\overline{V})`$, the orthogonal group of $`\overline{V}`$ (see ):
###### Lemma 2.1
Let $`g\overline{A}_{\mathrm{}}`$ be such that
$$gv_j=\underset{i+\frac{1}{2}}{}A_{ij}v_i$$
and let $`A=(A_{ij})_{i,j+\frac{1}{2}}`$ and $`A^1=(B_{ij})_{i,j+\frac{1}{2}}`$. Then
$$gv_j^{}=\underset{i+\frac{1}{2}}{}B_{ji}v_i^{}.$$
Proof. Suppose that $`gv_{\mathrm{}}^{}=_kB_{\mathrm{},k}v_k^{}`$, then it follows from
$$\begin{array}{ccc}\delta _{j,\mathrm{}}=(v_j,v_{\mathrm{}}^{})\hfill & =& (_iA_{ij}v_i,B_\mathrm{}kv_k^{})\hfill \\ & =& _iA_{ij}B_{\mathrm{},i}\hfill \end{array}$$
that $`B=A^1`$. $`\mathrm{}`$ This lemma defines the (transitive) action of $`\overline{A}_{\mathrm{}}`$ on $`Gr^{}`$ and hence on $`\overline{Gr}`$.
Let $`L=[t,t^1]`$ be the algebra of Laurent polynomials in $`t`$. The identification (2.1) gives us an embedding
$$\varphi :\text{Mat}_n(L)\overline{a}_{\mathrm{}},\varphi (e_{ij}t^{\mathrm{}})=\underset{k}{}E_{nk+i\frac{1}{2},n(k+\mathrm{})+j\frac{1}{2}},$$
(2.10)
where $`e_{ij}=(\delta _{ir}\delta _{js})_{1r,sn}\text{Mat}_n()`$ and $`E_{ij}=(\delta _{ir}\delta _{js})_{r,s\frac{1}{2}+}`$. This embedding gives rise to the embedding of the loop Lie algebra $`gl_n(L)`$ into $`\overline{a}_{\mathrm{}}`$ and the loop group $`GL_n(L)`$ in $`\overline{A}_{\mathrm{}}`$. Multiplication with $`t`$, i.e. with $`_{i=1}^ne_{ii}t`$, commutes with the action of $`GL_n(L)`$ and thus $`GL_n(L)`$ acts on $`gr`$. This action is transitive, see . Moreover, $`\varphi (t)=\varphi (_{i=1}^ne_{ii}t)=_{k\frac{1}{2}+}E_{k,k+n}`$.
Consider $`H_n^{(j)}`$, it has as a ”natural basis” the elements $`v_i`$ with $`i<j`$. Its dual space $`H_n^{(j)}`$ has $`v_i^{}`$ with $`i<j`$ as ”basis”. To these spaces we associate special vectors in two semi-infinite wedge spaces (see , )
$`|H_n^{(j)}`$ $`=v_{j\frac{1}{2}}v_{j\frac{3}{2}}v_{j\frac{5}{2}}v_{j\frac{7}{2}}\mathrm{},`$ (2.11)
$`H_n^{(j)^{}}|`$ $`=\mathrm{}v_{j\frac{7}{2}}^{}v_{j\frac{5}{2}}^{}v_{j\frac{3}{2}}^{}v_{j\frac{1}{2}}^{}.`$
In fact we can associate to any element $`WGr_j`$ and its dual space $`W^{}Gr_j^{}`$ elements in these wedge spaces, viz., we know that there exists an $`m<<0`$ such that $`H_n^{(m)}W`$ and $`H_n^{(m)}W^{}`$. Then let $`w_{j\frac{1}{2}},w_{j\frac{3}{2}},\mathrm{},w_{m+\frac{1}{2}}`$ be a basis of $`W\text{mod}H_n^{(m)}`$ and $`w_{j\frac{1}{2}}^{},w_{j\frac{3}{2}}^{},\mathrm{},w_{m+\frac{1}{2}}^{}`$ be a basis of $`W^{}\text{mod}H_n^{(m)}`$ then we put
$`|W`$ $`=w_{j\frac{1}{2}}w_{j\frac{3}{2}}\mathrm{}w_{m+\frac{1}{2}}v_{m\frac{1}{2}}v_{m\frac{3}{2}}v_{m\frac{5}{2}}\mathrm{},`$ (2.12)
$`W^{}|`$ $`=\mathrm{}v_{m\frac{5}{2}}^{}v_{m\frac{3}{2}}^{}v_{m\frac{1}{2}}^{}w_{m+\frac{1}{2}}^{}w_{m+\frac{3}{2}}^{}\mathrm{}w_{j\frac{1}{2}}^{}.`$
It is clear that upto a constant both $`|W`$ and $`W^{}|`$ are independent of the choice of basis of $`W\text{mod}H_n^{(m)}`$ and $`W^{}\text{mod}H_n^{(m)}`$. This gives us a map $`\mu `$ from $`Gr`$ and $`Gr^{}`$ into $`F`$ and $`F^{}`$, where $`F=_jF^{(j)}`$, $`F^{}=_jF^{(j)}`$ and the spaces $`F^{(j)}`$ and $`F^{(j)}`$ are the vector spaces generated by the elements on the right-hand side of (2.12) (see , or for more details). Since $`\overline{A}_{\mathrm{}}`$ and its subgroup $`GL_n(L)`$ act on $`Gr`$ and $`Gr^{}`$ we obtain a projective representation of these groups on $`F`$ and $`F^{}`$.
Let $`\omega `$ be the following orthogonal transformation of $`\overline{V}`$ with respect to the bilinear form (2.4):
$$\omega (v_m^{(a)})=()^{m+\frac{1}{2}}iv_m^{(a)^{}},\omega (v_m^{(a)})=()^{m+\frac{1}{2}}iv_m^{(a)},m\frac{1}{2}+,1an,$$
(2.13)
then $`\omega (\overline{H}_n^{(0)})=\overline{H}_n^{(0)}`$ and $`\omega (\overline{Gr}_m)=\overline{Gr}_m`$. If $`g\overline{A}_{\mathrm{}}`$ is such that
$$gv_j^{(b)}=\underset{a,k}{}A_{kj}^{(ab)}v_k^{(a)}\text{then }gv_j^{(b)}=\underset{k,a}{}B_{j,k}^{(b,a)}v_k^{(a)}\text{with }\underset{j,b}{}A_{kj}^{(ab)}B_j\mathrm{}^{bc}=\delta _k\mathrm{}\delta _{ab},$$
(2.14)
then
$`\omega (gv_j^{(b)})=`$ $`\omega (g)\omega (v_j^{(b)})=\omega (g)()^{j+\frac{1}{2}}iv_j^{(b)}`$
$`=`$ $`{\displaystyle \underset{a,k}{}}A_{kj}^{(ab)}\omega (v_k^{(a)})={\displaystyle \underset{a,k}{}}A_{kj}^{(ab)}()^{k+\frac{1}{2}}iv_k^{(a)}.`$
This and a similar calculation for $`gv_j^{(b)}`$ gives
$$\omega (g)v_j^{(b)}=\underset{a,k}{}()^{kj}B_{j,k}^{(ba)}v_k^{(a)},\omega (g)v_j^{(b)}=\underset{a,k}{}()^{kj}A_{kj}^{(ab)}v_k^{(a)},$$
(2.15)
and hence induces on $`GL_n(L)`$ the automorphism:
$$\omega (A(t))=((A(t))^T)^1,A(t)\text{Mat}_n(L),$$
(2.16)
where $`A^T`$ stand for the transposed of the matrix $`A`$. The fixed point set of this automorphism is the twisted loop group
$$GL_n(L)^{(2)}=\{A(t)GL_n(L)|A(t)=((A(t))^T)^1\}.$$
All this suggests to define a new skew- symmetric bilinear form $`,`$ on $`H_n`$
$$u,v=(u,\omega (v))u,vH_n,$$
or in other terms cf. (2.7)
$$u(t),v(t)=i\text{Res}_t(u(t)|v(t))u(t),v(t)H_n,$$
(2.17)
and new Grassmannians
$$Gr^{(2)}=\{WGr_0|u,v=0\text{for all }u,vW\},gr^{(2)}=grGr^{(2)}.$$
We now want to show that $`gr^{(2)}`$ the homogeneous space is of the twisted loop group $`GL_n(L)^{(2)}`$. It is obvious, from the above discussion, that $`gH_n^{(0)}gr^{(2)}`$ for any $`gGL_n(L)^{(2)}`$. To prove the converse, we will use the following Theorem of :
###### Theorem 2.1
Any loop in $`g(t)Gl_n(L)`$ can be uniquely factorized as $`g(t)=g_u(t)g_+(t)`$, with $`g_+(t)GL_n([t])`$ and $`g_u\mathrm{\Omega }U_n=\{h(t)U_n(L)|h(1)=I_n\}`$. Moreover, $`gr=\mathrm{\Omega }U_nH_n^{(0)}`$ and isotropy group of $`H_n^{(0)}`$ is the group $`U_n`$ of constant loops.
Now let $`Wgr^{(2)}`$, then we can write $`W=g(t)H_n^{(0)}`$ for certain $`g(t)\mathrm{\Omega }U_n`$. Clearly also $`\omega (g(t))H_n^{(0)}=W`$, hence $`(g(t))^1\omega (g(t))=uU_n`$. Since $`\omega `$ is an involution on $`GL_n(L)`$ we find that $`u\omega (u)=1`$, and since $`u`$ is unitairy we deduce that the following conditions hold for u:
$$u=u^T\text{and }u\overline{u}^T=I_n.$$
So we can find a real orthogonal matrix $`v`$ such that $`u=vdv^T`$, with $`d`$ a diagonal matrix. Let $`s`$ be a diagonal matrix that satisfies $`s^2=d`$, then $`w=w^T=vsv^T`$ is a unitary matrix that satisfies $`u=w^2`$ and
$$\omega (g(t)w)=g(t)u\omega (w)=g(t)w^2(w^T)^1=g(t)w^2w^1=g(t)w.$$
Since also $`W=g(t)H_n^{(0)}=g(t)wH_n^{(0)}`$, we find that $`WGL_n(L)^{(2)}H_n^{(0)}`$ and we have proven the following
###### Theorem 2.2
$$gr^{(2)}=GL_n(L)^{(2)}H_n^{(0)}.$$
## 3 The Clifford Algebra and Tau-Functions
Using (2.11), we can associate to any point $`WGr_j`$ and the corresponding $`W^{}Gr_j^{}`$ vectors $`|WF^{(j)}`$, $`W^{}|F^{(j)}`$ and hence upto a constants unique vectors $`|WF^{(j)}`$ and $`W^{}|F^{(j)}`$. Clearly $`\omega `$ defined in (2.13) is not well defined on these vectors in the spaces $`F`$ and $`F^{}`$. To solve this problem, we will define a Clifford algebra and its corresponding spin module and we will fix $`\omega `$ on one of the vectors of the spin module. In this construction we will only consider the space $`F`$, since the space $`W^{}`$ corresponding to a $`WGr`$ is always unique.
Recall that on the infinite space $`\overline{H}_n=H_nH_n^{}`$ we have a symmetric bilinear form $`(,)`$ given by (2.4). Let $`C\mathrm{}(\overline{H}_n)`$ be the Clifford algebra on this space, i.e. the associative algebra over $``$ with unity $`1`$ wich has the folowing defining relations
$$uv+vu=(u,v)1,u,v\overline{H}_n.$$
(3.1)
We obtain a obtain a representation $`\psi `$ of the clifford algebra on the space $`F`$ by defining it on wedges as follows ($`w,w_jH_nw^{}H_n^{}`$ with $`w_j=v_{jm\frac{1}{2}}`$ for certain $`m`$ and all $`j>>0`$)
$`\psi (w)(w_0w_1`$ $`w_2\mathrm{})=ww_0w_1w_2\mathrm{},`$ (3.2)
$`\psi (w^{})(w_0w_1`$ $`\mathrm{})={\displaystyle }_{i=0}^{\mathrm{}}()^j(w^{},w_j)w_0w_1\mathrm{}w_{j1}w_{j+1}\mathrm{}.`$
The space $`F`$ is the spin module for this Clifford algebra. Let
$$|0=v_{\frac{1}{2}}v_{\frac{3}{2}}v_{\frac{5}{2}}v_{\frac{7}{2}}\mathrm{},$$
then $`F`$ is the unique module generated by
$$\psi (v_k)|0=\psi (v_k^{})|0=0\text{for }k<0.$$
(3.3)
It is straightforward to check that (cf. (2.12))
$`w_{j\frac{1}{2}}w_{j\frac{3}{2}}`$ $`\mathrm{}w_{m+\frac{1}{2}}v_{m\frac{1}{2}}v_{m\frac{3}{2}}\mathrm{}`$ (3.4)
$`=\psi (w_{j\frac{1}{2}})`$ $`\psi (w_{j\frac{3}{2}})\mathrm{}\psi (w_{m+\frac{1}{2}})\psi (v_{m\frac{1}{2}}^{})\psi (v_{m\frac{3}{2}}^{})\mathrm{}\psi (v_{\frac{3}{2}}^{})\psi (v_{\frac{1}{2}}^{})|0.`$
Define the fermionic fields ($`z^\times `$)
$`\psi ^{+(j)}(z)={\displaystyle \underset{k+\frac{1}{2}}{}}\psi _k^{+(j)}z^{k\frac{1}{2}}`$ $`:={\displaystyle \underset{k+\frac{1}{2}}{}}\psi (v_k^{(j)})z^{k\frac{1}{2}},`$ (3.5)
$`\psi ^{(j)}(z)={\displaystyle \underset{k+\frac{1}{2}}{}}\psi _k^{(j)}z^{k\frac{1}{2}}`$ $`:={\displaystyle \underset{k+\frac{1}{2}}{}}\psi (v_k^{(j)})z^{k\frac{1}{2}}`$
and bosonic fields $`(1jn)`$ by
$$\alpha ^{(j)}(z)=\underset{k}{}\alpha _k^{(j)}z^{k1}=:\psi ^{+(j)}(z)\psi ^{(j)}(z):$$
(3.6)
where $`::`$ stands for the normal ordered product defined in the usual way $`(\lambda ,\mu =+`$ or $``$):
$$:\psi _k^{\lambda (i)}\psi _{\mathrm{}}^{\mu (j)}:=\{\begin{array}{cc}\psi _k^{\lambda (i)}\psi _{\mathrm{}}^{\mu (j)}\hfill & \text{if}\mathrm{}k,\hfill \\ \psi _{\mathrm{}}^{\mu (j)}\psi _k^{\lambda (i)}\hfill & \text{if}\mathrm{}<k.\hfill \end{array}$$
(3.7)
Notice that $`\psi ^{+(j)}(z)=\psi (\delta (zt)e_j)`$ and $`\psi ^{(j)}(z)=\psi (\delta (zt)e_j^{})`$, where $`\delta (zt)=z^1_k\left(\frac{z}{t}\right)^k`$. One checks (using e.g. the Wick formula) that these bosonic operators satisfy the canonical commutation relation of the associative oscillator algebra:
$$[\alpha _k^{(i)},\alpha _{\mathrm{}}^{(j)}]=k\delta _{ij}\delta _{k,\mathrm{}},$$
(3.8)
and one has
$$\alpha _k^{(i)}|0=0\text{for}k0.$$
(3.9)
In order to express the fermionic fields $`\psi ^{\pm (i)}(z)`$ in terms of the bosonic fields $`\alpha ^{(i)}(z)`$, we need some additional operators $`Q_i,i=1,\mathrm{},n`$, on $`F`$. These operators are uniquely defined by the following conditions:
$$Q_i|0=\psi _{\frac{1}{2}}^{+(i)}|0,Q_i\psi _k^{\pm (j)}=(1)^{\delta _{ij}+1}\psi _{k\delta _{ij}}^{\pm (j)}Q_i.$$
(3.10)
They satisfy the following commutation relations:
$$Q_iQ_j=Q_jQ_i\text{if}ij,[\alpha _k^{(i)},Q_j]=\delta _{ij}\delta _{k0}Q_j.$$
(3.11)
###### Theorem 3.1
(, )
$$\psi ^{\pm (i)}(z)=Q_i^{\pm 1}z^{\pm \alpha _0^{(i)}}\mathrm{exp}(\underset{k<0}{}\frac{1}{k}\alpha _k^{(i)}z^k)\mathrm{exp}(\underset{k>0}{}\frac{1}{k}\alpha _k^{(i)}z^k).$$
(3.12)
Proof. See .
The operators on the right-hand side of (3.12) are called vertex operators. They made their first appearance in string theory (cf. ).
We can describe now the $`n`$-component boson-fermion correspondence. Let $`[x]`$ be the space of polynomials in indeterminates $`x=\{x_k^{(i)}\},k=1,2,\mathrm{},i=1,2,\mathrm{},n`$. Let $`N`$ be a lattice with a basis $`\delta _1,\mathrm{},\delta _n`$ over $``$ and the symmetric bilinear form $`(\delta _i|\delta _j)=\delta _{ij}`$, where $`\delta _{ij}`$ is the Kronecker symbol. Let
$$\epsilon _{ij}=\{\begin{array}{cc}1\hfill & \text{if }i>j\hfill \\ 1\hfill & \text{if }ij\text{.}\hfill \end{array}$$
(3.13)
Define a bimultiplicative function $`\epsilon :N\times N\{\pm 1\}`$ by letting
$$\epsilon (\delta _i,\delta _j)=\epsilon _{ij}.$$
(3.14)
Let $`\delta =\delta _1+\mathrm{}+\delta _n,M=\{\gamma N|(\delta |\gamma )=0\}`$, $`\mathrm{\Delta }=\{\alpha _{ij}:=\delta _i\delta _j|i,j=1,\mathrm{},n,ij\}`$. Of course $`M`$ is the root lattice of $`s\mathrm{}_n()`$, the set $`\mathrm{\Delta }`$ being the root system.
Consider the vector space $`[N]`$ with basis $`e^\gamma `$$`\gamma L`$, and the following twisted group algebra product:
$$e^\alpha e^\beta =\epsilon (\alpha ,\beta )e^{\alpha +\beta }.$$
(3.15)
Let $`B=[x]_{}[N]`$ be the tensor product of algebras. Then the $`n`$-component boson-fermion correspondence is the vector space isomorphism
$$\sigma :FB,\text{with }\sigma :F^{(m)}B^{(m)}$$
(3.16)
given by
$$\sigma (\alpha _{m_1}^{(i_1)}\mathrm{}\alpha _{m_s}^{(i_s)}Q_1^{k_1}\mathrm{}Q_n^{k_n}|0)=m_1\mathrm{}m_sx_{m_1}^{(i_1)}\mathrm{}x_{m_s}^{(i_s)}e^{k_1\delta _1+\mathrm{}+k_n\delta _n}.$$
(3.17)
The transported action of the operators $`\alpha _m^{(i)}`$ and $`Q_j`$ looks as follows:
$$\{\begin{array}{cc}\sigma \alpha _m^{(j)}\sigma ^1(p(x)e^\gamma )=mx_m^{(j)}p(x)e^\gamma ,\text{if}m>0,\hfill & \\ \sigma \alpha _m^{(j)}\sigma ^1(p(x)e^\gamma )=\frac{p(x)}{x_m}e^\gamma ,\text{if}m>0,\hfill & \\ \sigma \alpha _0^{(j)}\sigma ^1(p(x)e^\gamma )=(\delta _j|\gamma )p(x)e^\gamma ,\hfill & \\ \sigma Q_j\sigma ^1(p(x)e^\gamma )=\epsilon (\delta _j,\gamma )p(x)e^{\gamma +\delta _j}.\hfill & \end{array}$$
(3.18)
The transported action of the fermionic fields is as follows:
$$\sigma \psi ^{\pm (j)}(z)\sigma ^1=e^{\pm \delta _j}z^{\pm \delta _j}\mathrm{exp}(\pm \underset{k=1}{\overset{\mathrm{}}{}}x_k^{(j)})\mathrm{exp}(\underset{k=1}{\overset{\mathrm{}}{}}\frac{}{x_k^{(j)}}\frac{z^k}{k})$$
(3.19)
We will now determine the second part of the boson–fermion correspondence, i.e., we want to determine $`\sigma `$ of the elements (3.4). Since for our purpose we are only interested in $`Gr^{0)}`$ we will assume that this element is a wedge in $`F^{(0)}`$, i.e. let
$$\tau =A_{\frac{1}{2}}A_{\frac{3}{2}}A_{\frac{5}{2}}\mathrm{}F^{(0)}\text{with }A_p=v_p\text{for all }p>P>>0.$$
(3.20)
To such an element we can associate an element in $`A=(A_{ij})\overline{A}_{\mathrm{}}`$ such that $`Av_k=A_k`$ for all $`k>0`$. Notice that $`A_{ij}=\delta _{ij}`$ for $`j<P`$. Then R. Martini and the author showed in the following
###### Proposition 3.1
Let $`\sigma (\tau )=_{\alpha M}\tau _\alpha (x)e^\alpha `$. Assume that $`\alpha =_{j=1}^nk_j\delta _j`$ and suppose that
$$Q_1^{k_1}Q_2^{k_2}\mathrm{}Q_n^{k_n}|0=\lambda _\alpha v_{j_{\frac{1}{2}}}v_{j_{\frac{3}{2}}}v_{j_{\frac{5}{2}}}\mathrm{},$$
with $`j_{\frac{1}{2}}>j_{\frac{3}{2}}>j_{\frac{5}{2}}\mathrm{}`$ and $`j_q=q`$ for all $`q>Q>>0`$ and $`\lambda _\alpha =\pm 1`$, then
$$\tau _\alpha (x)=\lambda _\alpha det\left(\underset{R<\mathrm{}<0}{}\underset{r=j_{\frac{1}{2}},j_{\frac{3}{2}},\mathrm{},j_{R+\frac{1}{2}}}{}\underset{\begin{array}{c}1jn,q+\frac{1}{2}\\ nq\frac{1}{2}(n2j+1)=r\end{array}}{}\left(\underset{k=0}{\overset{\mathrm{}}{}}A_{r+nk,\mathrm{}}S_k(x^{(j)})\right)E_{r,\mathrm{}}\right),$$
where $`R=\mathrm{max}(P,Q)`$ and $`S_k(y)`$ are the elementary Schur functions defined by $`_kS_k(y)z^k=\mathrm{exp}(_{k=1}^{\mathrm{}}y_kz^k)`$. In particular if $`1i<jn`$ and $`\alpha =0`$, $`\delta _i\delta _j`$, $`\delta _j\delta _i`$, respectively, then $`\lambda _0=1`$, $`\lambda _{\delta _i\delta _j}=(1)^{nj}`$, $`\lambda _{\delta _j\delta _i}=(1)^{ni+1}`$ and $`(j_{\frac{1}{2}},j_{\frac{3}{2}},\mathrm{})=(\frac{1}{2},\frac{3}{2},\frac{5}{2},\mathrm{})`$, $`=(i\frac{1}{2},\frac{1}{2},\frac{3}{2},\mathrm{},jn+\frac{1}{2},jn+\frac{3}{2}\mathrm{})`$, $`=(j\frac{1}{2},\frac{1}{2},\frac{3}{2},\mathrm{},in+\frac{1}{2},in+\frac{3}{2}\mathrm{})`$, respectively.
We now want to know what happens if we apply $`\omega `$ to such tau-functions. Since $`(\omega (u),\omega (v))=(u,v)`$, $`\omega `$ extends to an automorphism of order 4 on the Clifford algebra. Next notice that $`|0=|H_n^{(0)}F`$. Since $`\omega (H_n^{(0)})=H_n^{(0)}`$, it makes sense to extend $`\omega `$ to $`F`$ by fixing it on $`|0`$ as $`\omega (|0)=|0`$. It is then obvious that there is a one to one correspondence between elements $`wF^{(0)}`$ that satisfy $`\omega (w)=\lambda w`$ for certain $`\lambda `$ and points $`WGr^{(2)}`$. Let $`WGr^{(2)}`$ and
$$w=w_{\frac{1}{2}}w_{\frac{3}{2}}\mathrm{}w_{m+\frac{1}{2}}v_{m\frac{1}{2}}v_{m\frac{3}{2}}\mathrm{}F^{(0)}$$
be the corresponding vector. Since $`w_i,w_j=w_i,v_k=0`$, for all $`0>i,j>m`$ and all $`k<m`$ we obtain that all $`w_i`$ are of the form $`w_i=_{m<j<m}w_{ji}v_j`$. Hence we can find vectors $`w_{\frac{1}{2}},w_{\frac{3}{2}},\mathrm{}w_{m\frac{1}{2}}`$, of the same form such that $`w_i,w_j=v_i,v_j`$ for all $`m<i,j<m`$. Thus $`W=(w_{ij})_{m<i,j<m}`$ is a Symplectic matrix and must have determinant equal to 1. This makes it possible to write $`w`$ in two ways, viz,
$$w=\psi (w_{\frac{1}{2}})\psi (w_{\frac{3}{2}})\mathrm{}\psi (w_{m+\frac{1}{2}})\psi (v_{m\frac{1}{2}}^{})\psi (v_{m\frac{3}{2}}^{})\mathrm{}\psi (v_{\frac{1}{2}}^{})|0$$
(3.21)
and
$`w=\psi (\omega (w_{\frac{1}{2}}`$ $`,w_{\frac{1}{2}}^1w_{\frac{1}{2}}))\psi (\omega (w_{\frac{3}{2}},w_{\frac{3}{2}}^1w_{\frac{3}{2}}))\mathrm{}`$ (3.22)
$`\mathrm{}`$ $`\psi (\omega (w_{m\frac{1}{2}},w_{m+\frac{1}{2}}^1w_{m+\frac{1}{2}}))\psi (v_{m\frac{1}{2}})\psi (v_{m\frac{3}{2}})\mathrm{}\psi (v_{\frac{1}{2}})|0.`$
It is then straightforward to check that $`\omega (w)`$ of the representation (3.21) exactly gives (3.22). This means that $`\omega (w)=w`$ for all elements $`wF^{(0)}`$ corresponding to $`WGr^{(2)}`$.
Next notice that
$$\omega (\psi ^{\pm (j)}(z))=\psi ^{(j)}(z),$$
and hence
$$\omega (\alpha ^{(j)}(z))=\alpha ^{(j)}(z),$$
(3.23)
from which we deduce that
$$\omega (\delta _j)=\delta _j\text{and }\omega (x_k^{(j)})=()^{k+1}x_k^{(j)}.$$
Here we write, as an abuse of notation, $`\omega `$ for $`\sigma \omega \sigma ^1`$. Next we want to calculate what $`\omega `$ does with $`Q_j`$. Notice first,using (3.12), that
$$Q_j^{\pm 1}=\mathrm{exp}(\pm \underset{k<0}{}\frac{1}{k}\alpha _k^{(j)}z^k)\psi ^{\pm (j)}(z)\mathrm{exp}(\pm \underset{k>0}{}\frac{1}{k}\alpha _k^{(j)}z^k)z^{\alpha _0^{(j)}}$$
(3.24)
and that we may replace $`z`$ in this formula (3.24) by $`z`$, since the left-hand side is independent of $`z`$. So,
$`\omega (Q_j^{\pm 1})`$ $`=i\mathrm{exp}({\displaystyle \underset{k<0}{}}{\displaystyle \frac{1}{k}}\alpha _k^{(j)}(z)^k)\psi ^{(j)}(z)\mathrm{exp}({\displaystyle \underset{k>0}{}}{\displaystyle \frac{1}{k}}\alpha _k^{(j)}(z)^k)z^{\pm \alpha _0^{(j)}}`$
$`=iQ_j^1()^{\alpha _0^{(j)}}.`$
Thus we find for the operators $`e^{\pm \delta _j}`$:
$$\omega (e^{\pm \delta _j})=ie^{\delta _j}()^{\delta _j}.$$
So we conclude that
$`\omega `$ $`\left(\tau _0(x_k^{(a)})+{\displaystyle \underset{1i<jn}{}}\left(\tau _{\delta _i\delta _j}(x_k^{(a)})e^{\delta _i\delta _j}+\tau _{\delta _j\delta _i}(x_k^{(a)})e^{\delta _j\delta _i}\right)+\mathrm{}\right)`$ (3.25)
$`=\tau _0(()^{k+1}x_k^{(a)}){\displaystyle \underset{1i<jn}{}}\left(\tau _{\delta _j\delta _i}(()^{k+1}x_k^{(a)})e^{\delta _i\delta _j}+\tau _{\delta _i\delta _j}(()^{k+1}x_k^{(a)})e^{\delta _j\delta _i}\right)+\mathrm{}.`$
Next assume that $`Wgr_0`$. To this subspace corresponds an upto a multiple factor unique a vector $`wF^{(0)}`$. Since $`_{i=1}^nte_{ii}WW`$, we can find special linearly independent vectors $`w_1,w_2,\mathrm{},w_nW`$ such that $`t^{\mathrm{}}w_j=(_{i=1}^nte_{ii})^{\mathrm{}}w_jH_n^{(\mathrm{}n)}`$ for all $`1jn`$ and such that
$$w=w_1w_2\mathrm{}w_ntw_1tw_2\mathrm{}t^\mathrm{}1w_nv_{\mathrm{}n\frac{1}{2}}v_{\mathrm{}n\frac{3}{2}}\mathrm{}.$$
From this presentation of $`w`$ one easily sees that the action of
$$\underset{j=1}{\overset{n}{}}\alpha _k^{(j)}w=0\text{for all }k>0,$$
This leads to
$$\underset{j=1}{\overset{n}{}}\frac{\tau _\alpha (x)}{x_k^{(j)}}=0\text{for all }k>0$$
(3.26)
and hence to the following
###### Proposition 3.2
Tau-functions $`\tau _W(x)=_{\alpha M}\tau _\alpha (x_k^{(a)})e^\alpha `$ corresponding to $`Wgr^{(2)}`$ satisfy the following conditions:
$`(1)`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\tau _\alpha (x)}{x_k^{(j)}}}=0\text{for all }k>0,`$
$`(2)`$ $`\tau _0(()^{k+1}x_k^{(a)})=\tau _0(x_k^{(a)}),`$
$`(3)`$ $`\tau _{\delta _j\delta _i}(()^{k+1}x_k^{(a)})=\tau _{\delta _i\delta _j}(x_k^{(a)})\text{for all }1i,jn,ij.`$
## 4 The KP hierarchy as a dynamical system
It is well known, see e.g. , that $`\tau _W`$ corresponds to a $`WGr_0`$ if and only if $`\tau _W`$ satisfies the the KP hierarchy, i.e., the following equation:
$$\text{Res}_{z=0}\underset{j=1}{\overset{n}{}}\psi ^{+(j)}(z)\tau \psi ^{(j)}(z)\tau =0,\tau F^{(0)}.$$
(4.1)
Using the boson-fermion correspondence we can write this equation as a family of equation on certain $`n\times n`$ wave functions ($`\alpha \text{supp }\tau =\{\alpha M|\tau _\alpha 0\}`$)
$$V^\pm (\alpha ,x,z)=(V_{ij}^\pm (\alpha ,x,z))_{i,j=1}^n,$$
(4.2)
(see for more details) where
$`V_{ij}^\pm (\alpha ,x,z):=\epsilon (\delta _j,\alpha +\delta _i)z^{(\delta _j|\pm \alpha +\alpha _{ij})}`$ (4.3)
$`\times \mathrm{exp}(\pm {\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}x_k^{(j)}z^k)\mathrm{exp}({\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{}{x_k^{(j)}}}{\displaystyle \frac{z^k}{k}})\tau _{\alpha \pm \alpha _{ij}}(x)/\tau _\alpha (x).`$
The equations are:
$$\text{Res}_{z=0}V^+(\alpha ,x,z)V^{}(\beta ,x^{},z)^T=0\text{for all}\alpha ,\beta \text{supp }\tau .$$
(4.4)
Define $`n\times n`$ matrices $`W^{\pm (m)}(\alpha ,x)`$ by the following generating series (cf. (4.3)):
$$\underset{m=0}{\overset{\mathrm{}}{}}W_{ij}^{\pm (m)}(\alpha ,x)(\pm z)^m=\epsilon _{ji}z^{\delta _{ij}1}(\mathrm{exp}\underset{k=1}{\overset{\mathrm{}}{}}\frac{}{x_k^{(j)}}\frac{z^k}{k})\tau _{\alpha \pm \alpha _{ij}}(x))/\tau _\alpha (x).$$
(4.5)
Note that
$$W^{\pm (0)}(\alpha ,x)=I_n,$$
(4.6)
$$W_{ij}^{\pm (1)}(\alpha ,x)=\{\begin{array}{cc}\epsilon _{ji}\tau _{\alpha \pm \alpha _{ij}}/\tau _\alpha \hfill & \text{if}ij\hfill \\ \tau _\alpha ^1\frac{\tau _\alpha }{x_1^{(i)}}\hfill & \text{if}i=j,\hfill \end{array}$$
(4.7)
We see from (4.3) that $`V^\pm (\alpha ,x,z)`$ can be written in the following form:
$$V^\pm (\alpha ,x,z)=\underset{m=0}{\overset{\mathrm{}}{}}W^{\pm (m)}(\alpha ,x)(\pm z)^mR^\pm (\alpha ,\pm z)S^\pm (x,z),$$
(4.8)
where
$`R^\pm (\alpha ,z)`$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}\epsilon (\delta _i,\alpha )E_{ii}(\pm z)^{\pm (\delta _i|\alpha )},`$ (4.9)
$`S^\pm (x,z)`$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}e^{\pm _{j=1}^{\mathrm{}}x_j^{(i)}z^j}E_{ii}.`$
Here $`E_{ij}`$ stands for the $`n\times n`$ matrix whose $`(i,j)`$ entry is $`1`$ and all other entries are zero. Now let $`=_{j=1}^n\frac{}{x_1^{(j)}}`$, then $`V^\pm (\alpha ,x,z)`$ can be written in terms of formal pseudo-differential operators (see for more details). Let
$$P^\pm (\alpha )P^\pm (\alpha ,x,)=I_n+\underset{m=1}{\overset{\mathrm{}}{}}W^{\pm (m)}(\alpha ,x)^m,R^\pm (\alpha )=R^\pm (\alpha ,),$$
(4.10)
then
$$V^\pm (\alpha ,x,z)=P^\pm (\alpha )R^\pm (\alpha )S^\pm (x,z)$$
(4.11)
and one can prove that $`P^{}(\alpha )=P^+(\alpha )^1`$ and the following Lemma:
###### Proposition 4.1
Let $`\alpha ,\beta \text{supp }\tau `$, then $`P^+(\alpha )`$ satisfies the Sato equations:
$$\frac{P^+(\alpha )}{x_k^{(j)}}=(P^+(\alpha )E_{jj}^kP^+(\alpha )^1)_{}P^+(\alpha )$$
(4.12)
and $`P^+(\alpha ),P^+(\beta )`$ satisfy
$$(P^+(\alpha )R^+(\alpha \beta )P^+(\beta )^1)_{}=0\text{for all}\alpha ,\beta \text{supp}\tau .$$
(4.13)
This is another formulation of the $`n`$-component KP hierarchy (see ). Introduce the following formal pseudo-differential operators $`L(\alpha ),C^{(j)}(\alpha )`$:
$`L(\alpha )L(\alpha ,x,)`$ $`=P^+(\alpha )P^+(\alpha )^1,`$ (4.14)
$`C^{(j)}(\alpha )C^{(j)}(\alpha ,x,)`$ $`=P^+(\alpha )E_{jj}P^+(\alpha )^1,`$
then related to the Sato equation is the following linear system
$`L(\alpha )V^+(\alpha ,x,z)`$ $`=zV^+(\alpha ,x,z),`$ (4.15)
$`C^{(i)}(\alpha )V^+(\alpha ,x,z)`$ $`=V^+(\alpha ,x,z)E_{ii},`$
$`{\displaystyle \frac{V^+(\alpha ,x,z)}{x_k^{(i)}}}`$ $`=(L(\alpha )^kC^{(i)}(\alpha ))_+V^+(\alpha ,x,z).`$
To end this section we write down explicitly some of the Sato equations (4.12) on the matrix elements $`W_{ij}^{(s)}`$ of the coefficients $`W^{(s)}(x)`$ of the pseudo-differential operator
$$P=P^+(\alpha )=I_n+\underset{m=1}{\overset{\mathrm{}}{}}W^{(m)}(x)^m.$$
We shall write $`W=W^{(1)}`$ and $`W_{ij}`$ for $`W_{ij}^{(1)}`$ to simplify notation, then the simplest Sato equation is
$$\frac{P}{x_1^{(k)}}=[E_{kk},P]+[W,E_{kk}]P.$$
(4.16)
In particular we have for $`ik`$:
$$\frac{W_{ij}}{x_1^{(k)}}=W_{ik}W_{kj}\delta _{jk}W_{ij}^{(2)}.$$
(4.17)
The equation (4.16) is equivalent to the following equation for $`V=V^+(\alpha )`$:
$$\frac{V}{x_1^{(k)}}=(E_{kk}+[W,E_{kk}])V.$$
(4.18)
## 5 The Darboux-Egoroff system
Define
$$w_{ij}(x)=W_{ij}^{(1)}(0,x),$$
(5.1)
then from the previous section we know that $`w_{ij}(x)`$ satisfies
$$\frac{w_{ij}(x)}{x_1^{(k)}}=w_{ik}(x)w_{kj}(x)ikj.$$
(5.2)
If we moreover assume that the wave function corresponds to a point $`Wgr^{(2)}`$, we also have
$$\underset{k=1}{\overset{n}{}}\frac{w_{ij}(x)}{x_1^{(k)}}=0,$$
(5.3)
and
$$w_{ij}(()^{k+1}x_k^{(a)})=w_{ji}(x_k^{(a)}).$$
(5.4)
This makes it possible to obtain solutions of the Darboux-Egoroff system, viz define
$$\gamma _{ij}(x)=w_{ij}(x)|_{x_{2k}^{(\mathrm{})}=0\text{for all }k1,1\mathrm{}n},$$
(5.5)
then these $`\gamma _{ij}`$ satisfy the equations (1.7). Thus we have obtained the main theorem of this paper.
###### Theorem 5.1
Let $`Wgr^{(2)}=GL_n(L)^{(2)}H_N^{(0)}`$ and let $`\tau (x)=_{\alpha M}\tau _\alpha (x)e^\alpha `$ be the corresponding tau-function. Then the
$$\gamma _{ij}(x)=ϵ_{ji}\left(\frac{\tau _{\delta _i\delta _j}(x)}{\tau _0(x)}\right)_{x_{2k}^{(\mathrm{})}=0\text{for all }k1,1\mathrm{}n}$$
are solutions of the Darboux-Egoroff system (1.7) for $`u_i=x_1^{(i)}`$.
It is obvious that one can construct even more tau-functions that correspond to $`Wgr_0`$ and wich lead to solutions of the Darboux-Egoroff system. Namely, if we take a tau-function which comes from a $`Wgr^{(2)}`$, then we can always multiply it with an element $`e^\beta `$ for $`\beta M`$. The $`\tau _\beta (x)`$ and the $`\tau _{\beta +\delta _i\delta _j}(x)`$ of this new tau-function also lead to solutions of the Darboux-Egoroff system.
It is easy to see from theorem 3.2 that the wave functions satisfy
$$\underset{i=1}{\overset{n}{}}\frac{V^\pm (\alpha ,x,z)}{x_1^{(i)}}=zV^\pm (\alpha ,x,z)$$
(5.6)
This means that we do not really have formal pseudo-differential operators, but rather formal matrix-valued Laurent series in $`z^1`$. The Sato equation takes the following simple form. Let $`P(z)=P^+(\alpha ,x,z)=I+Wz^1=\mathrm{}`$ then
$$\frac{P(z)}{x_k^{(j)}}=(P(z)E_{jj}P(z)^1z^k)_{}P(z).$$
and equation (4.18) turns into
$$\frac{V^+(\alpha ,x,z)}{x_1^{(k)}}=(zE_{kk}+[W,E_{kk}])V^+(\alpha ,x,z).$$
(5.7)
Next let
$$\mathrm{\Phi }^\pm (x,z)=V^\pm (0,x,z)|_{x_{2k}^{(i)}=0\text{for all }k1,1in},$$
(5.8)
then it is straightforward to check that
$$\mathrm{\Phi }^{}(x,z)=\mathrm{\Phi }^+(x,z).$$
Thus
$$\text{Res}_z\mathrm{\Phi }^+(x,z)\mathrm{\Phi }^+(x^{},z)^T=0$$
from which one deduces, when one takes $`x=x^{}`$, that
$$\mathrm{\Phi }^+(x,z)\mathrm{\Phi }^+(x,z)^T=I_n.$$
Let $`\mathrm{\Gamma }(x)=(\gamma _{ij}(x))_{1i,jn}`$, then $`\mathrm{\Phi }(x,z):=\mathrm{\Phi }^+(x,z)`$ satisfies:
$`\mathrm{\Phi }(x,z)\mathrm{\Phi }(x,z)^T`$ $`=I_n,`$ (5.9)
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Phi }(x,z)}{x_1^{(j)}}}`$ $`=z\mathrm{\Phi }(x,z),`$
$`{\displaystyle \frac{\mathrm{\Phi }(x,z)}{x_1^{(k)}}}`$ $`=(zE_{kk}+[\mathrm{\Gamma }(x),E_{kk}])\mathrm{\Phi }(x,z),`$
Following , we want solutions of this system for $`z=0`$. However, just putting $`z=0`$ in (5.9) does not make sense. There is a way to construct such solutions, viz. let $`\tau =g(t)|0`$, with $`g(t)=_iA(i)t^iGL(L)^{(2)}`$, so in particular $`g(t)^T=g(t)^1`$, then
$$\psi (g(t)t^1e_j)\tau 0\text{ and }\psi (g(t)t^ke_j)\tau =0,\text{ for all }1jn,k0.$$
Using the fermionic fields we can rewrite this to
$$\text{Res}_z\underset{i}{}\underset{k=1}{\overset{n}{}}A(i)_{kj}z^{i1}\psi ^{+(k)}(z)\tau 0\text{ and }\text{Res}_z\underset{i}{}\underset{k=1}{\overset{n}{}}A(i)_{kj}z^{i+\mathrm{}}\psi ^{+(k)}(z)\tau =0$$
for all $`1jn`$, $`\mathrm{}0`$ and thus
$$\text{Res}_zz^1V^+(0,x,z)g(z)0\text{ and }\text{Res}_zz^{\mathrm{}}V^+(0,x,z)g(z)=0\text{for all }\mathrm{}0.$$
Now define
$$\mathrm{\Psi }(x,z):=z^1\mathrm{\Phi }(x,z)g(z),$$
then this satisfies
$`\mathrm{\Psi }(x,z)\mathrm{\Psi }(x,z)^T`$ $`=z^2I_n,`$ (5.10)
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Psi }(x,z)}{x_1^{(j)}}}`$ $`=z\mathrm{\Psi }(x,z),`$
$`{\displaystyle \frac{\mathrm{\Psi }(x,z)}{x_1^{(k)}}}`$ $`=(zE_{kk}+[\mathrm{\Gamma }(x),E_{kk}])\mathrm{\Psi }(x,z),`$
$`\text{Res}_zz^{\mathrm{}}\mathrm{\Psi }(x,z)`$ $`=0\text{for all }\mathrm{}>0.`$
We thus get (c.f. ,):
###### Proposition 5.1
Let $`\mathrm{\Psi }(x,z)`$ be constructed as above. Define
$$\psi (x)=(\psi _{ij}(x))_{1i,jn}:=\text{Res}_z\mathrm{\Psi }(x,z),$$
Then these $`\psi _{ij}`$’s satisfy the equations
$$\frac{\psi _{ij}}{x_1^{(k)}}=\gamma _{ik}\psi _{kj},ki,\underset{k=1}{\overset{n}{}}\frac{\psi _{ij}}{x_1^{(k)}}=0.$$
(5.11)
with $`\gamma _{ij}`$ given by (5.5) and the formula’s
$`h_i`$ $`=\psi _{i1},`$ (5.12)
$`\eta _{\alpha \beta }`$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}\psi _{i\alpha }\psi _{i\beta }=\delta _{\alpha \beta },`$
$`{\displaystyle \frac{t^\alpha }{x_1^{(i)}}}`$ $`=\psi _{i1}\psi _{i\alpha },`$
$`c_{\alpha \beta }^\gamma =c_{\alpha \beta \gamma }`$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{\psi _{i\alpha }\psi _{i\beta }\psi _{i\gamma }}{\psi _{i1}}},`$
determine (locally) a semisimple Frobenius manifold on the domain $`x_1^{(i)}x_1^{(j)}`$ and $`\psi _{11}\psi _{21}\mathrm{}\psi _{n1}0`$.
Define
$`\mathrm{\Theta }(x,z)^{}:=`$ $`z^2{\displaystyle \underset{i=1}{\overset{n}{}}}\psi _{i1}E_{ii}\mathrm{\Psi }(x,z),`$
$`\mathrm{\Theta }(x,z)=`$ $`\left(\theta _1(x,z)\theta _2(x,z)\theta _3(x,z)\mathrm{}\theta _n(x,z)\right):=z\left(\psi _{11}\psi _{21}\psi _{31}\mathrm{}\psi _{n1}\right)\mathrm{\Psi }(x,z),`$
then it is straightforward to check that
$`\text{Res}_zz^k\mathrm{\Theta }(x,z)`$ $`=0,\text{Res}_zz^{k1}\mathrm{\Theta }(x,z)^{}=0\text{for all }k0,`$ (5.13)
$`\mathrm{\Theta }(x,z)^{}`$ $`=\left({\displaystyle \frac{\theta _j(x,z)}{x_1^{(i)}}}\right)_{1i,jn},`$
$`\mathrm{\Theta }(x,z)^{}\mathrm{\Theta }(x,z)^{}^T`$ $`=z^2{\displaystyle \underset{i=1}{\overset{n}{}}}h_i^2(x)E_{ii},`$
$`z^1\mathrm{\Theta }(x,z)^{}\mathrm{\Theta }(x,z)^T`$ $`=\left(h_1^2(x)h_2^2(x)h_3^2(x)\mathrm{}h_n^2(x)\right)^T.`$
From which we deduce that the flat coordinates $`t^i`$ are given by
$$\theta _j(x,z)=\delta _{j,1}+t^j(x)z+\underset{k=2}{\overset{\mathrm{}}{}}\theta _j^{(k)}(x)z^k.$$
(5.14)
These $`\theta _j(x,z)`$ are the deformed flat coordinates, see e.g. . So we are in the situation of the paper and we can construct the prepotential $`F(t(x))`$. Using the formula’s (5.10) and (5.11) one calculates that
$`{\displaystyle \frac{^2\mathrm{\Theta }(x,z)}{x_1^{(i)}x_1^{(j)}}}`$ $`=\mathrm{\Gamma }_{ij}^i(x){\displaystyle \frac{\mathrm{\Theta }(x,z)}{x_1^{(i)}}}\mathrm{\Gamma }_{ji}^j(x){\displaystyle \frac{\mathrm{\Theta }(x,z)}{x_1^{(j)}}},ij`$ (5.15)
$`{\displaystyle \frac{^2\mathrm{\Theta }(x,z)}{x_1^{(i)2}}}`$ $`={\displaystyle \underset{j=1}{\overset{n}{}}}\mathrm{\Gamma }_{ii}^j(x){\displaystyle \frac{\mathrm{\Theta }(x,z)}{x_1^{(j)}}}+z{\displaystyle \frac{^2\mathrm{\Theta }(x,z)}{x_1^{(i)}x_1^{(j)}}},`$
where the Christoffel symbols are given by (1.10) and hence that
$$\frac{^2\mathrm{\Theta }(x,z)}{t^kt^{\mathrm{}}}=\underset{m=1}{\overset{n}{}}c_{k\mathrm{}m}(x)z\frac{\mathrm{\Theta }(x,z)}{t^m}$$
(5.16)
Since $`\frac{\mathrm{\Theta }(x,z)}{t^i}`$ is a linear combination of $`\frac{\mathrm{\Theta }(x,z)}{x_1^{(k)}}`$’s,
$$z^1\frac{\mathrm{\Theta }(x,z)}{t^k}\mathrm{\Theta }(x,z)^T$$
is independent of $`z`$, which means that all coefficients, except the constant coefficient, are zero. In particularly the coefficient of $`z^2`$ gives:
$$\theta _m^{(2)}(x)=\frac{\theta _1^{(3)}(x)}{t^m}+\underset{i=1}{\overset{n}{}}t^i\frac{\theta _i^{(2)}(x)}{t^m}.$$
The coefficient of $`z^2`$ of (5.16) leads to
$$\frac{^2\theta _m^{(2)}(x)}{t^kt^{\mathrm{}}}=c_{k\mathrm{}m}(x),$$
hence $`\frac{F(x)}{t^m}=\theta _m^{(2)}(x)`$ and we obtain the Theorem of for the $`GL(L)^{(2)}`$-group orbit:
###### Theorem 5.2
The function $`F(x)=F(t(x),x)`$ defined by
$$F(x)=\frac{1}{2}\theta _1^{(3)}(x)+\frac{1}{2}\underset{i=1}{\overset{n}{}}t^i(x)\theta _i^{(2)}(x)$$
satisfies equation (1.5).
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# On Unitarity Based Relations Between Various Lepton Family Violating Processes
## Basic Considerations
Let us assume that a vector boson $`V_i`$ \[ Here $`V_i`$ could be either a fundamental state, like the $`Z^o`$, or a quark-antiquark bound state like the $`\varphi ,J/\psi `$, or $`\mathrm{{\rm Y}}`$\] couples to $`\mu ^\pm e^{}`$. If it couples also to $`e^+e^{}`$— as all the states above do— then by unitarity its exchange contributes also to $`\mu 3e`$. Let us write the effective coupling between the vector boson $`V_i`$ and $`\mu ^\pm e^{}`$ as
$$_{\mathrm{eff}}=\stackrel{~}{g}_{V\mu e}\overline{\mu }\gamma _\alpha eV^\alpha +\mathrm{h}.\mathrm{c}..$$
(1)
This coupling, through the diagram of Fig. 1, contributes to the $`A(\mu 3e)`$ amplitude a term
$$A(\mu e)=\overline{u}_\mu (p)\gamma ^\alpha u_e(k_3)\overline{v}_e(k_1)\gamma _\alpha u_e(k_2)\frac{\stackrel{~}{g}_{V\mu e}g_{Vee}}{M_V^2s}.$$
(2)
Here $`g_{Vee}`$ is the effective coupling of the vector boson $`V_i`$ to $`e^+e^{}`$, while $`s(k_1+k_2)^2m_\mu ^2`$. <sup>*</sup><sup>*</sup>*There are, of course, also axial vector couplings of the $`Z^o`$ to $`e^+e^{}`$, which contribute to this process. These have not been included in the above, since they do not change our qualitative discussion. These couplings are, however, taken into account in the $`Z`$ bounds given in Eqs. (5) and (9) below. As a first approximation, it is sensible to neglect $`s`$ in comparison with $`M_V^2`$. Then comparing the above contribution to the $`\mu 3e`$ process to that of ordinary muon decay, $`\mu e\nu \overline{\nu }`$, which proceeds via $`W`$ exchange and (almost) identical kinematics, gives the relation:
$$\frac{[\mathrm{\Gamma }(\mu 3e)]_{V\mathrm{exch}.}}{\mathrm{\Gamma }(\mu e\nu \overline{\nu })}\frac{\stackrel{~}{g}_{V\mu e}^2g_{Vee}^2}{M_V^4}/\frac{g_W^4}{M_W^4}.$$
(3)
Since $`\mathrm{\Gamma }(Ve^+e^{})g_{Vee}^2M_V`$ and $`\mathrm{\Gamma }(V\mu e)\stackrel{~}{g}_{V\mu e}^2M_V`$, while $`\mathrm{\Gamma }(We\nu )g_W^2M_W`$, we can rewrite the last expression as
$$[\mathrm{BR}(\mu 3e)]_{V\mathrm{exch}.}\frac{\mathrm{\Gamma }(Ve^+e^{})\mathrm{\Gamma }(V\mu ^\pm e^{})}{\mathrm{\Gamma }^2(We\nu )}\left(\frac{M_W}{M_V}\right)^6.$$
(4)
Using BR$`(\mu 3e)10^{12}`$ and other data pertaining to the $`e^+e^{}`$ widths of the various vector mesons $`V_i`$, we find a set of bounds for the two-body LFV branching ratios of these vector bosons. These bounds are:
$`\mathrm{BR}(Z^o\mu e)`$ $``$ $`5\times 10^{13}`$ (5)
$`\mathrm{BR}(J/\psi \mu e)`$ $``$ $`4\times 10^{13}`$ (6)
$`\mathrm{BR}(\mathrm{{\rm Y}}\mu e)`$ $``$ $`2\times 10^9`$ (7)
$`\mathrm{BR}(\varphi \mu e)`$ $``$ $`4\times 10^{17}.`$ (8)
Likewise, the generic upper bounds on LFV tau decays BR$`(\tau \mathrm{}\mathrm{}^{}\overline{\mathrm{}}^{})10^6`$ yields
$`\mathrm{BR}(Z^o\tau \overline{\mathrm{}})`$ $``$ $`3\times 10^6`$ (9)
$`\mathrm{BR}(J/\psi \tau \overline{\mathrm{}})`$ $``$ $`6\times 10^7`$ (10)
$`\mathrm{BR}(\mathrm{{\rm Y}}\tau \overline{\mathrm{}})`$ $``$ $`10^2,`$ (11)
with $`\mathrm{}/\mathrm{}^{}=e/\mu `$. Except for (11), these inferred bounds are unlikely to be improved by future experimental data on two-body decays.
One can use similar considerations to obtain bounds on the LFV decays of pseudoscalar states. For these purposes, one considers instead of the $`\mu 3e`$ process the $`\mu ^{}e^{}\gamma \gamma `$ decay, which has a LFV bound BR$`(\mu e\gamma \gamma )10^{10}`$. For this latter process the LFV couplings of the $`\pi ^o/\eta ^o`$ contribute, due to the exchange of these particles in the s-channel. We can again utilize this fact to infer upper bounds on $`\pi ^o/\eta ^o\mu ^\pm e^{}`$. The $`\pi ^o/\eta ^o\gamma \gamma `$ vertex, because of gauge invariance, involves two derivatives:
$$_{\mathrm{eff}}=\frac{\varphi }{f_\varphi }F^{\mu \nu }\stackrel{~}{F}_{\mu \nu },$$
(12)
where $`\varphi =\pi ^o,\eta ^o`$. This derivative coupling, in contrast to the $`V\mu e`$ non-derivative coupling encountered earlier, kinematically suppresses the off-shell $`\pi ^o/\eta ^0\gamma \gamma `$ contribution at $`s=(k_1+k_2)^2m_\mu ^2`$ relative to what it would be for on-shell $`\pi ^o/\eta ^o`$ decay. Consequently the analog to Eq. (4) for the present case,
$$[\mathrm{BR}(\mu e\gamma \gamma )]_{\pi ^o/\eta ^o\mathrm{exch}.}\frac{\mathrm{\Gamma }(\pi ^o/\eta ^o\gamma \gamma )\mathrm{\Gamma }(\pi ^o/\eta ^o\mu ^\pm e^{})}{\mathrm{\Gamma }^2(We\nu )}\left(\frac{M_W}{m_{\pi ^/\eta }}\right)^6\left(\frac{s_{\gamma \gamma }}{m_{\pi /\eta }^2}\right)^2,$$
(13)
contains an extra factor
$$\frac{s_{\gamma \gamma }^2}{m_{\pi /\eta }^4}\left(\frac{m_\mu }{2m_{\pi /\eta }}\right)^4,$$
(14)
which tends to weaken the bounds one can derive. One finds, for pseudoscalar LFV decays the bounds:
$`\mathrm{BR}(\eta \mu e)`$ $``$ $`10^8`$ (15)
$`\mathrm{BR}(\pi ^o\mu e)`$ $``$ $`10^{10}.`$ (16)
In the discussion above we have obtained the quoted bounds purely by concentrating on the contribution of the exchanged state in question to the LFV process. One can imagine, however, additional LFV contributions. For example, for the $`\mu ^{}e^{}\gamma \gamma `$ decay, in addition to $`\pi ^o/\eta ^0`$ exchange in the $`s`$ channel, we have also the contribution of electron exchange in the $`t`$ and $`u`$ channels (see Fig. 2). In this case, however, the stringent bound on the $`\mu e\gamma `$ vertex coming from experiment \[BR ($`\mu e\gamma )5\times 10^{11}`$\] strongly suppresses these additional diagrams and causes negligible modifications to the bounds (15), (16). Even in the absence of a strong bound on the $`\mu e\gamma `$ coupling, we would like to note that cancellations between $`s`$ and $`t`$ channel contributions are in general expected to be at best rather partial. Unless all particles, both external and exchanged, are spinless any specific $`s`$ channel amplitude will have different $`\mathrm{cos}\theta _s`$ (or $`\mathrm{cos}\theta _t`$) dependence, and will contribute to different combinations of helicity amplitudes than the $`t`$ channel exchange contributions.
By the same token, it is clear that cancellations among different angular momentum states exchanged in the $`s`$ channel are also impossible. Indeed, for example, the total decay rate for $`\mu 3e`$ can be expressed as
$`\mathrm{\Gamma }(\mu 3e)`$ $`=`$ $`{\displaystyle _0^{\mu ^2}}𝑑s{\displaystyle \underset{\alpha \beta ,\gamma \delta }{}}(s4m_e^2)^{1/2}\left[{\displaystyle \frac{\mathrm{\Delta }(m_\mu ^2,s,m_e^2)}{s}}\right]^{1/2}`$ (18)
$`\times {\displaystyle \underset{J}{}}(2J+1)|A_{\alpha \beta ,\gamma \delta }^J(s)|^2`$
with $`\mathrm{\Delta }`$ the triangular function expressing the initial C.M. momentum in the $`s`$ channel, which here is that of $`\mu \overline{e}_1`$ or, equivalently, $`\overline{e}_3e_2`$. The $`A_{\alpha \beta ,\gamma \delta }^J(s)`$ are the partial waves in the Jacob-Wick expansion of the various $`s`$ channel helicity amplitudes. Note that for the $`\mu 3e`$ case, adding the $`t`$ channel amplitude amounts to enforcing the (anti) symmetrization between the $`e_3`$ and $`e_2`$ fermions. Since Fermi statistics does not preclude the vectorial coupling contributing to $`\mu 3e`$ considered here \[c.f. Eq. (2)\], no cancellation of $`s`$ and $`t`$ contributions should arise as well.
## Possible Limitations on the Derived LFV Bounds
Although we have called the bounds we obtained above unitarity bounds, in the strict sense the inferred bounds are not true unitarity bounds–as would be the case if the exchanged particle(s) were on mass shell. To illustrate this point, let us recall a well known example of a true unitarity bound arising in rare Kaon decays. This is the lower bound for the BR$`(K_L\mu ^+\mu ^{})`$ derived from the measured branching ratio of $`K_L\gamma \gamma `$. The $`K_L\gamma \gamma `$ process, with an on-shell $`\gamma \gamma `$ intermediate state, contributes to Im $`A(K_L\mu ^+\mu ^{})`$ via unitarity since Im $`A(K_L\mu ^+\mu ^{})`$ $`A(K_L\gamma \gamma )A(\gamma \gamma \mu ^+\mu ^{})`$. This contribution provides a strict lower bound to the BR$`(K_L\mu ^+\mu ^{})`$, so the apparent violation of this bound in early $`K_L\mu ^+\mu ^{}`$ data was a source for much concern. Modern day data, as expected, agrees with this bound.
In the present context, an example of a “pure” unitarity bound for a LFV process is provided by the “$`\tau `$ analog” of the $`\varphi \mu e`$ process. Because the decay $`\varphi \tau \overline{\mathrm{}}`$ is kinematically forbidden, what one should consider instead is $`\tau \varphi \mu `$ (or $`\tau \varphi e`$). The “on-shell” $`\varphi `$ emitted in this putative process propagates over a long distance, of order $`1/\mathrm{\Gamma }_\varphi 30`$ fm, before decaying into $`K\overline{K},\mu ^+\mu ^{},e^+e^{}`$ in a manner which is completely independent of its production. This will generate a distinct narrow contribution to the corresponding three body processes $`\tau K\overline{K}\mu `$, $`\tau \mu ^+\mu ^{}\mu `$, $`\tau e^+e^{}\mu `$ ,contributing to the imaginary part of these amplitudes at $`s=m_\varphi ^2`$. Hence, for example, there is an attendant lower bound on BR$`[\tau \mathrm{}\mathrm{}^{}\overline{\mathrm{}}^{}]`$ which is simply BR$`(\tau \varphi \overline{\mathrm{}})\mathrm{BR}(\varphi \mathrm{}^{}\overline{\mathrm{}}^{})3\times 10^4\times \mathrm{BR}(\tau \varphi \overline{\mathrm{}})`$. The resulting rigorous upper bound one obtains,
$$\mathrm{BR}(\tau \varphi \overline{\mathrm{}})3\times 10^3,$$
(19)
unfortunately happens to be rather weak.
All the vector (or pseudoscalars) used as intermediaries in deriving the bounds in Eqs (5)-(11) and Eqs. (15) and (16) are not on-shell. Thus we must entertain the possibility that their contribution to the three-body decays considered are reduced. This could void, or at least weaken the various strong bounds obtained above. In the rest of this note we will focus on possible mechanisms for such a reduction.
### Kinematical Suppression of the LFV Bounds
The size of the boson exchange contribution to the three-body decay amplitude can be reduced if there are kinematical suppressions. These arise when the effective boson couplings are not minimal, involving derivatives (or momentum factors). We already encountered one such case above, when we discussed the $`\pi ^o/\eta ^o`$ contribution to $`\mu e\gamma \gamma `$. We want to discuss here whether such kinematical suppressions may not also affect the vector exchange contributions.
It is clearly possible to imagine that the LFV $`V_i\mu ^\pm e^{}`$ vertex, instead of having the form of Eq. (1), involves an anomalous magnetic moment coupling:
$$_{\mathrm{eff}}^{\mathrm{Magnetic}}=\frac{1}{M_V}\overline{\mu }\sigma _{\alpha \beta }e(^\alpha V^\beta ^\beta V^\alpha )+\mathrm{h}.\mathrm{c}..$$
(20)
In this case the contribution of the virtual $`V_i`$ to $`\mu 3e`$ is reduced by
$$\frac{q^2}{M_V^2}\frac{s}{M_V^2}\frac{m_\mu ^2}{2M_V^2}=\left\{\begin{array}{ccc}3\times 10^3\hfill & & V=\varphi \hfill \\ 3\times 10^4\hfill & & V=J/\psi \hfill \\ 3\times 10^5\hfill & & V=\mathrm{{\rm Y}}\hfill \\ 3\times 10^7\hfill & & V=Z^o\hfill \end{array}\right\}$$
(21)
This would considerably weaken the bounds in Eqs. (6)-(8) and reduce the bound on $`Z`$ decay to only BR$`(Z^oe\mu )1.5\times 10^6`$.
It does not seem likely to us, however, that the strong suppression factors appearing in Eq. (20) obtain in practice. Indeed, various model calculation involving mixing among heavy neutrinos lead to $`Z^oe\mu `$ and effective $`c\overline{c}\mu e`$ flavor violating vertices, which involve non anomalous terms— terms like the $`Z\mu e`$ coupling of Eq. (1) and vectorial couplings like $`\overline{c}\gamma _\alpha c\overline{\mu }\gamma ^\alpha e`$. Hence we believe that the kinematic suppression given by Eq. (20) probably should not be included in our bounds.
### Dynamical Suppression of the LFV Bounds
There is another possible source of suppression which needs to be considered. This is connected to possible ”form factor” effects due to the dynamics which would, for example, reduce the contribution of the various $`V_i`$ states to $`\mu 3e`$ compared to the naive expectations. However, the effect of form factors should be minimal or controllable if the LFV is induced by physics at scales much higher than the EW scale or the $`Z`$ mass. The effects of dynamics are nicely illustrated in a recent paper by Ilana, Jack and Riemann . These authors find, in fact, an apparent mild enhancement when the $`Z^o\mu e`$ process is induced by relatively light ($`m_{\nu _i}45`$ GeV) neutrinos. Indeed, in this case the on-shell contribution of $`\nu \overline{\nu }`$ loops enhances the $`Z`$ decay rates relative to the $`s0`$ contribution by factors of 10-100. However, such light active neutrinos would contribute to the $`Z^o`$ width and are hence ruled out. Thus such an ehancement is not physically expected. The BR considered in Ref. for light neutrinos– i.e. neutrinos with masses in the eV range, as inferred from the SuperKamiokande data– are very much below our bounds
In terms of the dispersive approach adopted here, such a ”form factor” suppression would result from cancellations in the corresponding partial wave amplitudes. Consider, for example, the $`A^{J=1}`$ partial wave amplitude for the $`\mu 3e`$ process:
$$A_{\alpha \beta ,\gamma \delta }^{J=1}(s)=\underset{i}{}\frac{g_{V_iee}\stackrel{~}{g}_{V_i\mu e}}{M_i^2s}+\frac{ds^{}\rho _{\mathrm{LFV}}^{J=1}(s^{})}{s^{}s}.$$
(22)
To get a ”form factor” suppression, there must be a cancellation between the contributions of the various $`(V_i)`$ particles among themselves, or between these contributions and those of the continuum. Let us examine these possibilities.
For the case of quarkonium intermediate states, besides the lowest energy bound state there are towers of states of the same spin and parity. Thus, for example, in Eq.(21) besides the contribution of the $`J/\psi `$ one should also take into acount the exchange of the $`\psi ^{},\psi ^{\prime \prime },\mathrm{},\psi ^{(n)}`$ charmonium bound resonances. Is it possible that these additional contributions largely cancel the $`J/\mathrm{\Psi }`$ term in Eq. (21)? This is unlikely for the following reason. To get $`J/\psi `$ exchange to contribute to $`\mu 3e`$ in the first place, one needs to assume that the LFV physics at a high scale induces an effective four-Fermi coupling of the form:
$$_{\mathrm{eff}}=\stackrel{~}{G}_{c\overline{c}\mu e}\overline{c}\gamma _\alpha c\overline{\mu }\gamma ^\alpha e.$$
(23)
Such a coupling underlies all the other charmonium contributions. In fact, quark-hadron duality identifies the $`J/\psi ,\psi ^{},\mathrm{},\psi ^{(n)}`$ contributions to $`\mu 3e`$ as arising from specific portions of the $`s^{}M_{\overline{c}c}`$ integration region where due to non-perturbative QCD effects $`1^{}`$ $`c\overline{c}`$ bound states dominate, as shown schematically in Fig. 3. Both the $`g_{\psi ^{(n)}ee}`$ and $`\stackrel{~}{g}_{\psi ^{(n)}\mu e}`$ couplings appearing in Eq. (21) are proportional to the wave function of $`\psi ^{(n)}`$ at the origin, $`𝚿^{(𝐧)}(\mathrm{𝟎})`$. Thus all the terms in the sum share a common sign— fixed by the sign of $`\stackrel{~}{G}_{c\overline{c}\mu e}Q_c`$, with $`Q_c=2/3`$ being the charge of the charm quark— and cancellation cannot occur. Similar arguments apply against possible cancellations among the various states in the $`\mathrm{{\rm Y}}`$ sector.
The above discussion still leaves open the possibility of cancellations in the partial wave amplitude between different quarks-antiquark contributions ($`c\overline{c},b\overline{b},s\overline{s},\mathrm{}`$) or, equivalently, between the various resonant states ($`J/\psi ,\mathrm{{\rm Y}},\varphi ,\mathrm{}`$. While possible this seems highly unlikely. For example, even if all the effective couplings $`\stackrel{~}{G}_{q_i\overline{q}_i\mu e}`$ were equal due to some universality, and the bubble kinematics were identical, the net contribution would not vanish since the total contribution would be proportional to $`Q_{q_i}0`$. Furthermore, for the case of light quarks such cancellations cannot work even in principle. The $`s`$ dependence for $`0sm_\mu ^2`$ neglected above implies, for example, that a $`\omega /\rho `$ and $`\varphi `$ contribution to the total decay rate cancel only at the level of $`(m_\mu ^2/3)(m_\varphi ^2m_{\rho /\omega }^2)/(m_\varphi ^2)^210^3`$.
## Concluding Remarks
In general, lepton flavor violating processes have been analyzed within a specific theoretical framework. In this context, the restrictive role played by the low energy bounds ($`\mu 3e\mu e`$ conversion, etc.) has been noted by many authors. In this note instead we tried to present in a, relatively model-independent manner, the connections which unitarity implies between some two-body and three-body LFV decays. We have illustrated these connections by focusing on a few processes. Clearly, many other bounds can be obtained. Indeed, since the Particle Data Group lists altogether about one hundred LFV processes, many additional results can come from a more comprehensive analysis.
We have noted that the bounds that we derived can be avoided if one can kinematically suppress the small $`s`$ contributions (e.g. by having a purely anomalous magnetic $`Z^o\mu e`$ coupling), or as a result of some (rather unlikely) cancellations. Because we cannot rule out these possibilities with absolute certainty, we hope that our discussion will not dissuade future efforts to improve the bounds on LFV decays of the $`Z,J/\psi ,\mathrm{}`$. Such decays would not only signal new LFV physics but, because of our considerations, this physics must also naturally give cancellations among terms so as to lead to a small $`\mu 3e`$ branching ratio. Thus searching for $`V_i\mu ^\pm e^{}`$ decays at levels considerably higher than our suggested bounds remains a worthwhile experimental challenge.
## Acknowledgements
Both S. N. and X. M. Z. would like to acknowledge the hospitality of the department of Physics and Astronomy at UCLA, where this work was initiated. S. N. would like to acknowledge the support of USA-Israel Binational and Israeli Academy Grants. X. M. Z. thanks Z. H. Lin for discussions. The work of R. D. P. was supported in part by the Department of Energy under contract No. DE-FG03-91ER40662, Task C.
## References
## Figure Captions
Fig. 1. A vector exchange diagram contributing to $`\mu 3e`$.
Fig. 2. The $`\pi /\eta ^o`$ ($`s`$ channel) and $`e`$ ($`t`$ and $`u`$ channel) exchange contributing to $`\mu e\gamma \gamma `$.
Fig. 3. The $`\overline{c}c`$ bubble and its equivalent description in terms of $`\psi ,\psi ^{},\mathrm{},\psi ^n`$ exchanges.
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# The variational principle for a class of asymptotically abelian C∗-algebras
## 1 Introduction
The variational principle has over the years attracted much attention both in classical ergodic theory, see e.g. \[W\], and in the C-algebra setting of quantum statistical mechanics, see e.g. \[BR\]. In the years around 1970 there was much progress in the study of spin lattice systems. In that case one is given for each point $`x^\nu `$ ($`\nu `$) the algebra of all linear operators $`B()_x`$ on a finite dimensional Hilbert space $``$, and if $`\mathrm{\Lambda }^\nu `$ one defines the C-algebra $`A(\mathrm{\Lambda })=_{x\mathrm{\Lambda }}B()_x`$. One then considers the UHF-algebra $`A=\overline{_{\mathrm{\Lambda }^\nu }A(\mathrm{\Lambda })}`$ with space translations $`\alpha `$, and studies mean entropy $`s(\varphi )`$ of invariant states and the corresponding variational principle
$$P(H)=\underset{\varphi }{sup}(s(\varphi )\varphi (H))$$
(1.1)
together with the KMS-states defined by a natural derivation associated with a self-adjoint operator $`H`$, see \[BR, Chapter 6\].
With the development of dynamical entropy of automorphisms of C-algebras \[CS, CNT, V\] it was natural to replace the mean entropy $`s(\varphi )`$ by dynamical entropy $`h_\varphi (\alpha )`$. This was done by Narnhofer \[N\], who considered KMS-properties of the states on which the quantity $`h_\varphi (\alpha )\varphi (H)`$ attains its maximal value. Then Moriya \[M\] showed that one can replace $`s(\varphi )`$ by the CNT-entropy $`h_\varphi (\alpha )`$ in (1.1) and get the same result, i.e.
$$P(H)=\underset{\varphi }{sup}(h_\varphi (\alpha )\varphi (H)).$$
(1.2)
If one wants to study the variational principle and equilibrium states for more general C-dynamical systems, the mean entropy is usually not well defined, and it is necessary to consider dynamical entropy. However, in order to define time translations and extend the theory of spin lattice systems rather strong assumptions of asymptotic abelianness are needed.
In the present paper we shall study a restricted class of asymptotically abelian systems $`(A,\alpha )`$, namely we shall assume that $`A`$ is a unital separable C-algebra, and $`(A,\alpha )`$ is asymptotically abelian with locality, i.e. there exists a dense $`\alpha `$-invariant $``$-subalgebra $`𝒜`$ of $`A`$ such that for all pairs $`a,b𝒜`$ the C-algebra they generate is finite dimensional, and there is $`p=p(a,b)`$ such that $`[\alpha ^j(a),b]=0`$ for $`|j|p`$. In particular, $`A`$ is an AF-algebra. Examples of such systems are described in \[S\]. They are all different shifts, on infinite tensor products of the same AF-algebra with itself, on the Temperley-Lieb algebras, on towers of relative commutants, on binary shifts algebras defined by finite subsets of the natural numbers.
In Section 2 we define the pressure $`P_\alpha (H)`$ of $`\alpha `$ with respect to a self-adjoint operator $`HA`$. This can be done in any unital C-dynamical system $`(A,\alpha )`$ with $`A`$ a nuclear C-algebra, and follows closely Voiculescu’s definition of topological entropy $`ht(\alpha )`$ in \[V\]. The main difference is that he considered $`\mathrm{rank}B`$ of a finite dimensional C-algebra $`B`$, while we consider quantities of the form $`\mathrm{Tr}_B(e^K)`$ for $`K`$ self-adjoint, where $`\mathrm{Tr}_B`$ is the canonical trace on $`B`$ (so in particular we get $`\mathrm{rank}B`$ when $`K=0`$). We can then show the analogues of several of the classical results on pressure as presented in \[W\].
If $`(A,\alpha )`$ is asymptotically abelian with locality we show the variational principle (1.2) in Section 3. Furthermore, if we assume $`ht(\alpha )<\mathrm{}`$, $`H𝒜`$ and $`\varphi `$ is a $`\beta `$-equilibrium state at $`H`$, i.e. $`\varphi `$ is $`\alpha `$-invariant and $`P_\alpha (\beta H)=h_\varphi (\alpha )\beta \varphi (H)`$, then we show in Section 4 via a proof modelled on the corresponding proof based on (1.1) in \[BR\] for spin lattice systems, that $`\varphi `$ is a $`\beta `$-KMS state with respect to the one-parameter group defined by the derivation
$$\delta _H(x)=\underset{j}{}[\alpha ^j(H),x],x𝒜.$$
In particular, when $`H=0`$, so $`P_\alpha (0)=ht(\alpha )`$, we obtain
$$ht(\alpha )=\underset{\varphi }{sup}h_\varphi (\alpha ),$$
(1.3)
and if $`h_\varphi (\alpha )=ht(\alpha )`$ then $`\varphi `$ is a trace.
Equation (1.3) is false in general. Indeed in \[NST\] there is exhibited a (non asymptotically abelian) automorphism $`\alpha `$ of the CAR-algebra for which the trace $`\tau `$ is the unique invariant state, $`h_\tau (\alpha )=0`$, while $`ht(\alpha )\frac{1}{2}\mathrm{log}2`$. Furthermore, our assumption of locality is essential to conclude that $`\varphi `$ is a trace when $`h_\varphi (\alpha )=ht(\alpha )<\mathrm{}`$, even for asymptotically abelian systems. In Example 5.7 we show that there is a Bogoliubov automorphism $`\alpha `$ of the even CAR-algebra which is asymptotically abelian, $`ht(\alpha )=0`$, while there are an infinite number of non-tracial $`\alpha `$-invariant states.
In Section 5 we consider some other examples and special cases. First we apply our results to C-algebras associated with certain AF-groupoids arising naturally from expansive homeomorphisms of zero-dimensional compact spaces. We show that if $`H`$ lies in the diagonal then there is a one-to-one correspondence between equilibrium states on the algebra and equilibrium measures in the classical sense. We also consider unique ergodicity for non-abelian systems. If $`(A,\alpha )`$ is asymptotically abelian with locality, unique ergodicity turns out to be of marginal interest. Indeed, the unique invariant state $`\tau `$ is a trace, and the image of $`A`$ in the GNS-representation of $`\tau `$ is abelian.
Acknowledgement. The authors are indebted to A. Connes for suggesting to us to study the variational principle and equilibrium states in the setting of asymptotically abelian C-algebras.
## 2 Pressure
In order to define pressure of a $`C^{}`$-dynamical system we follow the setup of Voiculescu \[V\] for his definition of topological entropy.
Let $`A`$ be a nuclear $`C^{}`$-algebra with unit and $`\alpha `$ an automorphism. Let $`\mathrm{CPA}(A)`$ denote the set of triples $`(\rho ,\psi ,B)`$, where $`B`$ is a finite dimensional $`C^{}`$-algebra, and $`\rho :AB`$, $`\psi :BA`$ are unital completely positive maps. $`𝒫_f(A)`$ denotes the family of finite subsets of $`A`$. For $`\delta >0`$, $`\omega 𝒫_f(A)`$ and $`HA_{sa}`$, put
$$P(H,\omega ;\delta )=inf\{\mathrm{log}\mathrm{Tr}_B(e^{\rho (H)})|(\rho ,\psi ,B)\mathrm{CPA}(A),(\psi \rho )(x)x<\delta x\omega \},$$
where $`\mathrm{Tr}_B`$ is the canonical trace on $`B`$, i.e. the trace which takes the value $`1`$ on each minimal projection. Then set
$`P_\alpha (H,\omega ;\delta )`$ $`=`$ $`\underset{n\mathrm{}}{lim\; sup}{\displaystyle \frac{1}{n}}P({\displaystyle \underset{j=0}{\overset{n1}{}}}\alpha ^j(H),{\displaystyle \underset{j=0}{\overset{n1}{}}}\alpha ^j(\omega );\delta ),`$
$`P_\alpha (H,\omega )`$ $`=`$ $`\underset{\delta >0}{sup}P_\alpha (H,\omega ;\delta ).`$
###### Definition 2.1
The pressure of $`\alpha `$ at $`H`$ is
$$P_\alpha (H)=\underset{\omega 𝒫_f(A)}{sup}P_\alpha (H,\omega ).$$
We have chosen the minus sign $`e^{\rho (H)}`$ in the definition of $`P(H,\omega ;\delta )`$ because of its use in physical applications, see \[BR\], rather than the plus sign used in ergodic theory, see \[W\].
It is easy to see that if $`\omega _1\omega _2\mathrm{}`$ is an increasing sequence of finite subsets of $`A`$ such that the linear span of $`_n\omega _n`$ is dense in $`A`$, then $`P_\alpha (H)=lim_nP_\alpha (H,\omega _n)`$. If $`H=0`$ the pressure coincides with the topological entropy of Voiculescu \[V\].
Recall that $`h_\varphi (\alpha )`$ denotes the CNT-entropy of $`\alpha `$ with respect to an $`\alpha `$-invariant state $`\varphi `$ of $`A`$.
###### Proposition 2.2
For any $`\alpha `$-invariant state $`\varphi `$ of $`A`$ we have $`P_\alpha (H)h_\varphi (\alpha )\varphi (H)`$.
Proof. The proof is a rewording of the proof of \[V, Proposition 4.6\]. Let $`N`$ be a finite dimensional C-algebra, and $`\gamma :NA`$ a unital completely positive map. Let $`\omega 𝒫_f(A)`$ be such that $`H\omega `$ and $`\gamma (\{xN|x1\})`$ is contained in the convex hull of $`\omega `$. Then if $`(\rho ,\psi ,B)\mathrm{CPA}(A)`$ and
$$(\psi \rho )(a)a<\delta \text{for}a_{j=0}^{n1}\alpha ^j(\omega ),$$
we obtain by \[CNT, Proposition IV.3\]
$$|H_\varphi (\{\alpha ^j\gamma \}_{0jn1})H_\varphi (\{\psi \rho \alpha ^j\gamma \}_{0jn1})|<n\epsilon ,$$
where $`\epsilon =\epsilon (\delta )0`$ as $`\delta 0`$. If $`KB_{sa}`$ and $`\theta `$ is a state on $`B`$ then $`\mathrm{log}\mathrm{Tr}_B(e^K)S(\theta )\theta (K)`$, hence
$`H_\varphi (\{\psi \rho \alpha ^j\gamma \}_{0jn1})`$ $``$ $`H_\varphi (\psi )S(\varphi \psi )`$
$``$ $`\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(H)\right)}\right)+(\varphi \psi )\left(\rho \left({\displaystyle _{j=0}^{n1}}\alpha ^j(H)\right)\right)`$
$``$ $`\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(H)\right)}\right)+n\varphi (H)+n\delta .`$
Thus
$$\frac{1}{n}H_\varphi (\{\alpha ^j\gamma \}_{0jn1})\frac{1}{n}P(\underset{j=0}{\overset{n1}{}}\alpha ^j(H),\underset{j=0}{\overset{n1}{}}\alpha ^j(\omega );\delta )+\varphi (H)+\delta +\epsilon .$$
It follows that $`h_\varphi (\gamma ;\alpha )P_\alpha (H,\omega )+\varphi (H)`$, hence $`h_\varphi (\alpha )P_\alpha (H)+\varphi (H)`$.
###### Remark 2.3
If $`A`$ is abelian, $`A=C(X)`$, and $`P_\alpha ^{cl}(H)`$ denotes the pressure as defined in \[W\], then $`P_\alpha (H)=P_\alpha ^{cl}(H)`$. The inequality ’$``$’ can be proved just the same as in \[V, Proposition 4.8\]. The converse inequality follows from Proposition 2.2 and the classical variational principle. In the AF-case however, i.e. when $`X`$ is zero-dimensional, it is easy to give a direct proof. Indeed, if in the proof of Proposition 2.2 $`N`$ was the subalgebra of $`A`$ corresponding to a clopen partition $`P`$ of $`X`$, then we could conclude that
$$\frac{1}{n}H_\varphi (\{\alpha ^j(N)\}_{0jn1})\frac{1}{n}P(\underset{j=0}{\overset{n1}{}}\alpha ^j(H),\underset{j=0}{\overset{n1}{}}\alpha ^j(\omega );\delta )+\frac{1}{n}\varphi \left(\underset{j=0}{\overset{n1}{}}\alpha ^j(H)\right)+\delta +\epsilon $$
(2.1)
for any (not necessarily $`\alpha `$-invariant) state $`\varphi `$ of $`A`$, with $`\epsilon `$ independent of $`\varphi `$. Let $`T`$ be the homeomorphism corresponding to $`\alpha `$, so that $`\alpha (f)=fT`$. Suppose the points $`x_1,\mathrm{},x_m`$ lie in different elements of the partition $`_{j=0}^{n1}T^jP`$. Then inequality (2.1) for the measure
$$\varphi =\left(\underset{i}{}e^{(S_nH)(x_i)}\right)^1\underset{i}{}e^{(S_nH)(x_i)}\delta _{x_i},$$
where $`(S_nH)(x)=_{j=0}^{n1}H(T^jx)`$, means that
$$\frac{1}{n}\mathrm{log}\underset{i}{}e^{(S_nH)(x_i)}\frac{1}{n}P(\underset{j=0}{\overset{n1}{}}\alpha ^j(H),\underset{j=0}{\overset{n1}{}}\alpha ^j(\omega );\delta )+\delta +\epsilon .$$
Recalling the definition of pressure \[W, Definition 9.7\], we see that $`P_\alpha ^{cl}(H)P_\alpha (H)`$.
We list some properties of the function $`HP_\alpha (H)`$ on $`A_{sa}`$.
###### Proposition 2.4
The following properties are satisfied by $`P_\alpha `$ for $`H,KA_{sa}`$.
* (i) If $`HK`$ then $`P_\alpha (H)P_\alpha (K)`$.
* (ii) $`P_\alpha (H+c1)=P_\alpha (H)c`$, $`c`$.
* (iii) $`P_\alpha (H)`$ is either infinite for all $`H`$ or is finite valued.
* (iv) If $`P_\alpha `$ is finite valued then $`|P_\alpha (H)P_\alpha (K)|HK`$.
* (v) For $`k`$, $`P_{\alpha ^k}(_{j=0}^{k1}\alpha ^j(H))=kP_\alpha (H)`$.
* (vi) $`P_\alpha (H+\alpha (K)K)=P_\alpha (H)`$.
Proof. (i) Given $`HK`$ take $`\omega 𝒫_f(A)`$. If $`(\rho ,\psi ,B)\mathrm{CPA}(A)`$ we have
$$\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(H)\right)}\right)\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(K)\right)}\right),$$
see e.g. \[OP, Corollary 3.15\]. Thus (i) follows.
(ii) As in (i) we have
$$\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(H+c1)\right)}\right)=\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(H)\right)}\right)nc,$$
and (ii) follows.
(iii) By (i) and (ii) we have
$$P_\alpha (H)P_\alpha (H)=P_\alpha (0)H=ht(\alpha )H,$$
and similarly $`P_\alpha (H)ht(\alpha )+H`$. Thus (iii) follows.
(iv) For any $`(\rho ,\psi ,B)\mathrm{CPA}(A)`$ we have by the Peierls-Bogoliubov inequality \[OP, Corollary 3.15\]
$$\left|\frac{1}{n}\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(H)\right)}\right)\frac{1}{n}\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(K)\right)}\right)\right|HK.$$
Thus
$$P_\alpha (H,\omega ;\delta )HKP_\alpha (K,\omega ;\delta )P_\alpha (H,\omega ;\delta )+HK$$
for any $`\omega 𝒫_f(A)`$. Thus (iv) follows.
(v) Let $`(\rho ,\psi ,B)\mathrm{CPA}(A)`$ and $`\omega 𝒫_f(A)`$. Given $`n`$ choose $`m`$ such that $`mkn<(m+1)k`$. Set $`H_k=_{j=0}^{k1}\alpha ^j(H)`$ and $`\omega _k=_{j=0}^{k1}\alpha ^j(\omega )`$. Then
$`\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(H)\right)}\right)`$ $``$ $`\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{mk1}\alpha ^j(H)\right)kH}\right)`$
$`=`$ $`\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{m1}\alpha ^{jk}(H_k)\right)}\right)kH.`$
Similarly
$$\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(H)\right)}\right)\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^m\alpha ^{jk}(H_k)\right)}\right)+kH.$$
Since $`_{j=0}^{m1}\alpha ^{jk}(\omega _k)_{j=0}^{n1}\alpha ^j(\omega )_{j=0}^m\alpha ^{jk}(\omega _k)`$, it follows that
$$P_\alpha (H,\omega ;\delta )=\frac{1}{k}P_{\alpha ^k}(\underset{j=0}{\overset{k1}{}}\alpha ^j(H),\underset{j=0}{\overset{k1}{}}\alpha ^j(\omega );\delta ),$$
and hence $`P_\alpha (H)=\frac{1}{k}P_{\alpha ^k}(_{j=0}^{k1}\alpha ^j(H))`$.
(vi) Set $`H_k=_{j=0}^{k1}\alpha ^j(H)`$ and $`H_k^{}=_{j=0}^{k1}\alpha ^j(H+\alpha (K)K)=H_k+\alpha ^k(K)K`$. Then by (iv) and (v) we have
$$|P_\alpha (H)P_\alpha (H+\alpha (K)K)|=\frac{1}{k}|P_{\alpha ^k}(H_k)P_{\alpha ^k}(H_k^{})|\frac{2K}{k}.$$
Thus (vi) follows.
The next result is the analogue of \[W, Theorem 9.12\], see also \[R\].
###### Proposition 2.5
Suppose $`ht(\alpha )<\mathrm{}`$. Let $`\varphi `$ be a self-adjoint linear functional on $`A`$. Then $`\varphi `$ is an $`\alpha `$-invariant state if and only if $`\varphi (H)P_\alpha (H)`$ for all $`HA_{sa}`$.
Proof. If $`\varphi `$ is an $`\alpha `$-invariant state then by Proposition 2.2
$$\varphi (H)P_\alpha (H)h_\varphi (\alpha )P_\alpha (H)$$
for all $`HA_{sa}`$.
Conversely if $`\varphi (H)P_\alpha (H)`$ for all $`HA_{sa}`$ then by Proposition 2.4(vi)
$$\varphi (\alpha (H)H)=\frac{1}{n}\varphi (\alpha (nH)nH)\frac{1}{n}P_\alpha (\alpha (nH)nH)=\frac{1}{n}P_\alpha (0)=\frac{1}{n}ht(\alpha )\underset{n\mathrm{}}{}0.$$
Applying this also to $`H`$ we see that $`\varphi `$ is $`\alpha `$-invariant. Furthermore, by Proposition 2.4(i),(ii)
$$\varphi (H)=\frac{1}{n}\varphi (nH)\frac{1}{n}P_\alpha (nH)\frac{1}{n}ht(\alpha )+H\underset{n\mathrm{}}{}H,$$
so that $`\varphi 1`$. For $`c`$ we have
$$c\varphi (1)P_\alpha (c1)=ht(\alpha )c.$$
Hence $`\varphi (1)=1`$, and $`\varphi `$ is a state.
###### Definition 2.6
We say $`\varphi `$ is an equilibrium state at $`H`$ if
$$P_\alpha (H)=h_\varphi (\alpha )\varphi (H),$$
hence by Proposition 2.2
$$h_\varphi (\alpha )\varphi (H)=\underset{\psi }{sup}(h_\psi (\alpha )\psi (H)),$$
where the sup is taken over all $`\alpha `$-invariant states.
Recall that if $`F`$ is a real convex continuous function on a real Banach space $`X`$, then a linear functional $`f`$ on $`X`$ is called a tangent functional to the graph of $`F`$ at the point $`x_0X`$ if
$$F(x_0+x)F(x)f(x)xX.$$
In the sequel we will identify self-adjoint linear functionals on $`A`$ with real linear functionals on $`A_{sa}`$.
The next proposition is the analogue of \[W, Theorem 9.14\].
###### Proposition 2.7
Suppose $`ht(\alpha )<\mathrm{}`$ and the pressure is a convex function on $`A_{sa}`$. Then
* (i) if $`\varphi `$ is an equilibrium state at $`H`$ then $`\varphi `$ is a tangent functional for the pressure at $`H`$;
* (ii) if $`\varphi `$ is a tangent functional for the pressure at $`H`$ then $`\varphi `$ is an $`\alpha `$-invariant state.
Proof. (i) Let $`KA_{sa}`$. Then by Proposition 2.2
$$P_\alpha (H+K)P_\alpha (H)(h_\varphi (\alpha )\varphi (H+K))(h_\varphi (\alpha )\varphi (H))=\varphi (K),$$
so $`\varphi `$ is a tangent functional.
(ii) If $`KA_{sa}`$ then by Proposition 2.4(vi)
$$\varphi (\alpha (K)K)P_\alpha (H+\alpha (K)K)P_\alpha (H)=0.$$
Applying this also to $`K`$ we see that $`\varphi `$ is $`\alpha `$-invariant. Now note that $`\varphi 1`$ by Proposition 2.4(iv). By Proposition 2.4(ii) we have also $`cc\varphi (1)`$ for any $`c`$. Hence $`\varphi (1)=1`$, and $`\varphi `$ is a state.
## 3 The variational principle
We shall prove the variational principle for the following class of C-dynamical systems.
###### Definition 3.1
A unital C-dynamical system $`(A,\alpha )`$ is called asymptotically abelian with locality if there is a dense $`\alpha `$-invariant $``$-subalgebra $`𝒜`$ of $`A`$ such that for each pair $`a,b𝒜`$ the C-algebra generated by $`a`$ and $`b`$ is finite dimensional, and for some $`p=p(a,b)`$ we have $`[\alpha ^j(a),b]=0`$ whenever $`|j|p`$.
We call elements of $`𝒜`$ for local operators and finite dimensional C-subalgebras of $`𝒜`$ for local algebras.
Note that since we may add the identity operator to $`𝒜`$, we may assume that $`1𝒜`$. Since each finite dimensional C-algebra is singly generated, an easy induction argument shows that the C-algebra generated by a finite set of local operators is finite dimensional. In particular, $`A`$ is an AF-algebra. Note also that another easy induction argument shows that for each local algebra $`N`$ there is $`p`$ such that $`[\alpha ^j(a),b]=0`$ for all $`a,bN`$ whenever $`|j|p`$.
###### Theorem 3.2
Let $`(A,\alpha )`$ be a unital separable C-dynamical system which is asymptotically abelian with locality. Let $`HA_{sa}`$. Then
$$P_\alpha (H)=\underset{\varphi }{sup}(h_\varphi (\alpha )\varphi (H)),$$
where the sup is taken over all $`\alpha `$-invariant states of $`A`$. In particular, the topological entropy
$$ht(\alpha )=\underset{\varphi }{sup}h_\varphi (\alpha ).$$
Consider first the case when there exists a finite dimensional C-subalgebra $`N`$ of $`A`$ such that $`HN`$, $`\alpha ^j(N)`$ commutes with $`N`$ for $`j0`$, $`_j\alpha ^j(N)=A`$.
###### Lemma 3.3
Under the above assumptions there exists an $`\alpha `$-invariant state $`\varphi `$ such that
$$P_\alpha (H)=h_\varphi (\alpha )\varphi (H)=\underset{n\mathrm{}}{lim}\frac{1}{n}\mathrm{log}\mathrm{Tr}_{_{j=0}^{n1}\alpha ^j(N)}\left(e^{_{j=0}^{n1}\alpha ^j(H)}\right).$$
Proof. First note that if $`A_1`$ and $`A_2`$ are commuting finite dimensional C-algebras, and $`a_iA_i`$, $`a_i0`$, $`i=1,2`$, then
$$\mathrm{Tr}_{A_1A_2}(a_1a_2)\mathrm{Tr}_{A_1}(a_1)\mathrm{Tr}_{A_2}(a_2),$$
since if $`p_i`$ is a minimal projection in $`A_i`$, $`i=1,2`$, then $`p_1p_2`$ is either zero or minimal in $`A_1A_2`$. Hence the limit in the formulation of the lemma really exists. We denote it by $`\stackrel{~}{P}_\alpha (H)`$. It is easy to see that $`\stackrel{~}{P}_\alpha (H)P_\alpha (H)`$.
For each $`n`$ let $`N_n`$ be a copy of $`N`$. Consider the C-algebra $`M=_nN_n`$. Let $`\beta `$ be the shift to the right on $`M`$, and $`\pi :MA`$ the homomorphism which intertwines $`\beta `$ with $`\alpha `$, and identifies $`N_0`$ with $`N`$. Set $`I=\mathrm{Ker}\pi `$. For each $`n`$ let
$$M_n=N_0\mathrm{}N_{n1},I_n=IM_n,\pi _n=\pi |_{M_n}.$$
Identifying $`M`$ with $`M_n^{}`$ consider the $`\beta ^n`$-invariant state $`\psi _n=(f_n\pi _n)`$ on $`M`$, where $`f_n`$ is the state on $`_{j=0}^{n1}\alpha ^j(N)`$ with density operator
$$\left(\mathrm{Tr}_{_{j=0}^{n1}\alpha ^j(N)}e^{_{j=0}^{n1}\alpha ^j(H)}\right)^1e^{_{j=0}^{n1}\alpha ^j(H)}.$$
Set $`\varphi _n={\displaystyle \frac{1}{n}}{\displaystyle \underset{j=0}{\overset{n1}{}}}\psi _n\beta ^j`$. Then $`\varphi _n`$ is $`\beta `$-invariant. Using concavity of entropy we obtain
$`h_{\varphi _n}(\beta )`$ $`=`$ $`{\displaystyle \frac{1}{n}}h_{\varphi _n}(\beta ^n){\displaystyle \frac{1}{n^2}}{\displaystyle \underset{j=0}{\overset{n1}{}}}h_{\psi _n\beta ^j}(\beta ^n)={\displaystyle \frac{1}{n}}h_{\psi _n}(\beta ^n)={\displaystyle \frac{1}{n}}S(f_n\pi _n)={\displaystyle \frac{1}{n}}S(f_n)`$
$`=`$ $`{\displaystyle \frac{1}{n}}\mathrm{log}\mathrm{Tr}_{_{j=0}^{n1}\alpha ^j(N)}\left(e^{_{j=0}^{n1}\alpha ^j(H)}\right)+{\displaystyle \frac{1}{n}}f_n\left({\displaystyle \underset{j=0}{\overset{n1}{}}}\alpha ^j(H)\right)`$
$``$ $`\stackrel{~}{P}_\alpha (H)+{\displaystyle \frac{1}{n}}f_n\left({\displaystyle \underset{j=0}{\overset{n1}{}}}\alpha ^j(H)\right)`$
$`=`$ $`\stackrel{~}{P}_\alpha (H)+{\displaystyle \frac{1}{n}}\psi _n\left({\displaystyle \underset{j=0}{\overset{n1}{}}}\beta ^j(H)\right)=\stackrel{~}{P}_\alpha (H)+\varphi _n(H).`$
Let $`\stackrel{~}{\varphi }`$ be any weak limit point of the sequence $`\{\varphi _n\}_n`$. Then $`\stackrel{~}{\varphi }`$ is $`\beta `$-invariant. Let $`B`$ be a masa in $`N_0`$ containing $`H`$. Then $`B`$ is in the centralizer of the state $`\varphi _n`$, hence
$$h_{\varphi _n}(\beta )=h_{\varphi _n}(B;\beta ).$$
Since the mapping $`\psi h_\psi (B;\beta )`$ is upper semicontinuous, we conclude that
$$h_{\stackrel{~}{\varphi }}(\beta )\stackrel{~}{P}_\alpha (H)+\stackrel{~}{\varphi }(H).$$
Now note that $`\stackrel{~}{\varphi }`$ is zero on $`I`$. Indeed, if $`xI_n`$ then $`\beta ^j(x)I_m`$ for $`j=0,\mathrm{},mn`$ and $`mn`$, whence
$$|\varphi _m(x)|\frac{1}{m}\underset{j=mn+1}{\overset{m1}{}}|(\psi _m\beta ^j)(x)|\frac{n1}{m}x,$$
so $`\stackrel{~}{\varphi }(x)=0`$. Thus $`\stackrel{~}{\varphi }`$ defines a state $`\varphi `$ on $`A`$. We have
$$h_\varphi (\alpha )=h_{\stackrel{~}{\varphi }}(\beta )\stackrel{~}{P}_\alpha (H)+\varphi (H),$$
where the first equality follows from \[CNT, Theorem VII.2\]. Since by Proposition 2.2, $`h_\varphi (\alpha )\varphi (H)P_\alpha (H)\stackrel{~}{P}_\alpha (H)`$, the proof of the lemma is complete.
We shall reduce the general case to the case considered above by replacing $`\alpha `$ by its powers. For this suppose that $`N`$ is a local subalgebra of $`A`$, and $`HN`$. Choose $`p`$ such that $`\alpha ^j(N)`$ commutes with $`N`$ whenever $`|j|p`$. For $`kp`$ set $`M_k=_{j=0}^{kp}\alpha ^j(N)`$, $`H_k=_{j=0}^{kp}\alpha ^j(H)`$. Then $`H_kM_k`$, and $`\alpha ^{jk}(M_k)`$ commutes with $`M_k`$ for $`j0`$.
###### Lemma 3.4
For any finite subset $`\omega `$ of $`N`$ we have
$$P_\alpha (H,\omega )\underset{k\mathrm{}}{lim\; inf}\frac{1}{k}P_{\alpha ^k|_j\alpha ^{jk}(M_k)}(H_k).$$
Proof. The idea of the proof is to reduce to the situation of Lemma 3.3 by showing that the contribution of the indicies in the intervals $`[jkp+1,jk1]`$, $`j`$, becomes negligible for large $`k`$.
Fix $`\delta >0`$. Choose $`m_0`$ such that
$$\frac{2(p1)a}{m_0}<\delta \text{for}a\omega .$$
Take any $`km_0+p`$. Let $`n`$. Then $`(m1)kn<mk`$ for some $`m`$. Set $`B_0=_{j=0}^m\alpha ^{jk}(M_k)`$ and
$$B=\underset{m_0}{\underset{}{B_0\mathrm{}B_0}}.$$
Choose a conditional expectation $`E:AB_0`$, and define unital completely positive mappings $`\psi :BA`$ and $`\rho :AB`$ as follows:
$$\psi (b_1,\mathrm{},b_{m_0})=\frac{1}{m_0}\underset{i=1}{\overset{m_0}{}}\alpha ^{i+1}(b_i),$$
$$\rho (a)=(E(a),(E\alpha )(a),\mathrm{},(E\alpha ^{m_01})(a)).$$
For any $`aA`$ we have
$$(\psi \rho )(a)a\frac{2a}{m_0}\mathrm{\#}\{0im_01|\alpha ^i(a)B_0\},$$
where $`\mathrm{\#}S`$ means the cardinality of a set $`S`$. Let $`a=\alpha ^l(b)`$ for some $`b\omega `$ and $`l`$, $`0ln1`$. Then $`l=jk+r`$ for some $`j`$ and $`r`$, $`0jm1`$, $`0r<k`$. Since $`m_0kp`$, the interval $`[l,l+m_01]`$ is contained in $`[jk,(j+1)k+kp]`$. But for $`i[jk,(j+1)k+kp]\backslash [jk+kp+1,(j+1)k1]`$ we have $`\alpha ^i(N)B_0`$, so $`\mathrm{\#}\{0im_01|\alpha ^i(a)B_0\}p1`$, and
$$(\psi \rho )(a)a\frac{2(p1)b}{m_0}<\delta .$$
Hence
$$P(\underset{j=0}{\overset{n1}{}}\alpha ^j(H),\underset{j=0}{\overset{n1}{}}\alpha ^j(\omega );\delta )\mathrm{log}\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(H)\right)}\right).$$
Now note that for $`0im_01`$ the sets $`X_i=[i,i+n1]`$ and $`X=_{j=0}^m[jk,jk+kp]`$ are contained in $`Y=[0,mk+kp]`$, so
$`\mathrm{\#}(X_iX)`$ $``$ $`\mathrm{\#}(Y\backslash X_i)+\mathrm{\#}(Y\backslash X)=(mk+kp+1n)+m(p1)`$
$``$ $`mk+kp+1(m1)k+m(p1)mp+2k.`$
When $`jX_iX`$, $`\alpha ^j(H)B_0`$. Hence
$$(E\alpha ^i)\left(\underset{j=0}{\overset{n1}{}}\alpha ^j(H)\right)\underset{j=0}{\overset{m}{}}\alpha ^{jk}(H_k)=E\left(\underset{jX_i}{}\alpha ^j(H)\right)\underset{jX}{}\alpha ^j(H)(mp+2k)H.$$
By the Peierls-Bogoliubov inequality we obtain
$$\mathrm{Tr}_{B_0}\left(e^{(E\alpha ^i)\left(_{j=0}^{n1}\alpha ^j(H)\right)}\right)e^{(mp+2k)H}\mathrm{Tr}_{B_0}\left(e^{_{j=0}^m\alpha ^{jk}(H_k)}\right),$$
so
$$\mathrm{Tr}_B\left(e^{\rho \left(_{j=0}^{n1}\alpha ^j(H)\right)}\right)m_0e^{(mp+2k)H}\mathrm{Tr}_{B_0}\left(e^{_{j=0}^m\alpha ^{jk}(H_k)}\right).$$
Taking the $`\mathrm{log}`$, dividing by $`n`$, and letting $`n\mathrm{}`$, we obtain
$`P_\alpha (H,\omega ;\delta )`$ $``$ $`{\displaystyle \frac{pH}{k}}+{\displaystyle \frac{1}{k}}\underset{m\mathrm{}}{lim}{\displaystyle \frac{1}{m}}\mathrm{log}\mathrm{Tr}_{_{j=0}^{m1}\alpha ^{jk}(M_k)}\left(e^{_{j=0}^{m1}\alpha ^{jk}(H_k)}\right)`$
$`=`$ $`{\displaystyle \frac{pH}{k}}+{\displaystyle \frac{1}{k}}P_{\alpha ^k|_j\alpha ^{jk}(M_k)}(H_k),`$
where the last equality follows from Lemma 3.3.
We shall need also the following
###### Lemma 3.5
Let $`(A,\alpha )`$ be a C-dynamical system with A nuclear, $`B`$ an $`\alpha `$-invariant C-subalgebra of $`A`$, $`\varphi `$ an $`\alpha `$-invariant state on $`B`$. Then for any $`\epsilon >0`$ there exists an $`\alpha `$-invariant state $`\psi `$ on $`A`$ such that
$$\psi |_B=\varphi \text{and}h_\psi (\alpha )>h_\varphi (\alpha |_B)\epsilon .$$
Proof. Since the Sauvageot-Thouvenot entropy is not less than the CNT-entropy for general C-systems, there exist a commutative C-dynamical system $`(C,\beta ,\mu )`$, an $`(\alpha \beta )`$-invariant state $`\lambda `$ on $`BC`$, and a finite dimensional subalgebra $`P`$ of $`C`$ such that $`\lambda |_B=\varphi `$, $`\lambda |_C=\mu `$ and
$$h_\varphi (\alpha |_B)<H_\mu (P,P^{})H_\lambda (P|B)+\epsilon ,$$
see \[ST\] for notations. Extend $`\lambda `$ to an $`(\alpha \beta )`$-invariant state $`\mathrm{\Lambda }`$ on $`AC`$, and set $`\psi =\mathrm{\Lambda }|_A`$. Since the conditional entropy $`H_\lambda (P|B)`$ is decreasing in second variable, and ST-entropy coincides with CNT-entropy for nuclear algebras, we have
$$h_\psi (\alpha )H_\mu (P,P^{})H_\mathrm{\Lambda }(P|A)H_\mu (P,P^{})H_\lambda (P|B)>h_\varphi (\alpha |_B)\epsilon .$$
Proof of Theorem 3.2. The inequality ’$``$’ has been proved in Proposition 2.2. Since the pressure is continuous by Proposition 2.4(iv), to prove the converse inequality it suffices to consider local $`H`$. Then by Lemma 3.4 we have only to show that if $`H`$ is contained in a local algebra $`N`$ then
$$\underset{\varphi }{sup}(h_\varphi (\alpha )\varphi (H))\underset{k\mathrm{}}{lim\; inf}\frac{1}{k}P_{\alpha ^k|_j\alpha ^{jk}(M_k)}(H_k).$$
By Lemma 3.3, for each $`k`$ there exists an $`\alpha ^k`$-invariant state $`\psi _k`$ on $`_j\alpha ^{jk}(M_k)`$ such that
$$h_{\psi _k}(\alpha ^k|_{_j\alpha ^{jk}(M_k)})\psi _k(H_k)=P_{\alpha ^k|_j\alpha ^{jk}(M_k)}(H_k).$$
By Lemma 3.5 we may extend $`\psi _k`$ to an $`\alpha ^k`$-invariant state $`\stackrel{~}{\varphi }_k`$ on $`A`$ such that
$$h_{\stackrel{~}{\varphi }_k}(\alpha ^k)h_{\psi _k}(\alpha ^k|_{_j\alpha ^{jk}(M_k)})1.$$
Set $`\varphi _k={\displaystyle \frac{1}{k}}{\displaystyle \underset{j=0}{\overset{k1}{}}}\stackrel{~}{\varphi }_k\alpha ^j`$. Then as in the proof of Lemma 3.3
$$h_{\varphi _k}(\alpha )\frac{1}{k}h_{\stackrel{~}{\varphi }_k}(\alpha ^k)\frac{1}{k}h_{\psi _k}(\alpha ^k|_{_j\alpha ^{jk}(M_k)})\frac{1}{k}.$$
Since
$$\varphi _k(H)=\frac{1}{k}\underset{j=0}{\overset{k1}{}}\stackrel{~}{\varphi }_k(\alpha ^j(H))\frac{1}{k}\stackrel{~}{\varphi }_k(H_k)+\frac{p1}{k}H=\frac{1}{k}\psi _k(H_k)+\frac{p1}{k}H,$$
we get
$`h_{\varphi _k}(\alpha )\varphi _k(H)`$ $``$ $`{\displaystyle \frac{1}{k}}h_{\psi _k}(\alpha ^k|_{_j\alpha ^{jk}(M_k)}){\displaystyle \frac{1}{k}}\psi _k(H_k){\displaystyle \frac{1+(p1)H}{k}}`$
$`=`$ $`{\displaystyle \frac{1}{k}}P_{\alpha ^k|_j\alpha ^{jk}(M_k)}(H_k){\displaystyle \frac{1+(p1)H}{k}},`$
and the proof is complete.
###### Corollary 3.6
With our assumptions the pressure is a convex function of $`H`$.
Proof. Use the affinity of the function $`Hh_\varphi (\alpha )\varphi (H)`$.
###### Corollary 3.7
If $`(A_1,\alpha _1)`$ and $`(A_2,\alpha _2)`$ are asymptotically abelian systems with locality then
$$ht(\alpha _1\alpha _2)=ht(\alpha _1)+ht(\alpha _2).$$
Proof. If $`\varphi _i`$ is an $`\alpha _i`$-invariant state, $`i=1,2`$, then by \[SV, Lemma 3.4\] and \[V, Propositions 4.6 and 4.9\],
$$h_{\varphi _1}(\alpha _1)+h_{\varphi _2}(\alpha _2)h_{\varphi _1\varphi _2}(\alpha _1\alpha _2)ht(\alpha _1\alpha _2)ht(\alpha _1)+ht(\alpha _2).$$
Taking the sup over $`\varphi _i`$ we get the conclusion.
## 4 KMS-states
By Corollary 3.6 and Proposition 2.7 it follows that if $`(A,\alpha )`$ is asymptotically abelian with locality and $`ht(\alpha )<\mathrm{}`$, then for every equilibrium state $`\varphi `$ at $`H`$, $`\varphi `$ is a tangent functional for the pressure $`P_\alpha `$ at $`H`$. Furthermore, if $`\omega `$ is a tangent functional for $`P_\alpha `$ at $`H`$ then $`\omega `$ is an $`\alpha `$-invariant state.
If $`H`$ is local and $`I`$ is a subset then the derivation
$$\delta _{H,I}(x)=\underset{jI}{}[\alpha ^j(H),x],x𝒜,$$
defines a strongly continuous one-parameter automorphism group $`\sigma _t^{H,I}=\mathrm{exp}(it\delta _{H,I})`$ of $`A`$ (see \[BR, Theorem 6.2.6 and Example 6.2.8\]). We shall mainly be concerned with the case $`I=`$, and will write $`\delta _H=\delta _{H,}`$, $`\sigma _t^H=\sigma _t^{H,}`$. Recall that a state $`\varphi `$ is a $`(\sigma _t^H,\beta )`$-KMS state if $`\varphi (ab)=\varphi (b\sigma _{i\beta }^H(a))`$ for $`\sigma _t^H`$-analytic elements $`a,bA`$.
We say that an $`\alpha `$-invariant state $`\varphi `$ is an equilibrium state at $`H`$ at inverse temperature $`\beta `$ if
$$P_\alpha (\beta H)=h_\varphi (\alpha )\beta \varphi (H).$$
By Theorem 3.2, for systems which are asymptotically abelian with locality, this is equivalent to
$$h_\varphi (\alpha )\beta \varphi (H)=\underset{\psi }{sup}(h_\psi (\alpha )\beta \psi (H)).$$
The main result in this section is
###### Theorem 4.1
Suppose a unital separable C-dynamical system $`(A,\alpha )`$ is asymptotically abelian with locality, and $`ht(\alpha )<\mathrm{}`$. If $`H`$ is a local self-adjoint operator in $`A`$ and $`\varphi `$ is an equilibrium state at $`H`$ at inverse temperature $`\beta `$, then $`\varphi `$ is a $`(\sigma _t^H,\beta )`$-KMS state. In particular, if $`ht(\alpha )=h_\varphi (\alpha )`$ then $`\varphi `$ is a trace.
In order to prove the theorem we may replace $`H`$ by $`\beta H`$ and show that $`\varphi `$ is a $`(\sigma _t^H,1)`$-KMS state. We shall prove the following more general result.
###### Theorem 4.2
If $`\varphi `$ is a tangent functional for $`P_\alpha `$ at $`H`$ then $`\varphi `$ is a $`(\sigma _t^H,1)`$-KMS state.
We shall need an explicit formula for the pressure, which is a consequence of our proof of the variational principle.
###### Lemma 4.3
Let $`N`$ be a local algebra. Then there exist a sequence $`\{A_n\}_n`$ of local algebras containing $`N`$ and three sequences $`\{p_n\}_n`$, $`\{m_n\}_n`$, $`\{k_n\}_n`$ of positive integers such that
* (i) $`\alpha ^p(A_n)`$ commutes with $`A_n`$ whenever $`|p|p_n`$;
* (ii) $`{\displaystyle \frac{p_n}{k_n}}0`$ as $`n\mathrm{}`$;
* (iii) $`P_\alpha (H)=\underset{n\mathrm{}}{lim}{\displaystyle \frac{1}{k_nm_n}}\mathrm{log}\mathrm{Tr}_{_{jI_n}\alpha ^j(A_n)}\left(e^{_{jI_n}\alpha ^j(H)}\right)`$ for all $`HN_{sa}`$, where
$`I_n={\displaystyle \underset{j=0}{\overset{m_n1}{}}}[jk_n,jk_n+k_np_n]`$.
Proof. Let $`\{A_n\}_n`$ be an increasing sequence of local algebras containing $`N`$ such that $`_nA_n`$ is dense in $`A`$, $`\omega _n`$ a finite subset of $`A_n`$ such that $`\mathrm{span}(\omega _n)=A_n`$. Let $`\{p_n\}_n`$ be a sequence satisfying condition (i). By Lemma 3.4
$$P_\alpha (H,\omega _n)\underset{k\mathrm{}}{lim\; inf}\frac{1}{k}P_{\alpha ^k|_j\alpha ^{jk}(A_{n,k})}\left(\underset{j=0}{\overset{kp_n}{}}\alpha ^j(H)\right)HN_{sa},$$
where $`A_{n,k}=_{j=0}^{kp_n}\alpha ^j(A_n)`$. On the other hand, by the proof of Theorem 3.2
$$P_\alpha (H)\underset{k\mathrm{}}{lim\; sup}\frac{1}{k}P_{\alpha ^k|_j\alpha ^{jk}(A_{n,k})}\left(\underset{j=0}{\overset{kp_n}{}}\alpha ^j(H)\right)HN_{sa}.$$
Choose a countable dense subset $`X`$ of $`N_{sa}`$. Since $`P_\alpha (H,\omega _n)P_\alpha (H)`$ for any $`HX`$, we can find a sequence $`\{k_n\}_n`$ such that condition (ii) is satisfied and
$$P_\alpha (H)=\underset{n\mathrm{}}{lim}\frac{1}{k_n}P_{\alpha ^{k_n}|_j\alpha ^{jk_n}(A_{n,k_n})}\left(\underset{j=0}{\overset{k_np_n}{}}\alpha ^j(H)\right)HX.$$
Since by Lemma 3.3
$$P_{\alpha ^{k_n}|_j\alpha ^{jk_n}(A_{n,k_n})}\left(\underset{j=0}{\overset{k_np_n}{}}\alpha ^j(H)\right)=\underset{m\mathrm{}}{lim}\frac{1}{m}\mathrm{log}\mathrm{Tr}_{_{jI_{n,m}}\alpha ^j(A_n)}\left(e^{_{jI_{n,m}}\alpha ^j(H)}\right),$$
where $`I_{n,m}=_{j=0}^{m1}[jk_n,jk_n+k_np_n]`$, we can choose a sequence $`\{m_n\}_n`$ such that condition (iii) is satisfied for all $`HX`$. But then it is satisfied for all $`HN_{sa}`$ by Proposition 2.4(iv) and the Peierls-Bogoliubov inequality.
Every local operator is analytic for the dynamics, and $`\sigma _t^H`$ depends continuously on $`H`$ in a fixed local algebra. More precisely, we have
###### Lemma 4.4
(i) The series $`\sigma _\beta ^{H,I}(a)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(i\beta )^n}{n!}}\delta _{H,I}^n(a)`$ converges absolutely in norm for any $`\beta `$ and any local operator $`a`$.
(ii) Given a local algebra $`N`$, $`R>0`$, $`C>0`$ and $`\epsilon >0`$ there exist $`q`$ and $`\delta >0`$ such that
$$\sigma _\beta ^{H_1,I_1}(a)\sigma _\beta ^{H_2,I_2}(a)\epsilon a$$
$`aN`$, $`H_1,H_2N_{sa}`$ with $`H_1,H_2C`$ and $`H_1H_2<\delta `$, $`\beta `$ with $`|\beta |R`$, $`I_1,I_2`$ with $`[q,q]I_1I_2`$.
Proof. We shall use the arguments of Araki \[A, Theorem 4.2\]. Let $`H`$ and $`a`$ lie in a local algebra $`N`$. Choose $`p`$ such that $`\alpha ^j(N)`$ commutes with $`N`$ for $`|j|p`$. Then
$$\delta _{H,I}^m(a)=\underset{j_1,\mathrm{},j_m}{}[\alpha ^{j_m}(H),[\mathrm{},[\alpha ^{j_1}(H),a]\mathrm{}]],$$
where the sum is over all $`j_1,\mathrm{},j_m`$ such that
$$j_k[p,p]\left(\underset{l<k}{}[j_lp,j_l+p]\right)$$
(4.1)
for each $`k=1,\mathrm{},m`$. But as was already noted in \[GN\] condition (4.1) is equivalent to
$$[j_k,j_k+p]\left([0,p]\left(\underset{l<k}{}[j_l,j_l+p]\right)\right)\mathrm{}.$$
Thus the lemma follows from the proof of \[A, Theorem 4.2\] (with $`n=p`$ and $`r=p`$).
The following lemma contains the main technical result needed to prove Proposition 4.2.
###### Lemma 4.5
Let $`N`$ be a local algebra, $`HN_{sa}`$, $`\varphi N^{}`$ is a tangent functional to $`(P_\alpha )|_{N_{sa}}`$ at $`H`$. Let $`E:AN`$ be a conditional expectation. Then for any function $`f𝒟`$ (the space of $`C^{\mathrm{}}`$-functions with compact support) and any $`a,bN`$ we have
$`\left|{\displaystyle _{}}\widehat{f}(t)\varphi (aE(\sigma _t^H(b)))𝑑t{\displaystyle _{}}\widehat{f}(t+i)\varphi (E(\sigma _t^H(b))a)𝑑t\right|`$
$$a_{}(|\widehat{f}(t)|+|\widehat{f}(t+i)|)\sigma _t^H(b)E(\sigma _t^H(b))𝑑t.$$
Proof. First consider the case when $`(P_\alpha )|_{N_{sa}}`$ is differentiable at $`H`$, in other words $`\varphi `$ is the unique tangent functional. With the notations of Lemma 4.3 consider the state $`f_n`$ on $`_{jI_n}\alpha ^j(A_n)`$ with density operator
$$\left(\mathrm{Tr}_{_{jI_n}\alpha ^j(A_n)}\left(e^{_{jI_n}\alpha ^j(H)}\right)\right)^1e^{_{jI_n}\alpha ^j(H)}.$$
Then define a positive linear functional $`\varphi _n`$ on $`N`$ by
$$\varphi _n(x)=\frac{1}{k_nm_n}\underset{jI_n}{}f_n(\alpha ^j(x)).$$
Note that $`\varphi _n=\varphi _n(1)1`$. Since $`f_n`$ is a tangent functional to the convex function $`x\mathrm{log}\mathrm{Tr}_{_{jI_n}\alpha ^j(A_n)}(e^x)`$ on $`(_{jI_n}\alpha ^j(A_n))_{sa}`$ at the point $`_{jI_n}\alpha ^j(H)`$, $`\varphi _n`$ is a tangent functional to the function $`N_{sa}x\frac{1}{k_nm_n}\mathrm{log}\mathrm{Tr}_{_{jI_n}\alpha ^j(A_n)}\left(e^{_{jI_n}\alpha ^j(x)}\right)`$ at $`H`$. It follows that any limit point of the sequence $`\{\varphi _n\}_n`$ is a tangent functional to $`(P_\alpha )|_{N_{sa}}`$ at $`H`$. Since the latter is unique by assumption, $`\varphi _n\varphi `$ as $`n\mathrm{}`$.
Since $`f_n`$ is a $`(\sigma _t^{H,I_n},1)`$-KMS state, by \[BR, Proposition 5.3.12\] we have
$$_{}\widehat{f}(t)f_n(\alpha ^j(a)\sigma _t^{H,I_n}(\alpha ^j(b)))𝑑t=_{}\widehat{f}(t+i)f_n(\sigma _t^{H,I_n}(\alpha ^j(b))\alpha ^j(a))𝑑tjI_n.$$
(4.2)
Note that $`\sigma _t^{H,I_n}(\alpha ^j(b))=\alpha ^j(\sigma _t^{H,I_nj}(b))`$. Fix $`q`$, and set $`I_{n,q}=_{j=0}^{m_n1}[jk_n+q,jk_n+k_np_nq]`$. By Lemma 4.4, if $`q`$ is large enough then $`\sigma _t^{H,I_nj}(b)`$ is arbitrarily close to $`\sigma _t^H(b)`$ for any $`jI_{n,q}`$ and any $`t`$ in a fixed compact subset of $``$. But then $`\sigma _t^{H,I_n}(\alpha ^j(b))\alpha ^j(E(\sigma _t^H(b)))`$ is arbitrarily close to $`\alpha ^j(\sigma _t^H(b)E(\sigma _t^H(b)))`$. In other words,
$`\left|{\displaystyle _{}}𝑑t\widehat{f}(t){\displaystyle \frac{1}{k_nm_n}}{\displaystyle \underset{jI_{n,q}}{}}f_n\left(\alpha ^j(a)\sigma _t^{H,I_n}(\alpha ^j(b))\alpha ^j(a)\alpha ^j(E(\sigma _t^H(b)))\right)\right|`$
$$a_{}|\widehat{f}(t)|\sigma _t^H(b)E(\sigma _t^H(b))𝑑t+\epsilon (q)n,$$
where $`\epsilon (q)0`$ as $`q\mathrm{}`$. Since $`\mathrm{\#}I_{n,q}/\mathrm{\#}I_n1`$ as $`n\mathrm{}`$, letting $`n\mathrm{}`$ we may replace averaging over the set $`I_{n,q}`$ by averaging over $`I_n`$, and then obtain
$`\underset{n\mathrm{}}{lim\; sup}\left|{\displaystyle \frac{1}{k_nm_n}}{\displaystyle \underset{jI_n}{}}{\displaystyle _{}}\widehat{f}(t)f_n\left(\alpha ^j(a)\sigma _t^{H,I_n}(\alpha ^j(b))\right)𝑑t{\displaystyle _{}}\widehat{f}(t)\varphi (aE(\sigma _t^H(b)))𝑑t\right|`$
$$a_{}|\widehat{f}(t)|\sigma _t^H(b)E(\sigma _t^H(b))𝑑t.$$
Since an analogous estimate holds for $`_{}\widehat{f}(t+i)\varphi (E(\sigma _t^H(b))a)𝑑t`$, we obtain the conclusion of the lemma by virtue of (4.2).
If $`(P_\alpha )|_{N_{sa}}`$ is not differentiable at $`H`$ then by \[LR, Theorem 1\], $`\varphi `$ lies in the closed convex hull of those $`\stackrel{~}{\varphi }`$, for which there exists a sequence $`\{H_n\}_nN_{sa}`$ converging to $`H`$ such that $`(P_\alpha )|_{N_{sa}}`$ has a unique tangent functional $`\varphi _n`$ at $`H_n`$ and $`\varphi _n\stackrel{~}{\varphi }`$. Since for $`\varphi _n`$ the lemma is already proved (for $`H_n`$ instead of $`H`$), using Lemma 4.4 we conclude that the conclusion of the lemma is true for $`\stackrel{~}{\varphi }`$. But then it is true for any functional in the closed convex hull of the $`\stackrel{~}{\varphi }`$’s.
Proof of Theorem 4.2. If $`\varphi `$ is a tangent functional for $`P_\alpha `$ at $`H`$ then $`\varphi |_N`$ is a tangent functional for $`(P_\alpha )|_{N_{sa}}`$ at $`H`$ for any local algebra $`N`$ containing $`H`$. Thus by Lemma 4.5 the equality
$$_{}\widehat{f}(t)\varphi (a\sigma _t^H(b))𝑑t=_{}\widehat{f}(t+i)\varphi (\sigma _t^H(b)a)𝑑t$$
holds for all $`f𝒟`$ and all local $`a,b`$, hence for all $`a,bA`$. By \[BR, Proposition 5.3.12\] this is equivalent to the KMS-condition.
###### Remark 4.6
Under the assumptions of Theorem 4.1, if
$$\varphi \underset{\epsilon >0}{}\overline{\{\psi |P_\alpha (H)<h_\psi (\alpha )\psi (H)+\epsilon \}}$$
(weak closure), then $`\varphi `$ is a tangent functional for the pressure at $`H`$, hence $`\varphi `$ is a $`(\sigma _t^H,1)`$-KMS state. In other words, any weak limit point of a sequence on which the sup in the variational principle is attained, is a $`(\sigma _t^H,1)`$-KMS state. If $`ht(\alpha )=+\mathrm{}`$, this is of course false in general. Moreover, for any $`\alpha `$-invariant state $`\varphi `$ there exists a sequence $`\{\varphi _n\}_n`$ converging in norm to $`\varphi `$ such that $`h_{\varphi _n}(\alpha )=+\mathrm{}`$ for all $`n`$. Indeed, first note that by taking infinite convex combinations of states of large entropy we can find a state $`\psi `$ of infinite entropy. Then $`\varphi _n=\frac{1}{n}\psi +\frac{n1}{n}\varphi \underset{n\mathrm{}}{}\varphi `$ and $`h_{\varphi _n}(\alpha )=+\mathrm{}`$.
## 5 Examples
First we consider a class of systems arising naturally from systems of topological dynamics.
Let $`\sigma `$ be an expansive homeomorphism of a zero-dimensional compact space $`X`$, $`G`$ the group of uniformly finite-dimensional homeomorphisms of $`X`$ in the sense of Krieger \[K\]. By definition, a homeomorphism $`T`$ belongs to $`G`$ if
$$\underset{|n|\mathrm{}}{lim}\underset{xX}{sup}d(\sigma ^nTx,\sigma ^nx)=0,$$
where $`d`$ is a metric defining the topology of $`X`$. In other words, $`G`$ consists of those homeomorphisms $`T`$ of $`X`$, for which there exists a bound on the number of coordinates of any point that are changed under the action of $`T`$, when $`(X,\sigma )`$ is represented as a subshift by means of some generator. Since the group $`G`$ is locally finite, the orbit equivalence relation $`X\times X`$ has a structure of AF-groupoid \[Re\]. Consider the groupoid C-algebra $`A=C^{}()`$ and the automorphism $`\alpha `$ of $`A`$ defined by $`\alpha (f)=f(\sigma \times \sigma )`$. The algebra $`C(X)`$ is a subalgebra of $`A`$, and there exists a unique conditional expectation $`E:AC(X)`$. Let $`C_0(X)`$ be the $``$-subalgebra of $`C(X)`$ spanned by characteristic functions of clopen sets, and $`C_0(X,)`$ the subalgebra of $`C_0(X)`$ consisting of real functions. Every element $`gG`$ defines a canonical unitary $`u_gA`$ such that $`u_gfu_g^{}=fg^1`$ for $`fC(X)`$. The $``$-algebra generated by $`C_0(X)`$ and $`u_g`$, $`gG`$, is our algebra $`𝒜`$ of local operators.
For $`HC_0(X,)`$ consider the $`1`$-cocycle $`c_HZ^1(,)`$,
$$c_H(x,y)=\underset{j}{}(H(\sigma ^jx)H(\sigma ^jy)).$$
Recall \[Re, Definition 3.15\] that a measure $`\mu `$ on $`X=^{(0)}`$ satisfies the $`(c_H,1)`$-KMS condition if its modular function is equal to $`e^{c_H}`$. In other words,
$$\frac{dg_{}\mu }{d\mu }(x)=e^{c_H(g^1x,x)}.$$
###### Proposition 5.1
Let $`HC_0(X,)`$. Then
(i) Any measure $`\mu `$ on $`X`$ which is an equilibrium measure at $`H`$ satisfies the $`(c_H,1)`$-KMS condition. In particular, any measure of maximal entropy is $`G`$-invariant.
(ii) The mapping $`\mu \mu E`$ defines a one-to-one correspondence between equilibrium measures on $`X`$ at $`H`$ and equilibrium states on $`C^{}()`$ at $`H`$.
Proof. First note that if $`\varphi `$ is an $`\alpha `$-invariant state on $`C^{}()`$, and $`\mu =\varphi |_{C(X)}`$ then
$$h_\mu (\sigma )=h_{\mu E}(\alpha )h_\varphi (\alpha ).$$
The equality is proved by standard arguments using \[CNT, Corollary VIII.8\]. The inequality follows from the fact that if $`\psi `$ is a state on a finite dimensional C-algebra $`M`$ with a masa $`B`$ then $`S(\psi )S(\psi |_B)`$. It follows that if $`\mu `$ is an equilibrium measure at $`H`$ then $`\mu E`$ is an equilibrium state at $`H`$, and if $`\varphi `$ is an equilibrium state at $`H`$ then $`\varphi |_{C(X)}`$ is an equilibrium measure at $`H`$. By Theorem 4.1 any equilibrium state is a $`(\sigma _t^H,1)`$-KMS state. But by \[Re, Proposition 5.4\] any $`(\sigma _t^H,1)`$-KMS state has the form $`\mu E`$ for some measure $`\mu `$ satisfying the $`(c_H,1)`$-KMS condition. From this both assertions of the proposition follow.
###### Example 5.2
As an application of Proposition 5.1 consider a topological Markov chain $`(X,\sigma )`$ with transition matrix $`A_T`$. As is well-known, if $`A_T`$ is primitive then the Perron-Frobenius theorem implies the uniqueness of the trace on $`C^{}()`$. If $`A_T`$ is only supposed to be irreducible, then the traces of $`C^{}()`$ form a simplex with the number of vertices equal to the index of cyclicity of the matrix. The barycenter of this simplex is the unique $`\alpha `$-invariant trace. By Proposition 5.1 we conclude that if $`A_T`$ is irreducible then $`(X,\sigma )`$ has a unique measure of maximal entropy. Thus we have recovered a well-known result of Parry (see \[W, Theorem 8.10\]).
While in the abelian case uniquely ergodic systems are of great interest, they are not so for asymptotically abelian systems with locality. Indeed, we have
###### Proposition 5.3
Let $`(A,\alpha )`$ be a C-dynamical system which is asymptotically abelian with locality. If there is a unique invariant state $`\tau `$, then $`\tau `$ is a trace, and $`\pi _\tau (A)`$ is an abelian algebra.
For later use the main part of the proof will be given in a separate lemma.
###### Lemma 5.4
Let $`(A,\alpha )`$ be an asymptotically abelian system with locality, $`\tau `$ an $`\alpha `$-invariant ergodic trace on $`A`$, $`H`$ a local self-adjoint operator. Suppose for each $`H^{}`$ in the real linear span of $`\alpha ^j(H)`$, $`j`$, and for each $`k`$ there exists an $`\alpha ^k`$-invariant $`(\sigma _t^{H^{},k},1)`$-KMS state $`\varphi `$ such that $`\tau =\frac{1}{k}_{j=0}^{k1}\varphi \alpha ^j`$. Then $`\pi _\tau (H)`$ is central in $`\pi _\tau (A)`$.
Proof. Replacing $`A`$ by $`A/\mathrm{Ker}\pi _\tau `$ we may identify $`A`$ with $`\pi _\tau (A)B(_\tau )`$.
The automorphism $`\alpha `$ being extended to $`A^{\prime \prime }B(_\tau )`$ is strongly asymptotically abelian. Hence for any $`k`$ the fixed point algebra $`(A^{\prime \prime })^{\alpha ^k}`$ is central. Since $`\alpha `$ is ergodic, this algebra is $`k_0`$-dimensional for some $`k_0|k`$, and we may enumerate its atoms $`z_1,\mathrm{},z_{k_0}`$ in such a way that $`\alpha (z_1)=z_2,\mathrm{},\alpha (z_{k_01})=z_{k_0},\alpha (z_{k_0})=z_1`$. Now if $`\varphi `$ is an $`\alpha ^k`$-invariant state such that $`\tau =\frac{1}{k}_{j=0}^{k1}\varphi \alpha ^j`$ then $`\varphi k\tau `$, hence $`\varphi (x)=\tau (xa)`$ for some positive $`a(A^{\prime \prime })^{\alpha ^k}`$. In particular, $`\varphi `$ is a trace. So if in addition $`\varphi `$ is a $`(\sigma _t^{H^{},k},1)`$-KMS state then the dynamics $`\sigma _t^{H^{},k}`$ is trivial on $`(A^{\prime \prime })_{s(a)}`$, where $`s(a)`$ is the support of $`a`$. Hence $`\delta _{H^{},k}(y)s(a)=0`$ for all local $`y`$. Then $`\delta _{H^{},k}(y)z_i=0`$ for some $`z_i`$ ($`1ik_0`$) majorized by $`s(a)`$.
Fix a local $`xA`$. Choose $`p`$ such that $`\alpha ^j(H)`$ commutes with $`x`$ whenever $`|j|p`$. Pick any $`m>p`$ and set $`k=2m+1`$. For $`\lambda ^k`$ consider the operator
$$H(\lambda )=\underset{j=m}{\overset{m}{}}\lambda _j\alpha ^j(H).$$
Applying the result of the previous paragraph to $`H^{}=H(\lambda )`$, we find $`i`$, $`1ik_0`$, such that $`\delta _{H(\lambda ),k}(y)z_i=0`$ for all local $`y`$. Denote by $`X_i`$ the set of all $`\lambda ^k`$ satisfying the latter condition. Since $`^k=_iX_i`$, there exists $`i`$ for which $`^k`$ coincides with the linear span of $`X_i`$. Without loss of generality we may suppose that $`i=1`$. Since for any $`j[m+p,mp]`$, any $`j^{}0`$, and any $`\lambda X_1`$, the elements $`\alpha ^{j^{}k}(H(\lambda ))`$ and $`\alpha ^j(x)`$ commute, we obtain
$$0=\delta _{H(\lambda ),k}(\alpha ^j(x))z_1=[H(\lambda ),\alpha ^j(x)]z_1,$$
hence $`[\alpha ^j^{}(H),\alpha ^j(x)]z_1=0`$ for $`j^{}[m,m]`$. In particular,
$$\alpha ^j([H,x])z_1=0\text{for}j[m+p,mp].$$
If $`k_0k`$ and $`k2p\frac{k}{2}(>k_0)`$ then $`_{j=m+p}^{mp}\alpha ^j(z_1)=1`$, so $`[H,x]=0`$. If $`k_0=k`$ then $`[H,x]z=0`$, where
$$z=\underset{j=m+p}{\overset{mp}{}}\alpha ^j(z_1)=\underset{j=m+p}{\overset{mp}{}}\alpha ^j(z_1),\tau (z)=\frac{k2p}{k}.$$
Since $`\frac{k2p}{k}1`$ as $`m\mathrm{}`$, we conclude that $`[H,x]=0`$.
Proof of Proposition 5.3. Since $`A`$ is a unital AF-algebra, there exists a trace on $`A`$, hence there exists an $`\alpha `$-invariant trace. It follows that the unique $`\alpha `$-invariant state is a trace.
If $`H`$ is local then for any subset $`I`$ of $`Z`$ there exists a $`(\sigma _t^{H,I},1)`$-KMS state. Indeed, if we take an increasing sequence of finite subsets $`I_n`$ of $`I`$ such that $`_nI_n=I`$, an increasing sequence of local algebras $`A_n`$ such that $`\alpha ^j(H)A_n`$ for $`jI_n`$ and $`_nA_n`$ is dense in $`A`$, and a sequence of states $`\varphi _n`$ such that $`\varphi _n|_{A_n}`$ is a $`(\sigma _t^{H,I_n},1)`$-KMS state, then any weak limit point of the sequence $`\{\varphi _n\}_n`$ will be a $`(\sigma _t^{H,I},1)`$-KMS state. If in addition $`I+k=I`$, then the state can be chosen to be $`\alpha ^k`$-invariant (since the set of $`(\sigma _t^{H,I},1)`$-KMS states is $`\alpha ^k`$-invariant). But if $`\varphi `$ is an $`\alpha ^k`$-invariant state then the state $`\frac{1}{k}_{j=0}^{k1}\varphi \alpha ^j`$ is $`\alpha `$-invariant, hence it coincides with $`\tau `$. Thus the conditions of Lemma 5.4 are satisfied. Hence $`\pi _\tau (H)`$ is central in $`\pi _\tau (A)`$ for any local $`H`$, so $`\pi _\tau (A)`$ is abelian.
We consider two examples illustrating Proposition 5.3.
###### Example 5.5
Let $`U`$ be the bilateral shift on a separable Hilbert space $``$, and $`\alpha =\mathrm{Ad}U|_A`$, where $`A`$ is the C-algebra $`K()+1`$, $`K()`$ being the algebra of compact operators. Then the only $`\alpha `$-invariant state is the trace $`\tau `$, which annihilates $`K()`$. Then $`\pi _\tau (A)=1`$.
###### Example 5.6
More generally, consider a uniquely ergodic system $`(X,\sigma )`$ and construct a system $`(C^{}(),\alpha )`$ as above. Let $`\tau `$ be an $`\alpha `$-invariant trace. Then $`\tau =\mu E`$ for some measure $`\mu `$, and the unique ergodicity of $`(X,\sigma )`$ means that $`\mu `$ is the unique invariant measure. We check the conditions of Lemma 5.4 for any $`HC(X_0,)`$.
By the same reasons as in the proof of Lemma 5.4, the fixed point algebra $`(\pi _\tau (A)^{\prime \prime })^\alpha `$ is central. By \[FM\] the center of the algebra $`\pi _\tau (A)^{\prime \prime }`$ is isomorphic to $`L^{\mathrm{}}(X,\mu )`$. Since the measure $`\mu `$ is ergodic, we conclude that the trace $`\tau `$ is also ergodic.
Let $`HC(X_0,)`$, $`k`$, and $`\varphi `$ any $`\alpha ^k`$-invariant $`(\sigma _t^{H,kZ},1)`$-KMS state. Then $`\varphi =\nu E`$ for some $`\sigma ^k`$-invariant measure $`\nu `$. Since $`\mu =\frac{1}{k}_{j=0}^{k1}\nu \sigma ^j`$, we have $`\tau =\frac{1}{k}_{j=0}^{k1}\varphi \alpha ^j`$.
Thus we can apply Lemma 5.4, and conclude that $`\pi _\tau (C(X))C(\mathrm{supp}\mu )`$ is central in $`\pi _\tau (A)`$. This means that $`G`$ acts trivially on $`\mathrm{supp}\mu `$, and $`\pi _\tau (A)=C(\mathrm{supp}\mu )`$. By \[Re, Proposition 4.5\] the kernel of $`\pi _\tau `$ is the algebra corresponding to the groupoid $`_{X\backslash \mathrm{supp}\mu }`$. Since there is no non-zero finite $`\sigma `$-invariant measures on $`X\backslash \mathrm{supp}\mu `$, any $`\alpha `$-invariant state is zero on $`\mathrm{Ker}\pi _\tau `$. Thus the system $`(A,\alpha )`$ is uniquely ergodic and $`\pi _\tau (A)=C(\mathrm{supp}\mu )`$.
We next give an example of an asymptotically abelian C-dynamical system $`(A,\alpha )`$ with $`A`$ an AF-algebra, for which there exist non-tracial $`\alpha `$-invariant states with maximal finite entropy. Hence the assumption of locality in Theorem 4.1 is essential.
###### Example 5.7
Let $``$ be an infinite-dimensional Hilbert space, $`A`$ the even CAR-algebra over $``$, $`\alpha `$ the Bogoliubov automorphism corresponding to a unitary $`U`$. It easy to see that $`\alpha `$ is asymptotically abelian if and only if $`(U^nf,g)\underset{n\mathrm{}}{}0`$ for any $`f,g`$. If in addition $`U`$ has singular spectrum then by the proof of \[SV, Theorem 5.2\] we have $`ht(\alpha )=0`$, while there are many non-tracial $`\alpha `$-invariant states (for example, quasi-free states corresponding to scalars $`\lambda (0,1/2)`$). Unitaries with such properties can be obtained using Riesz products. We shall briefly recall the construction.
Let $`q>3`$ be a real number, $`\{n_k\}_{k=1}^{\mathrm{}}`$ a sequence of positive integers such that $`\frac{n_{k+1}}{n_k}q`$, $`\{a_k\}_{k=1}^{\mathrm{}}`$ a sequence of real numbers such that $`a_k(1,1)`$, $`a_k0`$ as $`k\mathrm{}`$, $`_ka_k^2=\mathrm{}`$. Then the sequence of measures
$$\frac{1}{2\pi }\left[\underset{k=1}{\overset{n}{}}(1+a_k\mathrm{cos}n_kt)\right]dt$$
on $`[0,2\pi ]`$ converges weakly to a probability measure $`\mu `$ with Fourier coefficients
$$\widehat{\mu }(n)=\mu (e^{int})=\{\begin{array}{cc}\underset{k=1}{\overset{\mathrm{}}{}}\left(\frac{a_k}{2}\right)^{|\epsilon _k|},\text{if}n=\underset{k}{}\epsilon _kn_k\text{with}\epsilon _k\{1,0,1\},\hfill & \\ \text{ }\hfill & \\ 0,\text{otherwise}.\hfill & \end{array}$$
The measure $`\mu `$ is singular by \[Z, Theorem V.7.6\]. We see also that $`\widehat{\mu }(n)0`$ as $`|n|\mathrm{}`$. Thus the operator $`U`$ of multiplication by $`e^{int}`$ on $`L^2([0,2\pi ],d\mu )`$ has the desired properties.
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# 1 Introduction
## 1 Introduction
Neutral current deep inelastic lepton–nucleon scattering and single nucleon inclusive production in $`e^+e^{}`$ pair–annihilation are formally related by crossing the kinematic channels. Already before the advent of Quantum Chromodynamics (QCD) Drell, Levy, and Yan mentioned the possibility that the deep inelastic scattering structure functions at the one side and the nucleon fragmentation functions in $`e^+e^{}`$ pair–annihilation on the other side may be related by an analytic continuation from the $`t`$– to the $`s`$–channel. The hadronic tensors for the space–like process of deep inelastic scattering (DIS) and the time–like single nucleon inclusive reaction would therefore be related by
$$W_{\mu \nu }^{(S)}(q,p)=W_{\mu \nu }^{(T)}(q,p).$$
(1)
Here $`p`$ denotes the nucleon momentum and $`q`$ is the 4–momentum transfer to the hadronic system, with $`q^2<0`$ for deep inelastic scattering and $`q^2>0`$ for $`e^+e^{}`$–annihilation.
At that time the physical nucleons were considered as bound states built up of bare nucleons and pions in the context of the Yukawa theory. The interactions were described by Bethe–Salpeter or Faddeev–type equations and their generalization <sup>1</sup><sup>1</sup>1Later on massive vector meson ladder–models were studied in , where the crossing relation Eq. (1) was verified for the respective kernels. aiming at a perturbative description of the structure and fragmentation functions. In these theories neither infra–red nor collinear singularities are occurring. One may think off a general representation of the structure and fragmentation functions in terms of current–current expectation values. However it was already shown that for $`e^+e^{}`$-annihilation also diagrams of distinct connectedness appear (cf. also ) which are absent in DIS so that a proof of Eq. (1) at the non–perturbative level becomes very difficult. The relation could be established for the aforementioned ladder–models <sup>2</sup><sup>2</sup>2For a review on the early developments see .. Within QCD the picture changes. Here it turns out that a thorough perturbative description of the structure functions and fragmentation functions is not possible. However, perturbation theory may be used to describe the scaling violations of these functions at large values of $`|q^2|`$ where the running coupling constant is small. The QCD–improved parton model, moreover, exhibits a similarity with the approach by Drell, Levy, and Yan, since in the range where single parton states are dominating, i.e. for the contributions of lowest twist, the non–perturbative contributions factorize. One therefore may calculate the respective one–particle evolution kernels and study their behavior under the crossing from the $`t`$– to the $`s`$–channel. A further complication in the case of a vector–theory as QCD is the emergence of infrared and also collinear singularities which have an essential impact on the crossing because of the behavior of the kernels at $`x=1`$ <sup>3</sup><sup>3</sup>3Eq. (1) was originally postulated assuming that those terms are absent . See also the subsequent discussion in pointing out that the exponent of the structure functions $`(1x)^p`$ near $`x=1`$ needs not to be integer.. Here $`x`$ denotes the Bjorken scaling variable which will be defined differently for the timelike and spacelike region except for $`x=1`$ where both definitions lead to the same value for $`x`$. As a consequence of the Bloch–Nordsiek theorem the kernels become distribution–valued for $`x=1`$ .
In leading order for unpolarized scattering the crossing relations mentioned above were given in . Similar relations hold in the polarized case. In this order these kernels are nothing but the lowest order splitting functions $`P_{kl}^{(0)}(x)`$, which are obtained from the inverse Mellin transforms of the anomalous dimensions $`\gamma _{kl}^{N(0)}`$ .
It is the aim of the present paper to investigate the validity of the Drell–Levy–Yan (DLY) relation, if applied to perturbatively calculable partonic structure functions and quantities related to them up to the level of two–loop order. To establish this crossing relation between space– and time–like processes one has to study scheme–invariant quantities which are the physical evolution kernels for specific choices of observables as the unpolarized and polarized structure and fragmentation functions or derivatives of them w.r.t. $`q^2`$. Furthermore, conditions are derived for the transformation of the splitting and coefficient functions from the space– to the time–like case. For the coefficient functions we extend the discussion to the NNLO level. Other relations between the splitting functions such as supersymmetric relations and relations due to conformal symmetry were discussed elsewhere, cf. e.g. .
The paper is organized as follows. Basic relations for the deep inelastic structure and fragmentation functions are summarized in section 2. In section 3 scheme–invariant combinations of coefficient and splitting functions are constructed for the space– and time–like processes both for unpolarized and polarized deep inelastic reactions where we consider two principal examples. The Drell–Levy–Yan relation is studied in detail in section 4. We also comment on a relation by Gribov and Lipatov which emerged in the same context. Section 5 contains the conclusions. In the appendix we present the differences between to the space– and time–like coefficient functions at $`O(\alpha _s^2)`$ as well as the convolution relations which are needed for the investigation of the DLY–relation.
## 2 Structure Functions and Fragmentation Functions
Deep inelastic scattering (DIS) of a lepton ($`l`$) off a hadron target ($`P`$) is described by the process
$$l(k_1)+P(p)l(k_2)+\mathrm{`}X^{},q=k_1k_2,q^2=Q^2<0.$$
(2)
where $`\mathrm{`}X^{}`$ represents an inclusive final state. When a single gauge boson is exchanged between the incoming lepton and the hadron the above process factorizes into the leptonic part and the remaining hadronic part. In the case of forward scattering the scattering matrix element can be written in terms of the leptonic tensor $`L_{\mu \nu }`$ and the hadronic tensor $`W_{\mu \nu }`$ by
$$|M|^2=L^{\mu \nu }W_{\mu \nu }.$$
(3)
The hadronic tensor contains the unpolarized and polarized deep inelastic structure functions $`F_i`$ and $`g_i`$. If the process is mediated by photon only we have $`i=1,2`$ in both the polarized and unpolarized case. Notice that instead of $`F_1`$ we can also take the longitudinal structure function $`F_L`$. At asymptotic values of the kinematic variables structure functions only depend on $`Q^2`$ and the Bjorken scaling variable
$$x_B=\frac{Q^2}{2p.q},0x_B1.$$
(4)
In QCD the $`Q^2`$ dependence of the structure functions is only logarithmic and it accounts for the violation of scaling. In the context of the parton model the structure functions can be expressed in terms of quark and gluon densities and the corresponding spacelike coefficient functions $`C_{i,k}^{(S)}`$ ($`k=q,g`$) <sup>4</sup><sup>4</sup>4Similar relations hold for the polarized structure functions $`g_1(x,Q^2)`$ and $`g_2(x,Q^2)`$ on the level of twist 2.
$`F_i^{(S)}(x_B,Q^2)`$ $`=`$ $`x_B{\displaystyle \underset{j=1}{\overset{N_f}{}}}e_j^2{\displaystyle _{x_B}^1}{\displaystyle \frac{dz}{z}}[{\displaystyle \frac{1}{N_f}}f_q^S({\displaystyle \frac{x_B}{z}},\mu _f^2)C_{i,q}^{(S)S}(z,{\displaystyle \frac{Q^2}{\mu _f^2}})+f_g({\displaystyle \frac{x_B}{z}},\mu _f^2)`$
$`\times C_{i,g}^{(S)}(z,{\displaystyle \frac{Q^2}{\mu _f^2}})+f_{q_j}^{NS}({\displaystyle \frac{x_B}{z}},\mu _f^2)C_{i,q}^{(S)NS}(z,{\displaystyle \frac{Q^2}{\mu _f^2}})],`$
$`i=2,L,`$ (5)
where $`e_j`$ denotes the charge of the $`j`$th quark flavor and $`N_f`$ represents the number of light flavors. The scale $`\mu _f`$, appearing in the above equation, denotes the factorization scale which is introduced while removing the collinear singularities from the partonic structure functions. In addition one encounters a dependence on the renormalization scale $`\mu _r`$ which arises in the renormalization procedure. For convenience this scale is put equal to the factorization scale in the following. Notice that the structure functions $`F_i`$ and $`g_i`$ do not depend on these scales. However the parton densities and the coefficient functions, which do depend on these scales, satisfy renormalization group equations which will be shown below.
In Eq. (2) the index $`(S)`$ in the structure functions indicates the space–like nature of the process ($`q^2<0`$). Furthermore in Eq. (2) appear the singlet ‘S’ and non-singlet ‘NS’ combinations of parton densities which are defined by
$`f_q^S(z,\mu _f^2)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N_f}{}}}\left[f_{q_i}(z,\mu _f^2)+f_{\overline{q}_i}(z,\mu _f^2)\right],`$ (6)
and
$`f_{q_i}^{NS}(z,\mu _f^2)`$ $`=`$ $`f_{q_i}(z,\mu _f^2)+f_{\overline{q}_i}(z,\mu _f^2){\displaystyle \frac{1}{N_f}}f_q^S(z,\mu _f^2),`$ (7)
respectively. Corresponding formulae hold for polarized scattering. In this case the polarized parton densities and polarized coefficient functions are denoted by $`\mathrm{\Delta }f_k(z,\mu _f^2)`$ ($`k=q,g`$) and $`\mathrm{\Delta }C_{i,k}(z,Q^2/\mu _f^2)`$ ($`i=1,2`$).
Whereas in deep inelastic scattering the constituent structure of the nucleons is studied, hadroproduction at $`e^+e^{}`$ colliders provides us with information about the fragmentation process of these constituents into the hadrons. This information is contained in the fragmentation functions observed in the reaction
$`l(k_1)+\overline{l}(k_2)\overline{P}(p)+\mathrm{`}X^{},q=k_1+k_2,q^2Q^2>0,`$ (8)
where the symbols have he same meaning as in Eq. (2). These fragmentation functions are the analogues of the DIS structure functions Therefore in the QCD improved parton model these functions can be expressed in a similar way in terms of parton fragmentation densities $`D_k`$ ($`k=q,g`$) multiplied by timelike coefficient functions i.e.
$`F_i^{(T)}(x_E,Q^2)`$ $`=`$ $`x_E{\displaystyle \underset{j=1}{\overset{n_f}{}}}e_j^2{\displaystyle _{x_E}^1}{\displaystyle \frac{dz}{z}}[{\displaystyle \frac{1}{N_f}}D_q^S({\displaystyle \frac{x_E}{z}},\mu _f^2)C_{i,q}^{(T)S}(z,{\displaystyle \frac{Q^2}{\mu _f^2}})+D_g({\displaystyle \frac{x_E}{z}},\mu _f^2)`$
$`\times C_{i,g}^{(T)}(z,{\displaystyle \frac{Q^2}{\mu _f^2}})+D_{q_j}^{NS}({\displaystyle \frac{x_E}{z}},\mu _f^2)C_{i,q}^{(T)NS}(z,{\displaystyle \frac{Q^2}{\mu _f^2}})],`$
$`i=2,L,`$ (9)
where the corresponding scaling variable for the process in Eq. (9) is defined by
$$x_E=\frac{2p.q}{Q^2},0x_E1,$$
(10)
The symbol $`T`$ appearing within parentheses in Eq. (2) denotes that the fragmentation functions are measured in time–like processes. The scales $`\mu _f`$ and $`\mu _r`$ are defined in the same way as in Eq. (2) where like in DIS we set the renormalization scale equal to the factorization scale. Furthermore the definitions for the singlet and non-singlet parton fragmentation functions are the same as those for the parton densities given in Eqs. (6, 7). Similarly as in DIS one can also study the annihilation processes in Eq. (8) where the hadron $`P`$ is polarized. This entails the definition of the polarized fragmentation functions denoted by $`g_1^{(T)}`$ and $`g_2^{(T)}`$ for which one can present a similar formula as in Eq. (2). Very often one also encounters the transverse structure function which in the timelike and spacelike case is given by
$`F_1^{(R)}(x,Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{2x}}\left[F_2^{(R)}(x,Q^2)F_L^{(R)}(x,Q^2)\right],`$
$`\text{with}R=S(x=x_B)`$ $`R=T(x=x_E).`$ (11)
## 3 Scheme–invariant Combinations
In this section we give a short outline of the origin of the factorization scheme dependence of the anomalous dimensions (splitting functions) and the coefficient functions. We also show how this dependence disappears in the evolution of the structure functions w.r.t. the kinematic variable $`Q^2`$. The discussion below deals with the DIS structure functions but the conclusions also hold for the fragmentation functions. The partonic structure functions denoted by $`\widehat{}_{i,k}`$ ($`i=1,2,L`$, $`k=q,g`$), representing the QCD radiative corrections, contain various divergences. First these divergences have to be regularized for which the most convenient way is to choose the method of $`n`$–dimensional regularization. Using this method the singularities reveal themselves in the form of pole terms of the type $`1/ϵ^j`$, with $`n=4+ϵ`$, in the quantity $`\widehat{}_{i,k}`$. The infrared divergences cancel between virtual and bremsstrahlung contributions by virtue of the Bloch–Nordsieck theorem . Due to the Kinoshita-Lee-Nauenberg theorem all the final state mass singularities are canceled too since the DIS structure function is an inclusive quantity. Then one is left with only two types of singularities. The first one originates from the ultraviolet region. This type of singularities is removed via a redefinition of the parameters appearing in the QCD Lagrangian. An example is the coupling constant which becomes equal to $`\alpha _s(\mu _r^2)`$ where $`\mu _r`$ is the renormalization scale. After coupling constant renormalization the hadronic structure function can be written as follows
$$F_i(x,Q^2)=\underset{k=q,g}{}\left(_{ik}(\alpha _s(\mu _r^2),\frac{Q^2}{\mu ^2},\frac{\mu ^2}{\mu _r^2},ϵ)\widehat{f}_k\right)(x),$$
(12)
where the symbol $``$ denotes the Mellin–convolution defined by
$$(fg)(z)=_0^1𝑑z_1_0^1𝑑z_2f(z_1)g(z_2)\delta (zz_1z_2).$$
(13)
Furthermore $`\widehat{f}_k`$ is defined as the bare parton density which is scale independent and is an unphysical object because of the singular behavior of $`_{ik}`$. Notice that the latter depends on the scale $`\mu _r`$ and therefore on the renormalization scheme w.r.t. the coupling constant. The parameter $`\mu `$ originates from $`n`$–dimensional regularization because in this method the coupling constant gets a dimension. The second type of singularity originates from the collinear region which can be attributed to the vanishing mass of the initial state parton represented by either the (anti-) quark or the gluon. Hence the $`ϵ`$ in Eq. (12) represents the collinear singularities which are removed from the partonic structure function via mass factorization and transferred to a transition function $`\mathrm{\Gamma }_{lk}`$ as follows
$$\widehat{}_{ik}(z,\alpha _s(\mu _r^2),\frac{Q^2}{\mu ^2},\frac{\mu ^2}{\mu _r^2},ϵ)=\underset{l=q,g}{}\left(C_{i,l}(\alpha _s(\mu _r^2),\frac{Q^2}{\mu _f^2},\frac{\mu _f^2}{\mu _r^2})\mathrm{\Gamma }_{lk}(\alpha _s(\mu _r^2),\frac{\mu _f^2}{\mu ^2},\frac{\mu _f^2}{\mu _r^2},ϵ)\right)(z).$$
(14)
This procedure provides us with the coefficient function denoted by $`C_{i,l}`$. Substitution of Eq. (14) into Eq. (12) leads to the result
$$F_i(x,Q^2)=\underset{l=q,g}{}\left(C_{i,l}(\alpha _s(\mu _r^2),\frac{Q^2}{\mu _f^2},\frac{\mu _f^2}{\mu _r^2})f_l(\alpha _s(\mu _r^2),\frac{\mu _f^2}{\mu ^2},\frac{\mu _f^2}{\mu _r^2})\right)(x),$$
(15)
where the renormalized parton density is defined as
$$f_l(z,\alpha _s(\mu _r^2),\frac{\mu _f^2}{\mu ^2},\frac{\mu _f^2}{\mu _r^2})=\underset{k=q,g}{}\left(\mathrm{\Gamma }_{lk}(\alpha _s(\mu _r^2),\frac{\mu _f^2}{\mu ^2},\frac{\mu _f^2}{\mu _r^2},ϵ)\widehat{f}_k\right)(z).$$
(16)
Since the mass factorization can be carried out in various ways one is left with an additional scheme dependence which comes on top of the renormalization scheme dependence entering the coupling constant in Eq. (12). The former only shows up in the parton densities and the coefficient functions and it only disappears in specific combinations representing physical quantities. Hence physical quantities are invariant under scheme transformation. Like in the case of renormalization, mass factorization leads to the introduction of a scale $`\mu _f`$ called mass factorization scale which is related to the factorization scheme dependence. Like in the latter case $`\mu _f`$ drops out in physical quantities as the DIS structure functions or fragmentation functions. The change of the parton densities and the coefficient functions with respect to a variation in the scales $`\mu _r`$ and $`\mu _f`$ is determined by the renormalization group equation (RGE) The renormalization group equation of the parton densities follow from the one presented for the transition functions $`\mathrm{\Gamma }_{lk}`$. The latter takes the following form
$`\left(\left[\left\{\mu _f^2{\displaystyle \frac{}{\mu _f^2}}+\beta (a_s(\mu _f^2)){\displaystyle \frac{}{a_s(\mu _f^2)}}\right\}\mathrm{𝟏}\delta _{lm}{\displaystyle \frac{1}{2}}P_{lm}(a_s(\mu _f^2),ϵ)\right]\mathrm{\Gamma }_{mk}(a_s(\mu _f^2),{\displaystyle \frac{\mu _f^2}{\mu ^2}},1,ϵ)\right)(z)=0,`$ (17)
$`a_s(\mu _f^2){\displaystyle \frac{\alpha _s(\mu _f^2)}{4\pi }},\mathrm{𝟏}=\delta (1z),`$
where we have set $`\mu _r=\mu _f`$ for simplicity. The functions $`P_{ij}(a_s,ϵ,z)`$ appearing in the above equation are the splitting functions. Furthermore the beta-function is defined by
$$\mu _r^2\frac{da_s(\mu _r^2)}{d\mu _r^2}=\beta _0a_s^2(\mu _r^2)\beta _1a_s^3(\mu _r^2)\mathrm{},$$
(18)
The same equation as in Eq. (17) also applies to the parton density because of the definition in Eq. (16). The scale dependence of the coefficient function in Eq. (15) is given by
$`\left(\left[\left\{\mu _f^2{\displaystyle \frac{}{\mu _f^2}}+\beta (a_s(\mu _f^2)){\displaystyle \frac{}{a_s(\mu _f^2)}}\right\}\mathrm{𝟏}\delta _{lm}+{\displaystyle \frac{1}{2}}P_{lm}(a_s(\mu _f^2),ϵ)\right]C_{i,m}(a_s(\mu _f^2),{\displaystyle \frac{Q^2}{\mu _f^2}},1)\right)(z)=0.`$
As has been mentioned above scheme transformations such as
$$\mathrm{\Gamma }_{lk}\underset{m=q,g}{}Z_{lm}\overline{\mathrm{\Gamma }}_{mk},C_{i,l}\underset{m=q,g}{}\overline{C}_{i,m}Z_{ml}^1,$$
(20)
will not alter the physical observable like e.g. the structure functions and fragmentation functions. The relation between the splitting and coefficient functions computed in two different schemes is found to be
$`P_{lk}`$ $`=`$ $`{\displaystyle \underset{\{m,n\}=q,g}{}}Z_{lm}\overline{P}_{mn}(Z^1)_{nk}2\beta (a_s){\displaystyle \underset{m=q,g}{}}Z_{lm}{\displaystyle \frac{d}{da_s}}(Z^1)_{mk},`$ (21)
$`C_{i,l}`$ $`=`$ $`{\displaystyle \underset{m=q,g}{}}\overline{C}_{i,m}(Z^1)_{ml}.`$ (22)
Below we present the relation between the coefficient functions computed in two different schemes up to order $`a_s^2(Q^2)`$. Notice that we have chosen here $`\mu _f^2=Q^2`$ in order to get rid off the logarithms $`\mathrm{ln}(Q^2/\mu _f^2)`$ which usually appear. Up to $`O(a_s^2)`$ one obtains
$`C_{i,q}`$ $`=`$ $`\delta (1z)+a_s(\overline{C}_{i,q}^{(1)}+Z_{qq}^{(1)})+a_s^2(\overline{C}_{i,q}^{(2)}+Z_{qq}^{(2)}+(Z_{qq}^{(1)})^2+Z_{qg}^{(1)}Z_{gq}^{(1)}`$ (23)
$`+\overline{C}_{i,q}^{(1)}Z_{qq}^{(1)}+\overline{C}_{i,g}^{(1)}Z_{gq}^{(1)})+\mathrm{},`$
$`C_{i,g}`$ $`=`$ $`a_s(\overline{C}_{i,g}^{(1)}+Z_{qg}^{(1)})+a_s^2(\overline{C}_{i,g}^{(2)}+Z_{qg}^{(2)}+Z_{qg}^{(1)}(Z_{gg}^{(1)}+Z_{qq}^{(1)})`$ (24)
$`+\overline{C}_{i,q}^{(1)}Z_{qg}^{(1)}+\overline{C}_{i,g}^{(1)}Z_{gg}^{(1)})+\mathrm{}.`$
In the subsequent part of the paper it is much more convenient to derive the expressions in the Mellin transform representation so that one can avoid the convolution symbol $``$. The Mellin transform of a function $`f(z)`$ is given by
$$f^{(N)}=_0^1𝑑zz^{N1}f(z)$$
(25)
In this way Eq. (13) can be written as
$$(fg)^N=_0^1𝑑zz^{N1}(fg)(z)=f^Ng^N.$$
(26)
Since the structure functions are scheme independent they become renormalization group invariants. Hence they satisfy the RG equation
$$\left[\mu _f^2\frac{}{\mu _f^2}+\beta (a_s(\mu _f^2))\frac{}{a_s(\mu _f^2)}\right]F_i^N(x,Q^2)=0$$
(27)
The equation above follows from combining Eqs. (17, 3) and (15). However the independence of the structure function on the scales $`\mu _f`$ and $`\mu _r`$ is not manifest when multiplying parton densities with coefficient functions. In particular when the perturbation series of a physical quantity is computed up to finite order there is a residual dependence on these unphysical scales (see e.g. . Their influence is expected to become smaller when higher order terms in the perturbation series are taken into account . To avoid the problem of the factorization scheme dependence of the structure function when the perturbation series is truncated up to fixed order it is better to study evolution equations for the structure functions with respect to a physical scale which is represented by a kinematic variable like $`Q^2`$. In these type of evolution equations the kernels are factorization scheme independent order by order in perturbation theory. However the dependence on the choice of renormalization scheme and therefore the dependence on $`\mu _r`$ remains so that one is able to obtain a better estimate of the theoretical error on $`\alpha _s`$. For such an equation one needs two different structure functions called $`F_A(x,Q^2)`$ and $`F_B(x,Q^2)`$. Examples are $`A=2`$ and $`B=L`$ or $`F_A`$ and $`Q^2dF_A/dQ^2`$. Limiting ourselves to the singlet case, the evolution equation for the non-singlet structure functions is even more simple, one can write
$`F_I^N(Q^2)=f_q^N(a_s(\mu _f^2),{\displaystyle \frac{\mu _f^2}{Q_0^2}})C_{I,q}^N(a_s(\mu _f^2),{\displaystyle \frac{Q^2}{\mu _f^2}})+f_g^N(a_s(\mu _f^2),{\displaystyle \frac{\mu _f^2}{Q_0^2}})C_{I,g}^N(a_s(\mu _f^2),{\displaystyle \frac{Q^2}{\mu _f^2}}),`$
$`I=A,B,`$ (28)
Here one can view the $`C_{I,l}^N`$ ($`I=A,B`$, $`l=q,g`$) as matrix elements so that the equation above has the form
$$\left(\begin{array}{c}F_A^N\\ F_B^N\end{array}\right)=\left(\begin{array}{cc}C_{Aq}^N& C_{Ag}^N\\ C_{Bq}^N& C_{Bg}^N\end{array}\right)\left(\begin{array}{c}f_q^N\\ f_g^N\end{array}\right).$$
(29)
The coefficient functions satisfy the RG-equation in Eq. (3) and the solution is given by the T-ordered exponential
$`C_{I,l}^N(a_s(\mu _f^2),{\displaystyle \frac{Q^2}{\mu _f^2}})=C_{I,m}^N(a_s(Q^2),1)\left(T_{a_s}\left[exp\left\{{\displaystyle _{a_s(\mu _f^2)}^{a_s(Q^2)}}𝑑x{\displaystyle \frac{\gamma ^N(x)}{2\beta (x)}}\right\}\right]\right)_{ml},`$ (30)
where $`\gamma ^N`$ is the anomalous dimension matrix defined by
$$\gamma _{lk}^N=_0^1𝑑zz^{N1}P_{lk}(z).$$
(31)
We will now differentiate the coefficient functions w.r.t. $`Q^2`$
$`Q^2{\displaystyle \frac{C_{I,k}^N(a_s(\mu _f^2),Q^2/\mu _f^2)}{Q^2}}=\beta (a_s(Q^2)){\displaystyle \frac{C_{I,k}^N(a_s(\mu _f^2),Q^2/\mu _f^2)}{a_s(Q^2)}}=`$
$`[\beta (a_s(Q^2){\displaystyle \frac{C_{I,m}^N(a_s(Q^2),1)}{a_s(Q^2)}}\left(C^N\right)_{m,J}^1(a_s(Q^2),1)`$
$`{\displaystyle \frac{1}{2}}C_{I,m}^N(a_s(Q^2),1)\gamma _{mn}^N(a_s(Q^2))\left(C^N\right)_{n,J}^1(a_s(Q^2),1)]C^N_{J,k}(a_s(\mu _f^2),{\displaystyle \frac{Q^2}{\mu _f^2}}).`$ (32)
One can show that the expression above is invariant under scheme transformations. The latter are given by
$`\gamma _{lk}^N`$ $`=`$ $`{\displaystyle \underset{\{m,n\}=q,g}{}}Z_{lm}^N\overline{\gamma }_{mn}^N\left(Z^N\right)_{nk}^1+2\beta (a_s){\displaystyle \underset{m=q,g}{}}Z_{lm}^N{\displaystyle \frac{}{a_s}}\left(Z^N\right)_{mk}^1,`$ (33)
$`C_{I,l}^N`$ $`=`$ $`{\displaystyle \underset{m=q,g}{}}\overline{C}_{I,m}^N\left(Z^N\right)_{ml}^1,\left(C^N\right)_{l,I}^1={\displaystyle \underset{m=q,g}{}}Z_{lm}^N\left(\overline{C}^N\right)_{m,I}^1.`$ (34)
Since the $`Q^2`$-dependence only resides in the coefficient function the same evolution equation as in Eq. (3) also applies to $`F_I^N`$ in Eq. (3). For a short-hand notation we introduce the evolution variable $`t`$
$$t=\frac{2}{\beta _0}\mathrm{ln}\left(\frac{a_s(Q^2)}{a_s(Q_0^2)}\right),$$
(35)
so that we obtain
$$\frac{}{t}\left(\begin{array}{c}F_A^N\\ F_B^N\end{array}\right)=\frac{1}{4}\left(\begin{array}{cc}K_{AA}^N& K_{AB}^N\\ K_{BA}^N& K_{BB}^N\end{array}\right)\left(\begin{array}{c}F_A^N\\ F_B^N\end{array}\right),$$
(36)
where the physical (scheme invariant) kernel is given by
$`K_{IJ}^N=\left[4{\displaystyle \frac{C_{I,m}^N(t)}{t}}\left(C^N\right)_{m,J}^1(t){\displaystyle \frac{\beta _0a_s(Q^2)}{\beta (a_s(Q^2))}}C_{I,m}^N(t)\gamma _{mn}^N(t)\left(C^N\right)_{n,J}^1(t)\right].`$ (37)
The kernels $`K_{IJ}^N`$ depend both on the anomalous dimensions $`\gamma _{lk}^N`$ and the coefficient functions $`C_{I,l}^N(Q^2)`$ but the latter two quantities are combined in a factorization scheme independent way. This factorization scheme independence of $`K_{IJ}^N`$ holds order by order in perturbation theory. Using the series expansions for the anomalous dimensions and coefficient functions in terms of the strong coupling constant
$$\gamma _{lk}^N=\underset{n=0}{\overset{\mathrm{}}{}}a_s^{n+1}(Q^2)\left(\gamma ^N\right)_{lk}^{(n)},C_{I,l}^N(Q^2)=\underset{n=0}{\overset{\mathrm{}}{}}a_s^n(Q^2)\left(C^N\right)_{I,l}^{(n)},l,k=q,g,I=A,B,$$
(38)
one can compute order by order the coefficients in the perturbation series of the kernel
$`K_{IJ}^N={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}a_s^n(Q^2)\left(K^N\right)_{IJ}^{(n)}.`$ (39)
Notice that the coefficients $`\left(K^N\right)_{IJ}^{(n)}`$ are not invariant with respect to a finite renormalization of the coupling constant. This dependence is removed when the perturbation series in Eq. (39) is resummed in all orders.
### 3.1 $`𝑭_\mathrm{𝟐}\mathbf{(}𝒙\mathbf{,}𝑸^\mathrm{𝟐}\mathbf{)}`$ and $`𝑭_𝑳\mathbf{(}𝒙\mathbf{,}𝑸^\mathrm{𝟐}\mathbf{)}`$
Let us consider now two specific examples, choosing the structure functions $`F_2(x,Q^2)`$ and $`F_L(x,Q^2)`$ or the structure function $`F_2(x,Q^2)`$ and its slope $`F_2(x,Q^2)/t`$ as the observables $`F_{A,B}(x,Q^2)`$. in this combination of observables it is convenient to normalize the structure function $`F_L(x,Q^2)`$ to its gluonic contribution in lowest order. This is because $`F_L(x,Q^2)`$ vanishes in zeroth order of $`\alpha _s`$ due to the Callan–Gross relation, cf. Eq. (2), contrary to the structure function $`F_2(x,Q^2)`$. Therefore this normalization accounts for keeping the same order in the coupling constant for the two quantities
$$F_A^N(Q^2)=F_2^{N(S)}(Q^2),F_B^N(Q^2)=\frac{F_L^N(Q^2)}{a_s(Q^2)C_{L,g}^{N(1)}}.$$
(40)
Since both the coefficient functions $`C_{Lq}^{(1)}`$ and $`C_{Lg}^{(1)}`$ are scheme invariants such a normalization is possible. We expand now the kernels $`K_{IJ}^N`$ for this choice of observables into a series in $`a_s`$. The lowest order contribution is well-known, cf. e.g. ,
$$\begin{array}{ccccccc}\hfill K_{22}^{N(0)}& =& \gamma _{qq}^{N(0)}\frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}\gamma _{qg}^{N(0)}\hfill & & \hfill K_{2L}^{N(0)}& =& \gamma _{qg}^{N(0)}\hfill \\ & & & & & & \\ \hfill K_{L2}^{N(0)}& =& \gamma _{gq}^{N(0)}\left(\frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}\right)^2\gamma _{qg}^{N(0)}\hfill & & \hfill K_{LL}^{N(0)}& =& \gamma _{gg}^{N(0)}+\frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}\gamma _{qg}^{N(0)}\hfill \\ & & +\frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}\left(\gamma _{qq}^{N(0)}\gamma _{gg}^{N(0)}\right).\hfill & & & & \end{array}$$
To next-to-leading order in $`a_s(Q^2)`$, one finds
$`K_{22}^{N(1)}`$ $`=`$ $`\gamma _{qq}^{N(1)}{\displaystyle \frac{\beta _1}{\beta _0}}\gamma _{qq}^{N(0)}{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\left(\gamma _{qg}^{N(1)}{\displaystyle \frac{\beta _1}{\beta _0}}\gamma _{qg}^{N(0)}\right)+{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}C_{2,g}^{N(1)}\gamma _{qq}^{N(0)}`$ (41)
$`\left[{\displaystyle \frac{C_{L,q}^{N(2)}}{C_{L,g}^{N(1)}}}+\left({\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\right)^2C_{2,g}^{N(1)}{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}{\displaystyle \frac{C_{L,g}^{N(2)}}{C_{L,g}^{N(1)}}}\right]\gamma _{qg}^{N(0)}+C_{2,g}^{N(1)}\gamma _{gq}^{N(0)}`$
$`{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}C_{2,g}^{N(1)}\gamma _{gg}^{N(0)}+2\beta _0\left(C_{2,q}^{N(1)}{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}C_{2,g}^{N(1)}\right)`$
$`K_{2L}^{N(1)}`$ $`=`$ $`\gamma _{qg}^{N(1)}{\displaystyle \frac{\beta _1}{\beta _0}}\gamma _{qg}^{N(0)}C_{2,g}^{N(1)}(\gamma _{qq}^{N(0)}\gamma _{gg}^{N(0)})+2\beta _0C_{2,g}^{N(1)}`$ (42)
$`+\left(C_{2,q}^{N(1)}+{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}C_{2,g}^{N(1)}{\displaystyle \frac{C_{L,g}^{N(2)}}{C_{L,g}^{N(1)}}}\right)\gamma _{qg}^{N(0)}`$
$`K_{L2}^{N(1)}`$ $`=`$ $`\gamma _{gq}^{N(1)}{\displaystyle \frac{\beta _1}{\beta _0}}\gamma _{gq}^{N(0)}+{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\left(\gamma _{qq}^{N(1)}{\displaystyle \frac{\beta _1}{\beta _0}}\gamma _{qq}^{N(0)}\right)`$ (43)
$`\left({\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\right)^2\left(\gamma _{qg}^{N(1)}{\displaystyle \frac{\beta _1}{\beta _0}}\gamma _{qg}^{N(0)}\right){\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\left(\gamma _{gg}^{N(1)}{\displaystyle \frac{\beta _1}{\beta _0}}\gamma _{gg}^{N(0)}\right)`$
$`+\left[{\displaystyle \frac{C_{L,q}^{N(2)}}{C_{L,g}^{N(1)}}}{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}C_{2,q}^{N(1)}+\left({\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\right)^2C_{2,g}^{N(1)}\right]\gamma _{qq}^{N(0)}`$
$`[\left({\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\right)^3C_{2,g}^{N(1)}+2{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}{\displaystyle \frac{C_{L,q}^{N(2)}}{C_{L,g}^{N(1)}}}\left({\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\right)^2{\displaystyle \frac{C_{L,g}^{N(2)}}{C_{L,g}^{N(1)}}}`$
$`\left({\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\right)^2C_{2,q}^{N(1)}]\gamma _{qg}^{N(0)}+({\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}C_{2,g}^{N(1)}C_{2,q}^{N(1)}+{\displaystyle \frac{C_{L,g}^{N(2)}}{C_{L,g}^{N(1)}}})\gamma _{gq}^{N(0)}`$
$`\left[{\displaystyle \frac{C_{L,q}^{N(2)}}{C_{L,g}^{N(1)}}}+\left({\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\right)^2C_{2,g}^{N(1)}{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}C_{2,q}^{N(1)}\right]\gamma _{gg}^{N(0)}`$
$`+2\beta _0\left({\displaystyle \frac{C_{L,q}^{N(2)}}{C_{L,g}^{N(1)}}}{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}{\displaystyle \frac{C_{L,g}^{N(2)}}{C_{L,g}^{N(1)}}}\right)`$
$`K_{LL}^{N(1)}`$ $`=`$ $`\gamma _{gg}^{N(1)}{\displaystyle \frac{\beta _1}{\beta _0}}\gamma _{gg}^{N(0)}+{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\left(\gamma _{qg}^{N(1)}{\displaystyle \frac{\beta _1}{\beta _0}}\gamma _{qg}^{N(0)}\right)`$ (44)
$`{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}C_{2,g}^{N(1)}\gamma _{qq}^{N(0)}+\left[{\displaystyle \frac{C_{L,q}^{N(2)}}{C_{L,g}^{N(1)}}}{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}{\displaystyle \frac{C_{L,g}^{N(2)}}{C_{L,g}^{N(1)}}}+\left({\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}\right)^2C_{2,g}^{N(1)}\right]\gamma _{qg}^{N(0)}`$
$`C_{2,g}^{N(1)}\gamma _{gq}^{N(0)}+{\displaystyle \frac{C_{L,q}^{N(1)}}{C_{L,g}^{N(1)}}}C_{2,g}^{N(1)}\gamma _{gg}^{N(0)}+2\beta _0{\displaystyle \frac{C_{L,g}^{N(2)}}{C_{L,g}^{N(1)}}}`$
It is evident that the representation in terms of Mellin–moments is of advantage when compared to the corresponding $`x`$–space expressions. In the latter case one has to find the inverse Mellin transforms of quantities where $`C_{i,k}^N`$ and $`\gamma _{kl}^N`$ appear in the denominators of the expressions above which in general is not possible.
### 3.2 $`𝑭_\mathrm{𝟐}\mathbf{(}𝒙\mathbf{,}𝑸^\mathrm{𝟐}\mathbf{)}`$ and $`\mathbf{}𝑭_\mathrm{𝟐}\mathbf{(}𝒙\mathbf{,}𝑸^\mathrm{𝟐}\mathbf{)}\mathbf{/}\mathbf{}𝒕`$
A second example concerns the structure function $`F_2(x,Q^2)`$ and its slope. Both quantities are well measurable in the present–day deep inelastic scattering experiments. The observables $`F_{A,B}(x,Q^2)`$ are here
$$F_A^N(Q^2)=F_2^{(S)N}(Q^2),F_B^N(Q^2)=\frac{}{t}F_2^{(S)N}(Q^2)$$
(45)
This example has been considered before in . In leading order one obtains
$$\begin{array}{ccccccc}\hfill K_{22}^{N(0)}& =& 0\hfill & & \hfill K_{2d}^{N(0)}& =& 4\hfill \\ & & & & & & \\ \hfill K_{d2}^{N(0)}& =& \frac{1}{4}\left(\gamma _{qq}^{N(0)}\gamma _{gg}^{N(0)}\gamma _{qg}^{N(0)}\gamma _{gq}^{N(0)}\right)\hfill & & \hfill K_{dd}^{N(0)}& =& \gamma _{qq}^{N(0)}+\gamma _{gg}^{N(0)}.\hfill \end{array}$$
The next-to-leading order kernels read :
$`K_{22}^{N(1)}`$ $`=`$ $`0`$ (46)
$`K_{2d}^{N(1)}`$ $`=`$ $`0`$ (47)
$`K_{d2}^{N(1)}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[\gamma _{gg}^{N(0)}\gamma _{qq}^{N(1)}+\gamma _{gg}^{N(1)}\gamma _{qq}^{N(0)}\gamma _{qg}^{N(1)}\gamma _{gq}^{N(0)}\gamma _{qg}^{N(0)}\gamma _{gq}^{N(1)}\right]`$ (48)
$`{\displaystyle \frac{\beta _1}{2\beta _0}}\left(\gamma _{qq}^{N(0)}\gamma _{gg}^{N(0)}\gamma _{gq}^{N(0)}\gamma _{qg}^{N(0)}\right)+{\displaystyle \frac{\beta _0}{2}}C_{2,q}^{N(1)}\left(\gamma _{qq}^{N(0)}+\gamma _{gg}^{N(0)}2\beta _0\right)`$
$`{\displaystyle \frac{\beta _0}{2}}{\displaystyle \frac{C_{2,g}^{N(1)}}{\gamma _{qg}^{N(0)}}}\left[(\gamma _{qq}^{N(0)})^2\gamma _{qq}^{N(0)}\gamma _{gg}^{N(0)}+2\gamma _{qg}^{N(0)}\gamma _{gq}^{N(0)}2\beta _0\gamma _{qq}^{N(0)}\right]`$
$`{\displaystyle \frac{\beta _0}{2}}\left(\gamma _{qq}^{N(1)}{\displaystyle \frac{\gamma _{qq}^{N(0)}\gamma _{qg}^{N(1)}}{\gamma _{qg}^{N(0)}}}\right)`$
$`K_{dd}^{N(1)}`$ $`=`$ $`\gamma _{qq}^{N(1)}+\gamma _{gg}^{N(1)}{\displaystyle \frac{\beta _1}{\beta _0}}\left(\gamma _{qq}^{N(0)}+\gamma _{gg}^{N(0)}\right)`$ (49)
$`{\displaystyle \frac{2\beta _0}{\gamma _{qg}^{N(0)}}}\left[C_{2,g}^{N(1)}\left(\gamma _{qq}^{N(0)}\gamma _{gg}^{N(0)}2\beta _0\right)\gamma _{qg}^{N(1)}\right]+4\beta _0C_{2,q}^{N(1)}2\beta _1.`$
For this combination in next-to-leading order the evolution depends on two evolution kernels only. In the case of polarized deep inelastic scattering similar relations apply considering the structure function $`g_1(x,Q^2)`$ and its slope. The anomalous dimensions and coefficient functions of the unpolarized case have to be substituted by those for polarized scattering.
Although enforced by Eq. (36) one has still to show that the kernels (3.1-44, 3.249) are scheme–independent by an explicit calculation, which we have done using Eqs. (33,34) for the next-to-leading order contributions. In leading order the scheme–invariance is visible explicitly, since the leading order anomalous dimensions and the lowest order coefficient functions for $`F_L(x,Q^2)`$ are scheme invariants.
With the help of the evolution equation (36) we are now prepared to ask for the validity of crossing relations between different space– and time–like quantities in perturbative QCD. Relations of this kind are henceforth called Drell–Levy–Yan (DLY) relations, although the original reasoning of these and other authors was quite different. One condition to ask such a question at all is that the behavior of all contributing parts under crossing from space– to time–like momentum transfer are controlled. At large momentum transfers $`|q^2|`$ the single parton picture applies and the non–perturbative parton densities factorize. This makes it possible to study the respective evolutions kernels without reference to the non–perturbative input densities. Even if a crossing relation for these quantities does not exist, one still may investigate whether it exists for the perturbative evolution kernels. A further condition for this investigation is that the latter quantities are scheme–invariant, as in Eq. (36).
## 4 Drell-Levy-Yan relations
In the following we study in detail an interesting relation between deep inelastic lepton hadron scattering and $`e^+e^{}`$ annihilation into a hadron and anything else, proposed by Drell, Levy and Yan . Here we first briefly illustrate the idea behind the work of DLY for completeness. In field theory, the deep inelastic $`e^+e^{}`$ annihilation can be related to matrix elements of hadronic electromagnetic current operators similar to that of deep inelastic lepton–hadron scattering. The crucial difference, apart from the ones which originate from the kinematics, is that the annihilation process is $`not`$ related to the forward Compton amplitude contrary to deep inelastic scattering because in the former process the hadron is observed in the final state. Nevertheless, both processes are related by crossing symmetry which any field theory enjoys. This motivated DLY to study the process in detail and then relate it to the deep inelastic scattering process. From the structure of the hadronic tensors $`W_{\mu \nu }^S(q,p)`$ (space–like) and $`W_{\mu \nu }^T(q,p)`$ (time–like) and using the standard reduction formalism one can infer that
$$W_{\mu \nu }^T(q,p)=W_{\mu \nu }^S(q,p),$$
(50)
where the momenta within the respective parentheses of the above quantities are the same as those defined in the beginning of the paper. In the Bjorken limit for both deep inelastic scattering and deep inelastic annihilation for $`q^2=Q^2,p.q\mathrm{}`$ and $`q^2=Q^2,p.q\mathrm{}`$, respectively the scaling structure and fragmentation functions satisfy the following relation <sup>5</sup><sup>5</sup>5Here we indicate the overall signs in case of the scattering of particles of different spin, cf. . In the original work of DLY the Yukawa-theory was discussed which does not contain gauge bosons.:
$$F_i^{(S)}(x_B)=(1)^{2(s_1+s_2)}x_EF_i^{(T)}\left(\frac{1}{x_E}\right),i=1,2,L.$$
(51)
Here it has been assumed that non–perturbative input parton densities can be decoupled trivially and are the same. In other words, the functions $`F_i^{(T)}(x_E)`$ are the analytic continuations of the corresponding functions $`F_i^{(S)}(x_B)`$ from $`0<x_B1`$ to $`1x_E<\mathrm{}`$. This is true only when the continuation is smooth, i.e. if there are no singularities for example at $`x=1`$ etc. This relation is called DLY–relation in the literature.
In this section, we study this property in more detail extending earlier work . It is particularly interesting to study the above transformation at the level of the splitting functions and coefficient functions which constitute the physical quantities such as the structure and fragmentation functions. Then we show how these relations are preserved for the physical quantities by looking at the kernels discussed in the previous section. Apart from scaling violation one also encounters distributions of the type
$$\delta (1z),\left(\frac{\mathrm{ln}^i(1z)}{1z}\right)_+,i=0,1,2,\mathrm{},$$
(52)
which destroy the continuation through $`z=1`$. Here the distribution $`(\mathrm{ln}^i(1z)/(1z))_+`$ is represented by
$$\left(\frac{\mathrm{ln}^i(1z)}{1z}\right)_+=\delta (1z)\frac{\mathrm{ln}^{i+1}\delta }{(i+1)}+\theta (1\delta z)\frac{\mathrm{ln}^i(1z)}{(1z)},$$
(53)
where $`\delta 1`$. It turns out that the DLY–relation is violated for the coefficient functions and splitting functions separately because both are scheme dependent. This in particular happens when we adopt the $`\overline{\mathrm{MS}}`$-scheme. Here the relation is already violated up to one-loop order for the coefficient functions. Although one can choose other schemes in which Eq. (51) is preserved (see ) up to one-loop order we do not know whether this will hold up to any arbitrary order in perturbation theory.
Let us start with the simplest examples and consider the scheme–invariant evolution kernels Eq. (3.1, 3.2).
### 4.1 The Drell-Levy-Yan Relations at Leading Order
In the case of the scheme–independent evolution kernels describing the evolution of $`F_2(x,Q^2)`$ and $`F_2(x,Q^2)/t`$, respectively, or its polarized counterpart for the structure function $`g_1(x,Q^2)`$, only the transformation of two combinations of the leading order anomalous dimensions has to be considered, cf. (3.2). These are the determinant and the trace of the singlet anomalous dimension matrix at leading order. In both quantities the color factors of the off–diagonal elements enter only as a product. The unpolarized and polarized leading order splitting functions read
$`P_{qq}^{(0)}(z)=\mathrm{\Delta }P_{qq}^{(0)}(z)`$ $`=`$ $`4C_F\left[{\displaystyle \frac{1+z^2}{(1z)_+}}+{\displaystyle \frac{3}{2}}\delta (1z)\right]`$ (54)
$`P_{qg}^{(0)}(z)`$ $`=`$ $`8T_RN_f\left[z^2+(1z)^2\right]`$ (55)
$`\mathrm{\Delta }P_{qg}^{(0)}(z)`$ $`=`$ $`8T_RN_f\left[z^2(1z)^2\right]`$ (56)
$`P_{gq}^{(0)}(z)`$ $`=`$ $`4C_F{\displaystyle \frac{1+(1z)^2}{z}}`$ (57)
$`\mathrm{\Delta }P_{gq}^{(0)}(z)`$ $`=`$ $`4C_F{\displaystyle \frac{1(1z)^2}{z}}`$ (58)
$`P_{gg}^{(0)}(z)`$ $`=`$ $`8C_A\left[{\displaystyle \frac{z}{(1z)_+}}+{\displaystyle \frac{1z}{z}}+z(1z)\right]+2\beta _0\delta (1z)`$ (59)
$`\mathrm{\Delta }P_{gg}^{(0)}(z)`$ $`=`$ $`8C_A\left[{\displaystyle \frac{1}{(1z)_+}}+12z\right]+2\beta _0\delta (1z).`$ (60)
The crossing relations of the leading order splitting functions are
$$\begin{array}{ccccccc}\hfill \overline{P}_{qq}^{(0)}& =& zP_{qq}^{(0)}\left(\frac{1}{z}\right)\hfill & & \hfill \overline{P}_{qg}^{(0)}& =& \frac{C_F}{2N_fT_f}zP_{qg}^{(0)}\left(\frac{1}{z}\right)\hfill \\ & & & & & & \\ \hfill \overline{P}_{gq}^{(0)}& =& \frac{2N_fT_f}{C_F}zP_{gq}^{(0)}\left(\frac{1}{z}\right)\hfill & & \hfill \overline{P}_{gg}^{(0)}& =& zP_{gg}^{(0)}\left(\frac{1}{z}\right),\hfill \end{array}$$
where one demands
$`\delta (1z)\delta (1z).`$ (61)
Eq. (4.1) is easily verified and implies the validity of the crossing relation from space– to time–like evolution kernels Eq. (3.2), i.e. the validity of the DLY–relation for this case. For the second set of physical evolution kernels the DLY–relation follows at leading order referring to the transformation relations for the leading order longitudinal coefficient functions, Eqs. (74, 75) in an analogous way.
### 4.2 NLO Splitting function
As we know, the splitting functions and coefficient functions are not physical quantities due to their factorization scheme dependence. Hence, the naive continuation rule for these quantities may be violated, which is indeed the case in most of the schemes, e.g. in the $`\overline{\mathrm{MS}}`$ scheme characteristic of $`n`$-dimensional regularization. It was demonstrated by Curci, Furmanski and Petronzio that by an appropriate modification of the continuation rule in the $`\overline{\mathrm{MS}}`$ scheme one can show that the time–like splitting functions are related to their space–like counter parts. Since the modification of the continuation rule has to do with the scheme one adopts, it simply amounts to finding finite renormalization factors. It was shown that the finite renormalization factors can be constructed from the $`ϵ`$–dependent part of the splitting function when computed in dimensional regularization . In addition to this, care should be taken when dealing with quark and gluon states which was not the case in the work by DLY, where a color and flavor neutral field theory was discussed. The transformation rules are :
* The diagonal elements of the space–like flavor singlet splitting functions $`P_{qq},P_{gg}`$ have to be multiplied by $`(1)`$.
* The off–diagonal elements of the singlet splitting functions matrix have to be multiplied by $`C_F/(2N_fT_f)`$ for $`P_{qg}`$ and $`2N_fT_f/C_F`$ for $`P_{gq}`$, respectively, accounting for the interchange of the initial and final state particles under crossing.
Note that these transformations are automatically accounted for in the case of the leading order physical evolution kernels discussed in the previous paragraph. Keeping this in mind and using the known splitting functions and the continuation rules
$`\mathrm{ln}(1z)\mathrm{ln}(1z)\mathrm{ln}(z)+i\pi ,`$ (62)
$`\mathrm{ln}(\delta )\mathrm{ln}(\delta )+i\pi ,`$ (63)
one finds that
$`\overline{P}_{qq}^{(1)(S)}P_{qq}^{(1)(T)}`$ $`=`$ $`2\beta _0Z_{qq}^{(T)(1)}+Z_{qg}^{(T)(1)}\overline{P}_{gq}^{(0)}Z_{gq}^{(T)(1)}\overline{P}_{qg}^{(0)},`$ (64)
$`\overline{P}_{qg}^{(1)S}P_{gq}^{(1)(T)}`$ $`=`$ $`2\beta _0Z_{qg}^{(T)(1)}+Z_{qg}^{(T)(1)}(\overline{P}_{gg}^{(0)}\overline{P}_{qq}^{(0)})+\overline{P}_{qg}^{(0)}(Z_{qq}^{(T)(1)}Z_{gg}^{(T)(1)}),`$ (65)
$`\overline{P}_{gq}^{(1)S}P_{qg}^{(1)(T)}`$ $`=`$ $`2\beta _0Z_{gq}^{(T)(1)}+Z_{gq}^{(T)(1)}(\overline{P}_{qq}^{(0)}\overline{P}_{gg}^{(0)})+\overline{P}_{gq}^{(0)}(Z_{gg}^{(T)(1)}Z_{qq}^{(T)(1)}),`$ (66)
$`\overline{P}_{gg}^{(1)S}P_{gg}^{(1)(T)}`$ $`=`$ $`2\beta _0Z_{gg}^{(T)(1)}+Z_{gq}^{(T)(1)}\overline{P}_{qg}^{(0)}Z_{qg}^{(T)(1)}\overline{P}_{gq}^{(0)},`$ (67)
where the quantities with a bar denote that they are continued from $`z1/z`$ with the appropriate factors in front. These quantities read in explicit form :
$$\begin{array}{ccccccc}\hfill \overline{P}_{qq}^{(n)}(z)& =& zP_{qq}^{(n)}\left(\frac{1}{z}\right)\hfill & & \hfill \overline{P}_{qg}^{(n)}(z)& =& \frac{C_F}{2N_fT_f}zP_{qg}^{(n)}\left(\frac{1}{z}\right)\hfill \\ & & & & & & \\ \hfill \overline{P}_{gq}(z)^{(n)}& =& \frac{2N_fT_f}{C_F}zP_{gq}^{(n)}\left(\frac{1}{z}\right)\hfill & & \hfill \overline{P}_{gg}(z)^{(n)}& =& zP_{gg}^{(n)}\left(\frac{1}{z}\right).\hfill \end{array}$$
The relations given in Eqs. (644.2) remain true for the polarized splitting functions as well. The renormalization factors appearing in the Eqs. (6467) are given by
$`Z_{ij}^{T(1)}`$ $`=`$ $`P_{ji}^{(0)}\left(\mathrm{ln}(z)+a_{ji}\right).`$ (68)
The constants $`a_{ij}`$ are different in the unpolarized and polarized case. For unpolarized scattering they read
$$a_{qq}=a_{gg}=0,a_{qg}=\frac{1}{2},a_{gq}=\frac{1}{2},$$
(69)
whereas in the polarized case
$$a_{ij}=0.$$
(70)
The logarithms in the renormalization factors originate from the kinematics. In dimensional regularization, when one continues the partonic structure function $`\widehat{}_{i,k}`$ in Eq. (12) from the space–like to the time–like region one obtains an additional factor $`z^ϵ`$ which when multiplied with the pole in $`ϵ`$ yields $`\mathrm{ln}(z)`$. Since the pole is always associated with the splitting functions, one has the function $`P_{ij}^{(0)}`$ along with $`\mathrm{ln}(z)`$. The $`z`$–independent constant $`a_{ij}`$, which is also multiplied by the splitting function, results from the polarization average. For deep inelastic scattering one averages the processes with one gluon in the initial state by a factor $`1/(ϵ+2)`$. Such an average is not needed for the annihilation process since here the gluon appears in the final state. Notice that the average over the polarization sum does not show up in the polarized structure functions. Therefore in this case the constants $`a_{ij}`$ are zero.
The transformation behavior of the non–singlet splitting functions in NLO have been worked out in where also the relations for the NLO non–singlet coefficient functions were presented.
### 4.3 NLO Coefficient Functions
Now, let us study how space–like and time–like coefficient functions are related. The coefficient functions are expected to violate the DLY–relation due to their scheme dependence. Here we first present the relations between the space–like and time–like coefficient functions $`C_{i,k}(z)`$ ($`i=1,L;k=q,g`$). The leading order transverse coefficient functions are identical. At next-to-leading order, in the $`\overline{\mathrm{MS}}`$ scheme , the coefficient functions are related by the $`Z`$–factors in Eq. (68) as follows :
$`C_{1,q}^{(T)(1)}(z)+\left\{zC_{1,q}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)\right\}`$ $`=`$ $`Z_{qq}^{(T)(1)}`$ (71)
$`{\displaystyle \frac{1}{2}}\left[C_{1,g}^{(T)(1)}(z){\displaystyle \frac{C_F}{2N_fT_f}}\left\{2zC_{1,g}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)\right\}\right]`$ $`=`$ $`Z_{qg}^{(T)(1)}.`$ (72)
Since the coefficient functions depend on the hard scale of the process, one has to replace the space–like $`q^2`$ by the time–like $`q^2`$ in addition to Eqs. (6263). This leads to the following continuation rule
$`\mathrm{ln}\left({\displaystyle \frac{Q^2}{\mu _f^2}}\right)_{\mathrm{space}\mathrm{like}}\mathrm{ln}\left({\displaystyle \frac{Q^2}{\mu _f^2}}\right)_{\mathrm{time}\mathrm{like}}i\pi .`$ (73)
The $`Z`$–factors get contributions from two sources. The first one is $`z`$–dependent and comes from the phase space integrals. The time–like phase space acquires an extra factor $`z^ϵ`$ which gives a finite contribution when being multiplied with the pole terms $`1/ϵ`$. The pole term originates from the collinear divergence in $`n`$–dimensional regularization. The second term originates from the polarization average which is again absent in the time–like case. The continuation rules given in Eqs. (6261, 73) are essential to get the constant $`\zeta (2)`$ right when one goes from the space–like to the time–like region. Notice that the space–like coefficient function contains $`4\zeta (2)\delta (1z)`$ and the time–like one contains $`8\zeta (2)\delta (1z)`$. The difference which is $`12\zeta (2)`$ can be understood to originate from the one-loop vertex correction when one continues from the space–like to time–like region in $`Q^2`$. The same also holds when other regularization methods for the collinear divergences are chosen. It is worth noticing that if one would replace $`\mathrm{ln}(1z)\mathrm{ln}(1z)\mathrm{ln}(z)`$ contrary to the prescription in Eq. (62) one would obtain an additional term $`12\zeta (2)`$ on the righthand side of Eq. (71).
The zeroth order longitudinal coefficient functions are identically zero so that the first order contributions are scheme independent. This implies that there are no pole terms in the corresponding partonic structure function $`\widehat{}_{L,k}`$. Hence, there is no left–over finite piece which could arise from the $`z^ϵ`$– or $`n`$dimensional polarization average. We find
$`C_{L,q}^{(T)(1)}(z){\displaystyle \frac{z}{2}}C_{L,q}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)`$ $`=`$ $`0,`$ (74)
$`{\displaystyle \frac{1}{2}}\left[C_{L,g}^{(T)(1)}(z)+{\displaystyle \frac{C_F}{2N_fT_f}}\left\{zC_{L,g}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)\right\}\right]`$ $`=`$ $`0.`$ (75)
### 4.4 NNLO Coefficient Functions
#### 4.4.1 Longitudinal Coefficient Functions
We consider the NNLO correction to the longitudinal coefficient function. We follow the results given in for the space–like and for the time–like case. It turns out that the coefficient functions are related by the $`Z`$–factors through the matrix–valued convolutions
$`C_{L,q}^{(T)(2)}(z)+\left\{{\displaystyle \frac{z}{2}}C_{L,q}^{(S)(2)}\left({\displaystyle \frac{1}{z}}\right)\right\}`$ $`=`$ $`Z_{qq}^{(T)(1)}{\displaystyle \frac{z}{2}}C_{Lq}^{(1)(S)}\left({\displaystyle \frac{1}{z}}\right)`$ (76)
$`+Z_{gq}^{(T)(1)}{\displaystyle \frac{C_F}{2N_fT_f}}\left\{{\displaystyle \frac{z}{2}}C_{L,g}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)\right\},`$
$`{\displaystyle \frac{1}{2}}\left[C_{L,g}^{(T)(2)}(z)+{\displaystyle \frac{C_F}{2N_fT_f}}\left\{zC_{L,g}^{(S)(2)}\left({\displaystyle \frac{1}{z}}\right)\right\}\right]`$ $`=`$ $`Z_{qg}^{(T)(1)}{\displaystyle \frac{z}{2}}C_{L,q}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)`$ (77)
$`+Z_{gg}^{(T)(1)}{\displaystyle \frac{C_F}{2N_fT_f}}\left\{{\displaystyle \frac{z}{2}}C_{L,g}^{(1)S}\left({\displaystyle \frac{1}{z}}\right)\right\}.`$
The right hand side of the above equation contains the convolutions of $`Z`$–factors with the continued NLO longitudinal space–like coefficient functions. We have found this pattern by comparing the scheme transformation which we derived in the last section. The reason for this structure relies on the fact that $`C_{L,i}`$ is obtained as the difference between $`C_{2,i}`$ and $`C_{1,i}`$. Since the NLO coefficient functions involve various Nielsen–integrals, we used the following identities to simplify the expressions :
$`\mathrm{Li}_2\left({\displaystyle \frac{1}{z}}\right)`$ $`=`$ $`\mathrm{Li}_2(z){\displaystyle \frac{1}{2}}\mathrm{ln}^2(z)\zeta (2),`$ (78)
$`\mathrm{Li}_2\left(1{\displaystyle \frac{1}{z}}\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}^2(z)\mathrm{Li}_2(1z),`$ (79)
$`\mathrm{S}_{1,2}\left(1{\displaystyle \frac{1}{z}}\right)`$ $`=`$ $`{\displaystyle \frac{1}{6}}\mathrm{ln}^3(z)+\mathrm{S}_{1,2}(1z),`$ (80)
$`\mathrm{Li}_3\left(1{\displaystyle \frac{1}{z}}\right)`$ $`=`$ $`{\displaystyle \frac{1}{6}}\mathrm{ln}^3(z)+\mathrm{S}_{1,2}(1z)\mathrm{Li}_3(1z)+\mathrm{ln}(z)\mathrm{Li}_2(1z),`$ (81)
$`\mathrm{Li}_3\left({\displaystyle \frac{1}{z}}\right)`$ $`=`$ $`\mathrm{Li}_3(z)+{\displaystyle \frac{1}{6}}\mathrm{ln}^3(z)+\zeta (2)\mathrm{ln}(z),`$ (82)
$`\mathrm{S}_{1,2}\left({\displaystyle \frac{1}{z}}\right)`$ $`=`$ $`\mathrm{S}_{1,2}(z)+\mathrm{Li}_3(z)\mathrm{ln}(z)\mathrm{Li}_2(z){\displaystyle \frac{1}{6}}\mathrm{ln}^3(z)+\zeta (3).`$ (83)
If we do not continue $`\mathrm{ln}(1z)`$ and replace $`\mathrm{ln}(1z)\mathrm{ln}(1z)\mathrm{ln}(z)`$, then terms proportional to $`\zeta (2)`$ are not compensated between space–like and time–like coefficient functions and hence the relations given in Eqs. (76, 77) are no longer true. Although formally of NNLO, the coefficient functions $`C_{Lq(G)}^{(2)S,T}(z)`$ may be combined to physical evolution kernels together with the NLO splitting functions as shown in section 3.1. In Section 4.5. we will show that because of the transformation in Eqs. (76, 77) the physical evolution kernels in sections 3.1 and 3.2 remain DLY–invariant.
#### 4.4.2 Transverse Coefficient Functions
In NNLO physical evolution kernels for the transverse structure and fragmentation function can only be constructed when the space–like and time–like three-loop splitting functions are known. If they become available one can extend Eqs. (3.249) up to second order. Here we consider the relation between the space– and time–like coefficient functions using the transformation relations (61, 62, 63) and (73) for unpolarized and polarized scattering.
The space–like coefficient functions for unpolarized scattering are computed in whereas the time–like ones can be found in . The transverse coefficient functions are related by (see appendix 6.1) :
$`C_{1,q}^{(T)(2)}(z)+\left\{zC_{1,q}^{(S)(2)}\left({\displaystyle \frac{1}{z}}\right)\right\}`$ $`=`$ $`{\displaystyle \frac{1}{4}}[2(z)P_{qq}^{(1)S}\left({\displaystyle \frac{1}{z}}\right)+2\beta _0Z_{qq}^{(T)(1)}+Z_{gq}^{(T)(1)}\overline{P}_{qg}^{(0)}`$ (84)
$`Z_{qg}^{(T)(1)}\overline{P}_{gq}^{(0)}]\mathrm{ln}(z)+{\displaystyle \frac{1}{2}}Z^{(T)(1)}_{qq}Z^{(T)(1)}_{qq}`$
$`+Z_{qq}^{(T)(1)}\left(zC_{1,q}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)\right)+{\displaystyle \frac{1}{2}}Z_{gq}^{(T)(1)}Z_{qg}^{(T)(1)}`$
$`+Z_{gq}^{(T)(1)}\left({\displaystyle \frac{C_F}{2N_fT_f}}zC_{1,g}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)\right)+{\displaystyle \frac{1}{8}}\overline{P}_{qg}^{(0)}\overline{P}_{gq}^{(0)}.`$
$`+12C_F^2\zeta (2)\left(2\mathrm{ln}\left({\displaystyle \frac{Q^2}{\mu _f^2}}\right)3\right)^2\delta (1z)`$
For the polarized NNLO coefficient functions which were derived in and , we find that the form of Eqs. (84) is the same but the term
$`{\displaystyle \frac{1}{8}}\overline{P}_{gq}^{(0)}\overline{P}_{qg}^{(0)}`$
does not occur.
Similarly for the gluonic coefficient functions we find
$`{\displaystyle \frac{1}{2}}\left[C_{1,g}^{(T)(2)}(z){\displaystyle \frac{C_F}{2N_fT_f}}\left\{2zC_{1,g}^{(S)(2)}\left({\displaystyle \frac{1}{z}}\right)\right\}\right]`$ $`=`$ $`{\displaystyle \frac{1}{4}}[{\displaystyle \frac{C_F}{2N_fT_f}}2zP_{qg}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)+2\beta _0Z_{qg}^{(T)(1)}`$ (85)
$`+Z_{qg}^{(T)(1)}\overline{P}_{qq}^{(0)}Z_{qq}^{(T)(1)}\overline{P}_{qg}^{(0)}+Z_{gg}^{(T)(1)}\overline{P}_{qg}^{(0)}`$
$`Z_{qg}^{(T)(1)}\overline{P}_{gg}^{(0)}\left]\right(\mathrm{ln}(z)+{\displaystyle \frac{1}{2}})+{\displaystyle \frac{1}{2}}Z^{(T)(1)}_{qg}Z^{(T)(1)}_{qq}`$
$`+Z_{qg}^{(T)(1)}\left(zC_{1,q}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)\right)+{\displaystyle \frac{1}{2}}Z_{qg}^{(T)(1)}Z_{gg}^{(T)(1)}`$
$`+Z_{gg}^{(T)(1)}\left({\displaystyle \frac{C_F}{2N_fT_f}}zC_{1,g}^{(S)(1)}\left({\displaystyle \frac{1}{z}}\right)\right)`$
$`{\displaystyle \frac{1}{8}}\beta _0\overline{P}_{qg}^{(0)}+{\displaystyle \frac{1}{16}}\overline{P}_{qg}^{(0)}\left(\overline{P}_{gg}^{(0)}\overline{P}_{qq}^{(0)}\right).`$
For the polarized case, the terms
$`{\displaystyle \frac{1}{8}}\beta _0\overline{P}_{qg}^{(0)}+{\displaystyle \frac{1}{16}}\overline{P}_{qg}^{(0)}\left(\overline{P}_{gg}^{(0)}\overline{P}_{qq}^{(0)}\right)`$
in Eq. (85) are absent. Since we do not have to average over the initial state gluon polarization in the case of polarized scattering the term $`\mathrm{ln}(z)+1/2`$ multiplying the first bracket in Eq. (85) is replaced by $`\mathrm{ln}(z)`$, cf. also Eq. (68, 70).
### 4.5 NLO Physical Evolution Kernels
After having found the relations between space–like and time–like splitting and coefficient functions, we investigate the DLY-transformation for the physical evolution kernels presented in sections 3.1 and 3.2. In order to do this, we define the difference between the time–like quantities $`K_{IJ}^T`$ and the continued space–like quantities $`\overline{K}_{ij}^S`$ by
$$\delta K_{IJ}=K_{IJ}^T\overline{K}_{IJ}^S,$$
(86)
where $`\overline{K}_{IJ}^S`$ is obtained by transforming $`K_{IJ}^S`$ to the time–like region using the continuation rules (62, 63, 61, 73). Application of the DLY–transformations provides us with the following results
$`\delta K_{22}^{N(1)}`$ $`=`$ $`\delta \gamma _{qq}^{N(1)}{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}\delta \gamma _{gq}^{N(1)}+{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}\delta C_{2,g}^{N(1)}\overline{\gamma }_{qq}^{N(0)}`$ (87)
$`\left[{\displaystyle \frac{\delta C_{L,q}^{N(2)}}{\overline{C}_{L,g}^{N(1)}}}+\left({\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}\right)^2\delta C_{2,g}^{N(1)}{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}{\displaystyle \frac{\delta C_{L,g}^{N(2)}}{\overline{C}_{L,g}^{N(1)}}}\right]\overline{\gamma }_{qg}^{N(0)}`$
$`+\overline{\gamma }_{gq}^{N(0)}\delta C_{2,g}^{N(1)}{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}\overline{\gamma }_{gg}^{N(0)}\delta C_{2,g}^{N(1)}+2\beta _0\left(\delta C_{2,q}^{N(1)}{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}\delta C_{2,g}^{N(1)}\right)`$
$`=`$ $`\delta \gamma _{qq}^{N(1)}2\beta _0Z_{qq}^{(T)N(1)}\overline{\gamma }_{gq}^{N(0)}Z_{qg}^{(T)N(1)}+\overline{\gamma }_{qg}^{N(0)}Z_{gq}^{(T)N(1)}`$
$`+{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}(\delta \gamma _{gq}^{N(1)}+2\beta _0Z_{qg}^{(T)N(1)}Z_{qg}^{(T)N(1)}\overline{\gamma }_{qq}^{N(0)}+Z_{qq}^{(T)N(1)}\overline{\gamma }_{qg}^{N(0)}`$
$`Z_{gg}^{(T)N(1)}\overline{\gamma }_{qg}^{N(0)}+Z_{qg}^{(T)N(1)}\overline{\gamma }_{gg}^{N(0)}).`$
Substituting the expressions for $`\delta \gamma _{qq}^{N(1)}`$ and $`\delta \gamma _{gq}^{N(1)}`$ using Eqs. (64, 66), we get
$$\delta K_{22}^{N(1)}=0.$$
(88)
For the remaining evolution kernels one obtains
$`\delta K_{2L}^{N(1)}`$ $`=`$ $`\delta \gamma _{gq}^{N(1)}2\beta _0Z_{qg}^{(T)N(1)}+Z_{qg}^{(T)N(1)}(\overline{\gamma }_{qq}^{N(0)}\overline{\gamma }_{gg}^{N(0)})\overline{\gamma }_{qg}^{N(0)}(Z_{qq}^{(T)N(1)}Z_{gg}^{(T)N(1)}),`$ (89)
$`\delta K_{LL}^{N(1)}`$ $`=`$ $`\delta \gamma _{gg}^{N(1)}2\beta _0\overline{\gamma }_{gg}^{N(0)}+Z_{qg}^{(T)N(1)}\overline{\gamma }_{gq}^{N(0)}Z_{gq}^{(T)N(1)}\overline{\gamma }_{qg}^{N(0)}`$ (90)
$`+{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}[\delta \gamma _{gq}^{N(1)}2\beta _0Z_{qg}^{(T)N(1)}+Z_{qg}^{(T)N(1)}\overline{\gamma }_{qq}^{N(0)}Z_{qg}^{(T)N(1)}\overline{\gamma }_{gg}^{N(0)}Z_{qq}^{(T)N(1)}\overline{\gamma }_{qg}^{N(0)}`$
$`+Z_{gg}^{(T)N(1)}\overline{\gamma }_{qg}^{N(0)}],`$
$`\delta K_{L2}^{N(1)}`$ $`=`$ $`\delta \gamma _{qg}^{N(1)}+{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}\delta \gamma _{qq}^{N(1)}\left({\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}\right)^2\delta \gamma _{gq}^{N(1)}{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}\delta \gamma _{gg}^{N(1)}`$ (91)
$`2\beta _0Z_{gq}^{(T)N(1)}Z_{gq}^{(T)N(1)}\overline{\gamma }_{qq}^{N(0)}+Z_{qq}^{(T)N(1)}\overline{\gamma }_{gq}^{N(0)}Z_{gg}^{(T)N(1)}\overline{\gamma }_{gq}^{N(0)}+Z_{gq}^{(T)N(1)}\overline{\gamma }_{gg}^{N(0)}`$
$`+{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}\left[2\beta _0Z_{qq}^{(T)N(1)}Z_{qg}^{(T)N(1)}\overline{\gamma }_{gq}^{N(0)}+Z_{gq}^{(T)N(1)}\overline{\gamma }_{qg}^{N(0)}\right]`$
$`+\left({\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}\right)^2[2\beta _0Z_{qg}^{(T)N(1)}Z_{qg}^{(T)N(1)}\overline{\gamma }_{qq}^{N(0)}Z_{gg}^{(T)N(1)}\overline{\gamma }_{qg}^{N(0)}+Z_{qq}^{(T)N(1)}\overline{\gamma }_{qg}^{N(0)}`$
$`+Z_{qg}^{(T)N(1)}\overline{\gamma }_{gg}^{N(0)}]+{\displaystyle \frac{\overline{C}_{L,q}^{N(1)}}{\overline{C}_{L,g}^{N(1)}}}[2\beta _0Z_{gg}^{(T)N(1)}Z_{qg}^{(T)N(1)}\overline{\gamma }_{gq}^{N(0)}+Z_{gq}^{(T)N(1)}\overline{\gamma }_{qg}^{N(0)}].`$
The explicit expressions for the differences in the coefficient function are given in appendix 6.1 as well as a series of involved Mellin–convolutions leading to Nielsen–integrals (see appendix 6.2), which are necessary in the explicit calculation.
Using Eqs. (6467), leads to
$`\delta K_{L2}^{N(1)}`$ $`=`$ $`0,`$ (92)
$`\delta K_{2L}^{N(1)}`$ $`=`$ $`0,`$ (93)
$`\delta K_{LL}^{N(1)}`$ $`=`$ $`0.`$ (94)
The physical evolution kernels $`K_{I,J}`$ for the evolution of the structure functions $`F_2`$ and $`F_L`$ are thus DLY–invariant to next-to-leading order if continued from the space–like to the time–like region.
We turn now to the physical evolution kernels in next-to-leading order where we choose the physical quantities $`F_2`$, $`F_2/t`$ as a basis. Here only two evolution kernels are contributing, which change under the DLY-transformation as follows :
$`\delta K_{d2}`$ $`=`$ $`{\displaystyle \frac{\beta _0}{2}}\left(\delta C_{2q}^{N(1)}Z_{qq}^{(T)N(1)}\right)\left(\overline{\gamma }_{qq}^{N(0)}+\overline{\gamma }_{gg}^{N(0)}2\beta _0\right)`$ (95)
$`{\displaystyle \frac{\beta _0}{2\overline{\gamma }_{gq}^{N(0)}}}(\delta C_{2g}^{N(1)}Z_{qg}^{(T)N(1)})((\overline{\gamma }_{qq}^{N(0)})^2\overline{\gamma }_{qq}^{N(0)}\overline{\gamma }_{gg}^{N(0)}`$
$`+2\overline{\gamma }_{qg}^{N(0)}\overline{\gamma }_{gq}^{N(0)}2\beta _0\overline{\gamma }_{qq}^{N(0)})`$
$`\delta K_{dd}`$ $`=`$ $`2{\displaystyle \frac{\beta _0}{\overline{\gamma }_{gq}^{N(0)}}}\left(\delta C_{2g}^{N(1)}Z_{qg}^{(T)N(1)}\right)\left(\overline{\gamma }_{qq}^{N(0)}\overline{\gamma }_{gg}^{N(0)}2\beta _0\right)`$ (96)
$`+4\beta _0\left(\delta C_{1q}^{N(1)}Z_{qq}^{(T)N(1)}\right).`$
From Eqs. (71, 72, 74, 75) we can derive that
$`\delta K_{d2}`$ $`=`$ $`0,`$ (97)
$`\delta K_{dd}`$ $`=`$ $`0.`$ (98)
From these results it is clear that the time–like physical evolution kernels $`K_{ij}^T`$ can be directly derived from the space–like physical evolution kernels using the continuations in Eqs. (62, 63, 61, 73) where one has to account for the corresponding changes in the overall color factors. The $`Z^T`$–factors which are needed for the transformation of the splitting and coefficient functions cancel in the expression above. In the future one can extend the investigation performed in this section to physical evolution kernels at the NNLO-level, provided the 3–loop anomalous dimensions are calculated. For the choice of observables $`(F_2,F_L)`$ one also needs the three-loop coefficient functions.
We finally would like to comment on a relation derived by Gribov and Lipatov in for the leading order kernels for a pseudoscalar and a vector field theory. <sup>6</sup><sup>6</sup>6 See also for related work. One may write it in the form
$$\overline{K}(x_E,Q^2)=K(x_B,Q^2),$$
(99)
where $`\overline{K}`$ and $`K`$ denote the time– and space–like evolution kernels, respectively, and $`x_B=1/x_E`$. One verifies, that this relation holds in leading order for the space– and time–like splitting functions of QCD, Eqs. (5460), without changing the $`\delta `$–function, Eq. (61).
Starting with next-to-leading order, this relation is not preserved. For the physical non–singlet evolution kernels this was shown in and for some singlet combinations in . We find, that also for the physical singlet combinations, Eqs. (4144, 4649), this relation is violated as well.
## 5 Conclusions
The old question, whether the scattering cross sections of deep inelastic scattering $`e^{}+Pe^{}+\mathrm{`}X^{}`$ are related to the annihilation cross section $`e^++e^{}\overline{P}+\mathrm{`}X^{}`$ by a crossing relation changing from $`t`$– to $`s`$–channel was newly discussed. Since in both reactions non–perturbative quantities such as the structure and fragmentation functions contribute the above question cannot be answered by means of perturbation theory for the process as a whole. However, since both the parton densities involved in the space– and time–like process factorize if the virtuality $`Q^2=|q^2|`$ of the four–momentum transfer is large a related question can be asked for the crossing behavior of the respective evolution kernels, which are computable within perturbation theory. In the calculation of both inclusive processes only two types of singularities occur, the collinear singularity and the ultraviolet singularity. These divergences are absorbed into the bare parton densities and the coupling constant, respectively. Two distinct renormalization group equations are implied. They quantify the impact of the factorization and the renormalization scale on the DIS structure functions and fragmentation functions when the perturbation series is truncated up to a given order. However one can construct factorization–scale independent evolution kernels which describe the scheme–invariant evolution of these physical quantities in terms of a kinematic variable given by $`Q^2`$. This scheme invariant evolution is guaranteed up to any finite order in perturbation theory. Notice that in finite order this method does not remove the dependence of the physical quantities on the renormalization scheme of the strong coupling constant or its scale $`\mu _r`$.
The first example of the application of the physical evolution kernels is the coupled structure functions $`F_2(x,Q^2)`$ and $`F_L(x,Q^2)`$ associated with the corresponding fragmentation functions in $`e^+e^{}`$–annihilation. A second example is given by $`F_2(x,Q^2)`$ and $`F_2(x,Q^2)/\mathrm{ln}(Q^2)`$. Contrary to the splitting functions (anomalous dimensions) and coefficient functions the evolution kernels of the examples above are factorization scheme independent. For that purpose transformation relations have been derived for the splitting functions up to NLO and the coefficient functions up to NNLO. We have also shown that these kernels are invariant under the Drell–Levy–Yan–transformation up to next–to–leading order. On the other hand the Gribov–Lipatov relation, which is valid in leading order, is already violated at next-to-leading order. It remains to be seen how the physical evolution kernels behave under the DLY crossing relation at NNLO, which presumes the knowledge of the yet unknown three-loop splitting functions (space– and time–like) as well the three–loop longitudinal coefficient functions in the first example above.
Acknowledgment. We would like to thank P. Menotti for providing us a reprint of . Discussions with S. Kurth in an early phase of this work are acknowledged. This work was supported in part by EU contract FMRX-CT98-0194 (DG 12-MIHT).
## 6 Appendix
### 6.1 Coefficient Functions
In this appendix we list the difference of the space- and time–like coefficient functions in the $`\overline{\mathrm{MS}}`$ scheme, which are used in section 4 to study the validity of the DLY–relation. Here the expressions also contain the logarithms
$$L_{\mu _f}=\mathrm{ln}(Q^2/\mu _f^2)$$
(100)
which arise when the factorization scale $`\mu _f^2`$ is chosen to be different from $`Q^2`$.
The difference between the longitudinal non–singlet coefficient functions corresponding to the processes $`\gamma ^{}+qq+g+g`$ and $`\gamma ^{}\mathrm{`}\overline{q}^{}+q+g+g`$ respectively are given by
$`\delta C_{L,q}^{(2)NS}`$ $`=`$ $`C_F^2\left[4\left(2z\mathrm{ln}(z)\right)\mathrm{ln}(z)16\mathrm{L}\mathrm{i}_2(1z)+88z\right]`$ (101)
where $`\mathrm{`}q^{}`$ denotes the quark in the final state which undergoes fragmentation into a hadron $`P`$ (see Eq. (8)). The difference between the longitudinal purely singlet coefficient functions corresponding to the processes $`\gamma ^{}+qq+q+\overline{q}`$ and $`\gamma ^{}\mathrm{`}\overline{q}^{}+q+q+\overline{q}`$ respectively are given by
$`\delta C_{L,q}^{(2)PS}`$ $`=`$ $`N_fT_fC_F[8(6+4z{\displaystyle \frac{4}{3}}z^2)\mathrm{ln}(z)+16\mathrm{ln}^2(z)16`$ (102)
$`{\displaystyle \frac{304}{9z}}+64z{\displaystyle \frac{128}{9}}z^2]`$
The same is done for the longitudinal gluonic coefficient functions corresponding to the processes $`\gamma ^{}+gg+q+\overline{q}`$ and $`\gamma ^{}\mathrm{`}g^{}+g+q+\overline{q}`$, respectively. The difference in the coefficient function is given by
$`\delta C_{L,g}^{(2)}`$ $`=`$ $`C_F^2\left[8\left(1+{\displaystyle \frac{2}{z}}z\right)\mathrm{ln}(z)+8\mathrm{ln}^2(z)28+{\displaystyle \frac{24}{z}}+4z\right]`$ (103)
$`+C_AC_F[16(4{\displaystyle \frac{2}{z}}+z{\displaystyle \frac{1}{3}}z^2)\mathrm{ln}(z)16(1+{\displaystyle \frac{1}{z}})\mathrm{ln}^2(z)`$
$`+32(1{\displaystyle \frac{1}{z}})\mathrm{Li}_2(1z)8+{\displaystyle \frac{248}{9z}}24z+{\displaystyle \frac{40}{9}}z^2].`$
Notice that for the computation of the coefficients functions above and the ones following hereafter one also needs the virtual contributions to the zeroth and first order partonic processes. The differences between the transverse coefficient functions emerge from the same processes as mentioned above Eqs. (101, 102, 103). In the non–singlet case we have
$`\delta C_{1,q}^{(2)NS}`$ $`=`$ $`(C_F^2{\displaystyle \frac{1}{2}}C_AC_F)[8{\displaystyle \frac{\mathrm{ln}(z)}{1+z}}(2\zeta (2)4\mathrm{ln}(z)\mathrm{ln}(1+z)+\mathrm{ln}^2(z))`$ (104)
$`+4(2(1z)\zeta (2)+4(1z)+4(1z)\mathrm{ln}(z)\mathrm{ln}(1+z)`$
$`+2(1+z)\mathrm{ln}(z)(1z)\mathrm{ln}^2(z))\mathrm{ln}(z)16\mathrm{L}\mathrm{i}_2(z)\mathrm{ln}(z)`$
$`\times ({\displaystyle \frac{2}{1+z}}1+z)]`$
$`+N_fT_fC_F\left[{\displaystyle \frac{8}{9}}\left({\displaystyle \frac{10}{1z}}1+11z\right)\mathrm{ln}(z)\right]`$
$`+C_AC_F[4{\displaystyle \frac{\mathrm{ln}(z)}{1z}}({\displaystyle \frac{67}{9}}+\mathrm{ln}^2(z)2\zeta (2))+2({\displaystyle \frac{53}{9}}{\displaystyle \frac{187}{9}}z`$
$`+2(1+z)\mathrm{ln}(z)(1+z)\mathrm{ln}^2(z))\mathrm{ln}(z)+4\zeta (2)(1+z)\mathrm{ln}(z)]`$
$`+C_F^2[4{\displaystyle \frac{\mathrm{ln}(z)}{1z}}((8L_{\mu _f}6+4\mathrm{ln}(1z))\mathrm{ln}(1z)+6L_{\mu _f}18`$
$`+(4L_{\mu _f}+6{\displaystyle \frac{16}{3}}\mathrm{ln}(z))\mathrm{ln}(z))+4(4(1+z)L_{\mu _f}+1+5z`$
$`2(1+z)\mathrm{ln}(1z))\mathrm{ln}(z)\mathrm{ln}(1z)+2(2(5+z)L_{\mu _f}+14+40z`$
$`+2(1+z)\mathrm{ln}(z)\mathrm{ln}(1z)+6(1+z)L_{\mu _f}\mathrm{ln}(z)8(2+z)\mathrm{ln}(z)`$
$`+7(1+z)\mathrm{ln}^2(z))\mathrm{ln}(z)+2\mathrm{L}\mathrm{i}_2(1z)(4({\displaystyle \frac{12}{1z}}+7+7z)\mathrm{ln}(z)`$
$`10+22z)+8({\displaystyle \frac{24}{1z}}+13+13z)\mathrm{S}_{1,2}(1z)+36(1z)`$
$`+32\zeta (2)({\displaystyle \frac{2}{1z}}1z)\mathrm{ln}(z)+12\zeta (2)(912L_{\mu _f}+4L_{\mu _f}^2)\delta (1z)].`$
For the purely singlet difference we obtain
$`\delta C_{1,q}^{(2)PS}`$ $`=`$ $`N_fT_fC_F[(8(4+6z+{\displaystyle \frac{8}{3}}z^2)(L_{\mu _f}+\mathrm{ln}(1z))16(1{\displaystyle \frac{2}{3z}}+2z)`$ (105)
$`\times \mathrm{ln}(z)160{\displaystyle \frac{160}{9z}}112z{\displaystyle \frac{368}{9}}z^2+4(1+z)(4\mathrm{ln}(1z)+4L_{\mu _f}`$
$`+{\displaystyle \frac{10}{3}}\mathrm{ln}(z))\mathrm{ln}(z))\mathrm{ln}(z)8(1+{\displaystyle \frac{38}{9z}}z{\displaystyle \frac{38}{9}}z^2\left)\right(L_{\mu _f}+\mathrm{ln}(1z))`$
$`+16\left(2(1+z)\mathrm{ln}(z)23z{\displaystyle \frac{4}{3}}z^2\right)\mathrm{Li}_2(1z)+32(1+z)\mathrm{S}_{1,2}(1z)`$
$`{\displaystyle \frac{1168}{9}}{\displaystyle \frac{224}{27z}}+{\displaystyle \frac{640}{9}}z+{\displaystyle \frac{1808}{27}}z^2].`$
The difference between the gluonic coefficient functions equals
$`\delta C_{1,g}^{(2)}`$ $`=`$ $`2C_AC_F[(188+{\displaystyle \frac{704}{9z}}+66z+{\displaystyle \frac{184}{9}}z^2+(48{\displaystyle \frac{24}{z}}+16z+{\displaystyle \frac{32}{3}}z^2)L_{\mu _f}`$ (106)
$`+\left(4{\displaystyle \frac{100}{3z}}+10z\right)\mathrm{ln}(z)+\left(40{\displaystyle \frac{16}{z}}+12z+{\displaystyle \frac{32}{3}}z^2\right)\mathrm{ln}(1z)`$
$`+4\left(2+{\displaystyle \frac{2}{z}}+z\right)\mathrm{ln}(1+z)16\left(1+{\displaystyle \frac{1}{z}}+z\right)\mathrm{ln}(z)L_{\mu _f}`$
$`+8\left(2+{\displaystyle \frac{2}{z}}+z\right)\mathrm{ln}(z)\mathrm{ln}(1+z)16\left(1+{\displaystyle \frac{1}{z}}+z\right)\mathrm{ln}(z)\mathrm{ln}(1z)`$
$`({\displaystyle \frac{40}{3}}+{\displaystyle \frac{64}{3z}}+{\displaystyle \frac{52}{3}}z)\mathrm{ln}^2(z)+4(2{\displaystyle \frac{2}{z}}z)\mathrm{ln}^2(1z))\mathrm{ln}(z)`$
$`+\left(32+{\displaystyle \frac{356}{9z}}+4z{\displaystyle \frac{104}{9}}z^2\right)\mathrm{ln}(1z)+2\left(2{\displaystyle \frac{2}{z}}z\right)\mathrm{ln}^2(1z)`$
$`+(24+{\displaystyle \frac{284}{9z}}+4z{\displaystyle \frac{104}{9}}z^2)L_{\mu _f}+(32(1{\displaystyle \frac{3}{z}}2z)\mathrm{ln}(z)`$
$`+16(2{\displaystyle \frac{2}{z}}z)(\mathrm{ln}(1z)+L_{\mu _f})+16+{\displaystyle \frac{8}{z}}+16z+{\displaystyle \frac{32}{3}}z^2)\mathrm{Li}_2(1z)`$
$`+4\left(2+{\displaystyle \frac{2}{z}}+z\right)\left(1+2\mathrm{ln}(z)\right)\mathrm{Li}_2(z)+16\left(2+{\displaystyle \frac{2}{z}}+z\right)\mathrm{Li}_3(1z)`$
$`+8\left(8{\displaystyle \frac{16}{z}}10z\right)\mathrm{S}_{1,2}(1z){\displaystyle \frac{418}{9}}+{\displaystyle \frac{4168}{27z}}{\displaystyle \frac{668}{9}}z{\displaystyle \frac{856}{27}}z^2`$
$`+8\zeta (2)(3+{\displaystyle \frac{4}{z}}+2z)(1+2\mathrm{ln}(z))]`$
$`+2C_F^2[(42+{\displaystyle \frac{8}{z}}29z(1411z)\mathrm{ln}(z)+16(3{\displaystyle \frac{3}{z}}z)\mathrm{ln}(1z)`$
$`+4(2z)\mathrm{ln}(z)L_{\mu _f}4\left(2{\displaystyle \frac{4}{z}}z\right)\mathrm{ln}(z)\mathrm{ln}(1z)`$
$`16\left(2{\displaystyle \frac{2}{z}}z\right)\mathrm{ln}(1z)L_{\mu _f}+{\displaystyle \frac{10}{3}}(2z)\mathrm{ln}^2(z)`$
$`12(2{\displaystyle \frac{2}{z}}z)\mathrm{ln}^2(1z))\mathrm{ln}(z)+(64{\displaystyle \frac{52}{z}}18z`$
$`8(2{\displaystyle \frac{2}{z}}z)L_{\mu _f}6(2{\displaystyle \frac{2}{z}}z)\mathrm{ln}(1z))\mathrm{ln}(1z)`$
$`+2(16{\displaystyle \frac{10}{z}}3z)L_{\mu _f}+8((2z)\mathrm{ln}(z)2(2{\displaystyle \frac{2}{z}}z)(\mathrm{ln}(1z)`$
$`+L_{\mu _f})+6{\displaystyle \frac{6}{z}}3z)\mathrm{Li}_2(1z)+16(2{\displaystyle \frac{2}{z}}z)\mathrm{Li}_3(1z)`$
$`+8(10{\displaystyle \frac{8}{z}}5z)\mathrm{S}_{1,2}(1z)169+{\displaystyle \frac{106}{z}}+50z].`$
We have computed the same differences between the coefficient functions corresponding to the structure function $`g_1(x,Q^2)`$ which describes polarized scattering. The analogues of Eqs. (104, 105, 106) are given by
$`\delta \mathrm{\Delta }C_{1,q}^{(2)NS}`$ $`=`$ $`(C_F^2{\displaystyle \frac{1}{2}}C_AC_F)[8{\displaystyle \frac{\mathrm{ln}(z)}{1+z}}(2\zeta (2)4\mathrm{ln}(z)\mathrm{ln}(1+z)+\mathrm{ln}^2(z))`$ (107)
$`+4(2(1z)\zeta (2)+4(1z)+4(1z)\mathrm{ln}(z)\mathrm{ln}(1+z)`$
$`+2(1+z)\mathrm{ln}(z)(1z)\mathrm{ln}^2(z))\mathrm{ln}(z)16\mathrm{L}\mathrm{i}_2(z)\mathrm{ln}(z)`$
$`\times ({\displaystyle \frac{2}{1+z}}1+z)]`$
$`+N_fT_fC_F\left[{\displaystyle \frac{8}{9}}\left({\displaystyle \frac{10}{1z}}1+11z\right)\mathrm{ln}(z)\right]`$
$`+C_AC_F[4{\displaystyle \frac{\mathrm{ln}(z)}{1z}}({\displaystyle \frac{67}{9}}+\mathrm{ln}^2(z)2\zeta (2))+2({\displaystyle \frac{53}{9}}{\displaystyle \frac{187}{9}}z`$
$`+2(1+z)\mathrm{ln}(z)(1+z)\mathrm{ln}^2(z))\mathrm{ln}(z)+4\zeta (2)(1+z)\mathrm{ln}(z)]`$
$`+C_F^2[4{\displaystyle \frac{\mathrm{ln}(z)}{1z}}((8L_{\mu _f}6+4\mathrm{ln}(1z))\mathrm{ln}(1z)+6L_{\mu _f}18`$
$`+(4L_{\mu _f}+6{\displaystyle \frac{16}{3}}\mathrm{lg}(z))\mathrm{ln}(z))+4(4(1+z)L_{\mu _f}+1+5z`$
$`2(1+z)\mathrm{ln}(1z))\mathrm{ln}(z)\mathrm{ln}(1z)+2(2(5+z)L_{\mu _f}+18+36z`$
$`+2(1+z)\mathrm{ln}(z)\mathrm{ln}(1z)+6(1+z)L_{\mu _f}\mathrm{ln}(z)2(7+5z)\mathrm{ln}(z)`$
$`+7(1+z)\mathrm{ln}^2(z))\mathrm{ln}(z)+2\mathrm{L}\mathrm{i}_2(1z)(4({\displaystyle \frac{12}{1z}}+7+7z)\mathrm{ln}(z)`$
$`2+14z)+8({\displaystyle \frac{24}{1z}}+13+13z)\mathrm{S}_{1,2}(1z)+36(1z)`$
$`+32\zeta (2)({\displaystyle \frac{2}{1z}}1z)\mathrm{ln}(z)+12\zeta (2)(912L_{\mu _f}+4L_{\mu _f}^2)\delta (1z)]`$
$`\delta \mathrm{\Delta }C_{1,q}^{(2)PS}`$ $`=`$ $`N_fT_fC_F[(16(14z)(L_{\mu _f}+\mathrm{ln}(1z))184+24z16(2+3z)`$ (108)
$`\times \mathrm{ln}(z)+4(1+z)(4L_{\mu _f}+4\mathrm{ln}(1z)+{\displaystyle \frac{10}{3}}\mathrm{ln}(z))\mathrm{ln}(z))\mathrm{ln}(z)`$
$`48(1z)\left(L_{\mu _f}+\mathrm{ln}(1z)\right)+16\left(2(1+z)\mathrm{ln}(z)+14z\right)\mathrm{Li}_2(1z)`$
$`+32(1+z)\mathrm{S}_{1,2}(1z)112+112z]`$
$`\delta \mathrm{\Delta }C_{1,g}^{(2)}`$ $`=`$ $`2C_AC_F[(212+48z32(12z)L_{\mu _f}+4(8+5z)\mathrm{ln}(z)16(13z)`$ (109)
$`\times \mathrm{ln}(1z)8(4+z)\mathrm{ln}(z)L_{\mu _f}+8(2+z)\mathrm{ln}(z)\mathrm{ln}(1+z)`$
$`8(4+z)\mathrm{ln}(z)\mathrm{ln}(1z){\displaystyle \frac{4}{3}}(26+5z)\mathrm{ln}^2(z)4(2z)`$
$`\times \mathrm{ln}^2(1z))\mathrm{ln}(z)+32(1z)(\mathrm{ln}(1z)+L_{\mu _f})`$
$`+16\left(2+z(2z)\left(L_{\mu _f}+\mathrm{ln}(1z)\right)(8z)\mathrm{ln}(z)\right)\mathrm{Li}_2(1z)`$
$`+8(2+z)\mathrm{ln}(z)\mathrm{Li}_2(z)+16(2z)\mathrm{Li}_3(1z)32(5z)\mathrm{S}_{1,2}(1z)`$
$`+224224z+8\zeta (2)(83z)\mathrm{ln}(z)]`$
$`+2C_F^2[(22+40z+4(2z)L_{\mu _f}4(149z)\mathrm{ln}(1z)+4(85z)\mathrm{ln}(z)`$
$`+16(2z)\mathrm{ln}(1z)L_{\mu _f}4(2z)\mathrm{ln}(z)L_{\mu _f}`$
$`+4(2z)\mathrm{ln}(z)\mathrm{ln}(1z){\displaystyle \frac{10}{3}}(2z)\mathrm{ln}^2(z)`$
$`+12(2z)\mathrm{ln}^2(1z))\mathrm{ln}(z)8(1z)(\mathrm{ln}(1z)+L_{\mu _f})`$
$`+4\left(10+5z+4(2z)\left(\mathrm{ln}(1z)+L_{\mu _f}\right)2(2z)\mathrm{ln}(z)\right)`$
$`\times \mathrm{Li}_2(1z)16(2z)\mathrm{Li}_3(1z)40(2z)\mathrm{S}_{1,2}(1z)+9696z].`$
### 6.2 Convolutions
Here we list the convolutions of a series of functions, which are needed for the investigation of the DLY–relation in section 4. Using the definition in Eq. (13) we obtain
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}`$ $`=`$ $`{\displaystyle \frac{1}{(1z)}}[4\mathrm{S}_{1,2}(1z)2\mathrm{ln}(z)\mathrm{Li}_2(1z)`$ (110)
$`{\displaystyle \frac{1}{6}}\mathrm{ln}^3(z)]`$
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}{\displaystyle \frac{\mathrm{ln}(z)}{z}}`$ $`=`$ $`{\displaystyle \frac{1}{z}}\left[2\mathrm{S}_{1,2}(1z)+\mathrm{ln}(z)\mathrm{Li}_2(1z)\right]`$ (111)
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}z^2\mathrm{ln}(z)`$ $`=`$ $`{\displaystyle \frac{1}{12}}[324z+27z^224\mathrm{S}_{1,2}(1z)z^2`$ (112)
$`3(1+4z+5z^2)\mathrm{ln}(z)2z^2\mathrm{ln}^3(z)`$
$`12z^2\mathrm{ln}(z)\mathrm{Li}_2(1z)]`$
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}z\mathrm{ln}(z)`$ $`=`$ $`2+2z2z\mathrm{S}_{1,2}(1z)\mathrm{ln}(z)z\mathrm{ln}(z)`$ (113)
$`{\displaystyle \frac{z}{6}}\mathrm{ln}^3(z)z\mathrm{ln}(z)\mathrm{Li}_2(1z)`$
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}\mathrm{ln}(z)`$ $`=`$ $`2\mathrm{S}_{1,2}(1z){\displaystyle \frac{1}{6}}\mathrm{ln}^3(z)\mathrm{ln}(z)\mathrm{Li}_2(1z)`$ (114)
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}{\displaystyle \frac{\mathrm{ln}(1z)}{z}}`$ $`=`$ $`{\displaystyle \frac{1}{z}}\left[\mathrm{S}_{1,2}(1z)\mathrm{ln}(1z)\mathrm{Li}_2(1z)+\mathrm{Li}_3(1z)\right]`$ (115)
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}z^2\mathrm{ln}(1z)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\{45z+9z^28\mathrm{S}_{1,2}(1z)z^2+\mathrm{ln}(1z)`$ (116)
$`+4z\mathrm{ln}(1z)5z^2\mathrm{ln}(1z)3\mathrm{ln}(z)4z\mathrm{ln}(z)`$
$`+2\mathrm{ln}(1z)\mathrm{ln}(z)+4z\mathrm{ln}(1z)\mathrm{ln}(z)`$
$`2z^2\mathrm{ln}(1z)\mathrm{ln}^2(z)+[2+4z4z^2\mathrm{ln}(1z)`$
$`4z^2\mathrm{ln}(z)]\mathrm{Li}_2(1z)+4z^2\mathrm{Li}_3(1z)\}`$
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}z\mathrm{ln}(1z)`$ $`=`$ $`2+2z2z\mathrm{S}_{1,2}(1z)+\mathrm{ln}(1z)z\mathrm{ln}(1z)`$ (117)
$`\mathrm{ln}(z)+\mathrm{ln}(1z)\mathrm{ln}(z){\displaystyle \frac{z}{2}}\mathrm{ln}(1z)\mathrm{ln}^2(z)`$
$`\left[1+z\mathrm{ln}(1z)+z\mathrm{ln}(z)\right]\mathrm{Li}_2(1z)`$
$`+z\mathrm{Li}_3(1z)`$
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}\mathrm{ln}(1z)`$ $`=`$ $`2\mathrm{S}_{1,2}(1z){\displaystyle \frac{1}{2}}\mathrm{ln}(1z)\mathrm{ln}^2(z)`$ (118)
$`\left[\mathrm{ln}(1z)+\mathrm{ln}(z)\right]\mathrm{Li}_2(1z)+\mathrm{Li}_3(1z)`$
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_+`$ $`=`$ $`{\displaystyle \frac{1}{(1z)}}[{\displaystyle \frac{1}{2}}\mathrm{ln}(z)\mathrm{ln}^2(1z)2\mathrm{S}_{1,2}(1z)`$ (119)
$`\mathrm{ln}(z)\mathrm{Li}_2(1z){\displaystyle \frac{1}{2}}\mathrm{ln}^2(z)\mathrm{ln}(1z)]`$
$`{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}{\displaystyle \frac{1}{(1z)_+}}`$ $`=`$ $`{\displaystyle \frac{1}{(1z)}}\left[\mathrm{ln}(z)\mathrm{ln}(1z){\displaystyle \frac{1}{2}}\mathrm{ln}^2(z)\right]`$ (120)
$`{\displaystyle \frac{\mathrm{ln}(z)}{z}}{\displaystyle \frac{\mathrm{ln}(1z)}{z}}`$ $`=`$ $`{\displaystyle \frac{1}{2z}}[2\mathrm{S}_{1,2}(1z)\mathrm{ln}(1z)\mathrm{ln}^2(z)`$ (121)
$`2\mathrm{ln}(z)\mathrm{Li}_2(1z)]`$
$`{\displaystyle \frac{\mathrm{ln}(z)}{z}}\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_+`$ $`=`$ $`{\displaystyle \frac{1}{z}}[{\displaystyle \frac{1}{2}}\mathrm{ln}^2(1z)\mathrm{ln}(z)+\mathrm{ln}(1z)\mathrm{Li}_2(1z)`$ (122)
$`\mathrm{Li}_3(1z)]`$
$`{\displaystyle \frac{\mathrm{ln}(z)}{z}}{\displaystyle \frac{1}{(1z)_+}}`$ $`=`$ $`{\displaystyle \frac{1}{z}}[\mathrm{Li}_2(z)+\zeta (2)]`$ (123)
$`z^2\mathrm{ln}(z)\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_+`$ $`=`$ $`z^2[{\displaystyle \frac{3}{4}}\mathrm{S}_{1,2}(1z)+{\displaystyle \frac{3}{4z}}+{\displaystyle \frac{5}{4}}\mathrm{ln}(1z)`$ (124)
$`{\displaystyle \frac{1}{4z^2}}\mathrm{ln}(1z){\displaystyle \frac{1}{z}}\mathrm{ln}(1z){\displaystyle \frac{3}{4}}\mathrm{ln}(z)`$
$`{\displaystyle \frac{3}{2}}\mathrm{ln}(1z)\mathrm{ln}(z)+{\displaystyle \frac{1}{2}}\mathrm{ln}^2(1z)\mathrm{ln}(z)`$
$`+{\displaystyle \frac{3}{4}}\mathrm{ln}^2(z){\displaystyle \frac{1}{2}}\mathrm{ln}(1z)\mathrm{ln}^2(z){\displaystyle \frac{3}{2}}\mathrm{Li}_2(1z)`$
$`+\mathrm{ln}(1z)\mathrm{Li}_2(1z)\mathrm{ln}(z)\mathrm{Li}_2(1z)`$
$`\mathrm{Li}_3(1z)]`$
$`z^2\mathrm{ln}(z)z^2\mathrm{ln}(1z)`$ $`=`$ $`{\displaystyle \frac{z^2}{2}}(2\mathrm{S}_{1,2}(1z)+\mathrm{ln}(1z)\mathrm{ln}^2(z)`$ (125)
$`+2\mathrm{ln}(z)\mathrm{Li}_2(1z))`$
$`z^2\mathrm{ln}(z){\displaystyle \frac{1}{(1z)_+}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}z+{\displaystyle \frac{5}{4}}z^2{\displaystyle \frac{3}{2}}z^2\mathrm{ln}(z)`$ (126)
$`{\displaystyle \frac{z^2}{2}}\mathrm{ln}^2(z)+z^2\mathrm{ln}(z)\mathrm{ln}(1z)+z^2\mathrm{Li}_2(1z)`$
$`z\mathrm{ln}(z)\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_+`$ $`=`$ $`z[\mathrm{S}_{1,2}(1z)+\mathrm{ln}(1z){\displaystyle \frac{1}{z}}\mathrm{ln}(1z)\mathrm{ln}(z)`$ (127)
$`+{\displaystyle \frac{1}{2}}\mathrm{ln}^2(1z)\mathrm{ln}(z)\mathrm{ln}(z)\mathrm{ln}(1z)+{\displaystyle \frac{1}{2}}\mathrm{ln}^2(z)`$
$`{\displaystyle \frac{1}{2}}\mathrm{ln}(1z)\mathrm{ln}^2(z)\mathrm{Li}_2(1z)+\mathrm{ln}(1z)`$
$`\times \mathrm{Li}_2(1z)\mathrm{ln}(z)\mathrm{Li}_2(1z)\mathrm{Li}_3(1z)]`$
$`z\mathrm{ln}(z){\displaystyle \frac{1}{(1z)_+}}`$ $`=`$ $`z\mathrm{ln}(z)1+z{\displaystyle \frac{z}{2}}\mathrm{ln}^2(z)+z\mathrm{Li}_2(1z)`$ (128)
$`+z\mathrm{ln}(z)\mathrm{ln}(1z)`$
$`z\mathrm{ln}(z)z\mathrm{ln}(1z)`$ $`=`$ $`z\mathrm{S}_{1,2}(1z){\displaystyle \frac{z}{2}}\mathrm{ln}^2(z)\mathrm{ln}(1z)`$ (129)
$`z\mathrm{ln}(z)\mathrm{Li}_2(1z)`$
$`\mathrm{ln}(z)\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_+`$ $`=`$ $`\mathrm{S}_{1,2}(1z)+{\displaystyle \frac{1}{2}}\mathrm{ln}^2(1z)\mathrm{ln}(z)`$ (130)
$`{\displaystyle \frac{1}{2}}\mathrm{ln}(1z)\mathrm{ln}^2(z)+\mathrm{ln}(1z)\mathrm{Li}_2(1z)`$
$`\mathrm{ln}(z)\mathrm{Li}_2(1z)\mathrm{Li}_3(1z)`$
$`\mathrm{ln}(z)\mathrm{ln}(1z)`$ $`=`$ $`\mathrm{S}_{1,2}(1z){\displaystyle \frac{1}{2}}\mathrm{ln}^2(z)\mathrm{ln}(1z)`$ (131)
$`\mathrm{ln}(z)\mathrm{Li}_2(1z)`$
$`\mathrm{ln}(z){\displaystyle \frac{1}{(1z)_+}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}^2(z)+\mathrm{Li}_2(1z)+\mathrm{ln}(z)\mathrm{ln}(1z)`$ (132)
$`{\displaystyle \frac{1}{(1z)_+}}{\displaystyle \frac{1}{(1z)_+}}`$ $`=`$ $`2{\displaystyle \frac{\mathrm{ln}(1z)}{(1z)}}{\displaystyle \frac{\mathrm{ln}(z)}{(1z)}}\delta (1z)\zeta (2)`$ (133)
$`{\displaystyle \frac{1}{z}}\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_+`$ $`=`$ $`{\displaystyle \frac{1}{2z}}\mathrm{ln}^2(1z)`$ (134)
$`{\displaystyle \frac{1}{z}}{\displaystyle \frac{1}{(1z)_+}}`$ $`=`$ $`{\displaystyle \frac{1}{z}}\mathrm{ln}(1z)`$ (135)
$`z^2\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_+`$ $`=`$ $`\left[{\displaystyle \frac{1}{2}}+z{\displaystyle \frac{3}{2}}z^2\right]\mathrm{ln}(1z){\displaystyle \frac{1}{2}}z(1z)`$ (136)
$`+{\displaystyle \frac{3}{2}}z^2\mathrm{ln}^2(z)+{\displaystyle \frac{1}{2}}z^2\mathrm{ln}^2(1z)z^2\mathrm{ln}(z)\mathrm{ln}(1z)`$
$`z^2\mathrm{Li}_2(1z)`$
$`z^2{\displaystyle \frac{1}{(1z)_+}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+z{\displaystyle \frac{3}{2}}z^2z^2\mathrm{ln}(z)+z^2\mathrm{ln}(1z)`$ (137)
$`z\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_+`$ $`=`$ $`(1z)\mathrm{ln}(1z)+z\mathrm{ln}(z)z\mathrm{ln}(z)\mathrm{ln}(1z)`$ (138)
$`z\mathrm{Li}_2(1z)+{\displaystyle \frac{1}{2}}z\mathrm{ln}^2(1z)`$
$`z{\displaystyle \frac{1}{(1z)_+}}`$ $`=`$ $`1z+z\mathrm{ln}(1z)z\mathrm{ln}(z)`$ (139)
$`1\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_+`$ $`=`$ $`\mathrm{Li}_2(1z)\mathrm{ln}(z)\mathrm{ln}(1z)+{\displaystyle \frac{1}{2}}\mathrm{ln}^2(1z)`$ (140)
$`1{\displaystyle \frac{1}{(1z)_+}}`$ $`=`$ $`\mathrm{ln}(1z)\mathrm{ln}(z)`$ (141)
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# Reduction of the two-body dynamics to a one-body description in classical electrodynamics
## I Introduction
Recently, a novel approach to the two-body problem in general relativity has been introduced . The main motivation of that investigation rests on better understanding the late dynamical evolution of a coalescing binary system made of compact bodies of comparable masses, such as black holes and/or neutron stars. In fact, these astrophysical systems are among the most promising candidate sources for the detection of gravitational-waves with the future terrestrial interferometers such as the Laser Interferometric Gravitational Wave Observatory (LIGO) and Virgo. The basic idea pursued in , in part inspired by some results obtained in quantum electrodynamics , was to map the conservative two-body dynamics (henceforth denoted as the “real” dynamics) onto an effective one-body one, where a test particle moves in an effective external metric. As long as radiation reaction effects are not taken into account, the effective metric is just a deformation of the Schwarzschild metric with deformation parameter $`\nu =\mu /M`$, where $`\mu `$ is the reduced mass of the binary system and $`M`$ its total mass. The “effective” description should be viewed as a way of re-summing in a non-perturbative manner the badly convergent post-Newtonian-expanded dynamics of the “real” description. The results in were restricted to the second post-Newtonian level (2PN) and the analysis was mainly focused on the conservative part of the dynamics. More recently, a feasible way of incorporating radiation reaction effects has been proposed and the extension of the aforesaid approach to 3PN order has been investigated .
The purpose of the present paper is to test the robustness of the basic idea underlying the mapping of the two-body problem onto an effective one-body one, by applying it to classical electrodynamics. We limit to the conservative part of the dynamics of the bound states of two charged particles, up to second post-Coulombian order (2PC), and we take into account recoil effects. We investigate the possibility of describing the exchange of energies between the two bodies in the “real” problem through an “effective” auxiliary description, where a test particle moves in some external effective electromagnetic field. Generically, we expect that this electromagnetic field will be a deformation of the Coulomb potential with deformation parameter $`\nu =\mu /M`$, where $`\mu `$ is the usual reduced mass of the two charged particles and $`M`$ the total mass of the system. We shall see that the matching is also possible introducing in the effective description either a $`\nu `$-dependent vector potential or a deformed flat metric with deformation parameter $`\nu `$.
As already mentioned, the idea of reducing the relativistic two-body dynamics onto a relativistic one-body one was originally introduced in quantum electrodynamics. In particular, in the authors, taking into account recoil effects, resummed in the eikonal approximation the “crossed-ladder” Feynman diagrams for the scattering of two relativistic particles and mapped the one-body relativistic Balmer formula onto the two-body relativistic one. This method gives the correct quantum energy levels at least up to 1PC order, but some of the centrifugal barrier effects have to be added by hand. Todorov et al. developed a more systematic approach, based on the Lyppmann-Schwinger quasi-potential equation, which also gives correct results for the quantum energy levels, including the main parts of the radiative effects of the Lamb shift. Nevertheless, this last approach rests on some choices for the quasi-potential equation which are not very well justified and introduces in the effective description various energy-dependent quantities. In the following, whenever it is possible, we will compare our results in classical electrodynamics with the previous analysis for the corresponding quantum problem. Finally, note that, the aim of this paper is not to obtain new results with respect to the quantum energy-levels of the bound states of a two-body charged system, which is well known to be a hard problem . On the other hand, the present work wants to investigate, in the context of classical electrodynamics, the basic idea of reducing the two-body dynamics onto a one-body one, recently introduced in general relativity .
The outline of the paper is as follows. In Section II we review the relativistic two-body problem up to 2PC order and summarize its dynamics in a coordinate-invariant manner evaluating , within the Hamilton-Jacobi framework, the “energy-levels” of the bound states. In Section III we introduce the “effective” one-body description and define the “rules” needed to map the “real” onto the “effective” problem. Then, in Sections III A, III B and III C we analyze three feasible manners of implementing the matching. Finally, Section IV summarizes our main conclusions.
## II Two-body dynamics up to second post-Coulombian order
It was realized long ago that, in relativistic dynamics, if the position variables that are used to describe a system of charged interacting particles are the coordinates associated to a Lorentz frame <sup>*</sup><sup>*</sup>*For coordinates belonging to a Lorentz frame we mean coordinates which transform as linear representation of the Poincarè group ., then all higher time derivatives must appear in the Lagrangian . To get an “ordinary” Lagrangian it is necessary to introduce canonical position variables different from the Lorentz coordinates . At 2PC order the acceleration dependent Lagrangian was originally derived by Golubenkov and Smorodinskii . If one eliminates in that Lagrangian the higher time derivatives by using the equation of motion of lower orders then, as pointed out in , one does not obtain the correct equations of motion in a Lorentz frame. To eliminate correctly the accelerations one can use the method of “redefinition of position variables”, introduced by Damour and Schäfer in , which consists in appealing to a contact transformation induced by a change of coordinates from the Wheeler-Feynman coordinate system (Lorentz frame) to a well defined asymptotically inertial frame . More explicitly, the acceleration dependent Lagrangian at 2PC order is given by :
$$\stackrel{~}{}(𝒛_1,𝒛_2,𝒗_1,𝒗_2,𝒂_1,𝒂_2)=\stackrel{~}{}_0+\frac{1}{c^2}\stackrel{~}{}_2+\frac{1}{c^4}\stackrel{~}{}_4,$$
(1)
with
$`\stackrel{~}{}_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}m_1𝒗_1^2+{\displaystyle \frac{1}{2}}m_2𝒗_2^2{\displaystyle \frac{e_1e_2}{R}},`$ (2)
$`\stackrel{~}{}_1`$ $`=`$ $`{\displaystyle \frac{1}{8}}m_1𝒗_1^4+{\displaystyle \frac{1}{8}}m_2𝒗_2^4+{\displaystyle \frac{e_1e_2}{2R}}\left[𝒗_1𝒗_2+(\stackrel{~}{𝒏}𝒗_1)(\stackrel{~}{𝒏}𝒗_2)\right],`$ (3)
$`\stackrel{~}{}_4`$ $`=`$ $`{\displaystyle \frac{1}{16}}m_1𝒗_1^6+{\displaystyle \frac{1}{16}}m_2𝒗_2^6{\displaystyle \frac{e_1e_2}{8}}\{R[3(𝒂_1𝒂_2)(\stackrel{~}{𝒏}𝒂_1)(\stackrel{~}{𝒏}𝒂_2)]+2[(𝒗_1𝒂_2)(\stackrel{~}{𝒏}𝒗_1)`$ (6)
$`(𝒗_2𝒂_1)(\stackrel{~}{𝒏}𝒗_2)]+(\stackrel{~}{𝒏}𝒂_1)[𝒗_2^2(\stackrel{~}{𝒏}𝒗_2)^2](\stackrel{~}{𝒏}𝒂_2)[𝒗_1^2(\stackrel{~}{𝒏}𝒗_1)^2]`$
$`+{\displaystyle \frac{1}{R}}[𝒗_1^2𝒗_2^22(𝒗_1𝒗_2)^2𝒗_1^2(\stackrel{~}{𝒏}𝒗_2)^2𝒗_2^2(\stackrel{~}{𝒏}𝒗_1)^2+3(\stackrel{~}{𝒏}𝒗_1)^2(\stackrel{~}{𝒏}𝒗_2)^2]\},`$
where $`𝑹=𝒛_1𝒛_2`$, $`\stackrel{~}{𝒏}=𝑹/R`$, $`𝒗_i=\dot{𝒛}_i`$ and $`𝒂_i=\dot{𝒗}_i`$. In , Damour and Schäfer after having critically discussed and clarified the various results previously derived in the literature , worked out the contact transformations,
$`𝒒_1=𝒛_1{\displaystyle \frac{1}{c^4}}{\displaystyle \frac{e_1e_2}{4m_1}}\left\{(\stackrel{~}{𝒏}𝒗_2)𝒗_2+\stackrel{~}{𝒏}\left[{\displaystyle \frac{1}{2}}((\stackrel{~}{𝒏}𝒗_2)^2𝒗_2^2)+{\displaystyle \frac{e_1e_2}{m_2R}}\right]\right\},`$ (7)
$`𝒒_2=𝒛_2+{\displaystyle \frac{1}{c^4}}{\displaystyle \frac{e_1e_2}{4m_2}}\left\{(\stackrel{~}{𝒏}𝒗_1)𝒗_1+\stackrel{~}{𝒏}\left[{\displaystyle \frac{1}{2}}((\stackrel{~}{𝒏}𝒗_1)^2𝒗_1^2)+{\displaystyle \frac{e_1e_2}{m_1R}}\right]\right\},`$ (8)
which allow to eliminate the accelerations appearing in Eqs. (2)–(6). Hence, the final acceleration independent Lagrangian at 2PC order is given by :
$$(𝒒_1,𝒒_2,\dot{𝒒}_1,\dot{𝒒}_2)=_0+\frac{1}{c^2}_2+\frac{1}{c^4}_4,$$
(9)
with
$`_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}m_1\dot{𝒒}_1^2+{\displaystyle \frac{1}{2}}m_2\dot{𝒒}_2^2{\displaystyle \frac{e_1e_2}{q}},`$ (10)
$`_2`$ $`=`$ $`{\displaystyle \frac{1}{8}}m_1\dot{𝒒}_1^4+{\displaystyle \frac{1}{8}}m_2\dot{𝒒}_2^4+{\displaystyle \frac{e_1e_2}{2q}}\left[\dot{𝒒}_1\dot{𝒒}_2+(𝒏\dot{𝒒}_1)(𝒏\dot{𝒒}_2)\right],`$ (11)
$`_4`$ $`=`$ $`{\displaystyle \frac{1}{16}}m_1\dot{𝒒}_1^6+{\displaystyle \frac{1}{16}}m_2\dot{𝒒}_2^6{\displaystyle \frac{e_1e_2}{8q}}\{\dot{𝒒}_1^2\dot{𝒒}_2^22(\dot{𝒒}_1\dot{𝒒}_2)^2+3(𝒏\dot{𝒒}_1)^2(𝒏\dot{𝒒}_2)^2`$ (14)
$`(𝒏\dot{𝒒}_1)^2\dot{𝒒}_2^2(𝒏\dot{𝒒}_2)^2\dot{𝒒}_1^2+{\displaystyle \frac{e_1e_2}{m_2q}}\left[\dot{𝒒}_1^23(𝒏\dot{𝒒}_1)^2\right]`$
$`+{\displaystyle \frac{e_1e_2}{m_1q}}[\dot{𝒒}_2^23(𝒏\dot{𝒒}_2)^2]{\displaystyle \frac{2(e_1e_2)^2}{m_1m_2q^2}}\},`$
where $`𝒒=𝒒_1𝒒_2`$ and $`𝒏=𝒒/q`$. Applying the Legendre transformation to $``$, we derive (in full agreement with )
$$(𝒒_1,𝒒_2,𝒑_1,𝒑_2)=_0+\frac{1}{c^2}_2+\frac{1}{c^4}_4,$$
(15)
where
$`_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{𝒑_1^2}{m_1}}+{\displaystyle \frac{𝒑_2^2}{m_2}}\right)+{\displaystyle \frac{e_1e_2}{q}},`$ (16)
$`_2`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left({\displaystyle \frac{𝒑_1^4}{m_1^3}}+{\displaystyle \frac{𝒑_2^4}{m_2^3}}\right){\displaystyle \frac{e_1e_2}{2m_1m_2q}}\left[𝒑_1𝒑_2+(𝒏𝒑_1)(𝒏𝒑_2)\right],`$ (17)
$`_4`$ $`=`$ $`{\displaystyle \frac{1}{16}}({\displaystyle \frac{𝒑_1^6}{m_1^5}}+{\displaystyle \frac{𝒑_2^6}{m_2^5}})+{\displaystyle \frac{e_1e_2}{m_1m_2q}}\{{\displaystyle \frac{3(𝒏𝒑_1)^2(𝒏𝒑_2)^2}{8m_1m_2}}{\displaystyle \frac{𝒑_1^2(𝒏𝒑_2)^2}{8m_1m_2}}{\displaystyle \frac{𝒑_2^2(𝒏𝒑_1)^2}{8m_1m_2}}`$ (20)
$`+{\displaystyle \frac{1}{4}}\left[(𝒏𝒑_1)(𝒏𝒑_2)+(𝒑_1𝒑_2)\right]\left({\displaystyle \frac{𝒑_1^2}{m_1^2}}+{\displaystyle \frac{𝒑_2^2}{m_2^2}}\right){\displaystyle \frac{(𝒑_1𝒑_2)^2}{4m_1m_2}}+{\displaystyle \frac{𝒑_1^2𝒑_2^2}{8m_1m_2}}`$
$`+{\displaystyle \frac{e_1e_2}{q}}({\displaystyle \frac{𝒑_1^2}{m_1}}+{\displaystyle \frac{𝒑_2^2}{m_2}}){\displaystyle \frac{(e_1e_2)^2}{4q^2}}\}.`$
Let us denote
$$M=m_1+m_2,\mu =\frac{m_1m_2}{M},\nu =\frac{\mu }{M},$$
(21)
where the parameter $`\nu `$ takes values between 0 and 1/4, corresponding to the test mass limit and the equal mass case, respectively. Henceforth, we shall limit to the dynamics of the bound states generated by the two charged bodies, therefore $`e_1e_2<0`$ and we pose the coupling constant $`\alpha =e_1e_2>0`$. In the center of mass frame we have $`𝑷=𝒑_1=𝒑_2`$ and introducing the following reduced variables
$$\widehat{}=\frac{}{\mu },𝒑=\frac{𝑷}{\mu },\widehat{t}=\frac{\mu t}{\alpha },r=\frac{\mu q}{\alpha },$$
(22)
we can re-write the Hamiltonian, Eq. (15), in the more convenient form
$`\widehat{}(𝒓,𝒑)`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒑^2{\displaystyle \frac{1}{r}}{\displaystyle \frac{1}{8c^2}}(13\nu )𝒑^4{\displaystyle \frac{1}{2c^2}}{\displaystyle \frac{\nu }{r}}[𝒑^2+(𝒏𝒑)^2]`$ (25)
$`{\displaystyle \frac{1}{8c^4}}{\displaystyle \frac{1}{r}}\left[3\nu ^2(𝒏𝒑)^4+\nu \left(3\nu 2\right)𝒑^4+2\nu \left(\nu 1\right)𝒑^2(𝒏𝒑)^2\right]`$
$`+{\displaystyle \frac{1}{16c^4}}(15\nu +5\nu ^2)𝒑^6+{\displaystyle \frac{1}{4c^4}}{\displaystyle \frac{\nu }{r^2}}𝒑^2+{\displaystyle \frac{1}{4c^4}}{\displaystyle \frac{\nu }{r^3}}.`$
The above Hamiltonian is invariant under time translations and space rotations. We denote the two conserved quantities, that is the centre-of-mass non-relativistic energy and angular momentum, by
$$\widehat{}(𝒓,𝒑)=\widehat{}^{\mathrm{NR}}=\frac{_{\mathrm{c}.\mathrm{m}.}^{\mathrm{NR}}}{\mu },𝒓𝒑=𝒋=\frac{𝓙_{\mathrm{c}.\mathrm{m}.}}{\alpha }.$$
(26)
In the following we pose $`^{\mathrm{NR}}_{\mathrm{c}.\mathrm{m}.}^{\mathrm{NR}}`$ and $`𝓙𝓙_{\mathrm{c}.\mathrm{m}.}`$. Using the Hamilton-Jacobi formalism, we can summarize in a coordinate-invariant manner the two-charge dynamics by evaluating the “energy-levels” of the system. Introducing the reduced Hamilton principal-function $`\widehat{S}`$, defined by $`(\widehat{S}/𝒓)=𝒑`$, separating the time and angular coordinates and restricting to the planar motion, we can write
$$\widehat{S}=\widehat{}^{\mathrm{NR}}\widehat{t}+j\phi +\widehat{S}_r(r,\widehat{}^{\mathrm{NR}},j).$$
(27)
Solving the Hamilton-Jacobi equation $`\widehat{}(𝒓,𝒑)=\widehat{}^{\mathrm{NR}}`$ with respect to $`(d\widehat{S}_r/dr)=p_r=𝒏𝒑`$, using $`𝒑^2=(𝒏𝒑)^2+𝒋^2/r^2`$, we get
$$\widehat{S}_r(r,\widehat{}^{\mathrm{NR}},j)=𝑑r\sqrt{(r,\widehat{}^{\mathrm{NR}},j)},$$
(28)
where $``$ is a polynomial of the fifth order in $`1/r`$, explicitly given by:
$$(r,\widehat{}^{\mathrm{NR}},j)=A+\frac{2B}{r}+\frac{C}{r^2}+\frac{D_1}{r^3}+\frac{D_2}{r^4}+\frac{D_3}{r^5},$$
(29)
with
$`A`$ $`=`$ $`2\widehat{}^{\mathrm{NR}}+{\displaystyle \frac{1}{c^2}}(13\nu )(\widehat{}^{\mathrm{NR}})^2+{\displaystyle \frac{1}{c^4}}\nu (4\nu 1)(\widehat{}^{\mathrm{NR}})^3,`$ (30)
$`B`$ $`=`$ $`1+{\displaystyle \frac{1}{c^2}}(1\nu )\widehat{}^{\mathrm{NR}}+{\displaystyle \frac{1}{c^4}}{\displaystyle \frac{\nu }{2}}\left(2\nu 1\right)(\widehat{}^{\mathrm{NR}})^2,`$ (31)
$`C`$ $`=`$ $`j^2+{\displaystyle \frac{1}{c^2}}(1+\nu ),`$ (32)
$`D_1`$ $`=`$ $`{\displaystyle \frac{1}{c^2}}\nu j^2{\displaystyle \frac{1}{c^4}}\nu ^2j^2\widehat{}^{\mathrm{NR}}+{\displaystyle \frac{1}{c^4}}{\displaystyle \frac{\nu }{2}}\left(4\nu 1\right),`$ (33)
$`D_2`$ $`=`$ $`{\displaystyle \frac{3}{c^4}}\nu ^2j^2,`$ (34)
$`D_3`$ $`=`$ $`+{\displaystyle \frac{3}{4c^4}}\nu ^2j^4.`$ (35)
For our purposes we need to compute the reduced radial action variable
$$i_r^{\mathrm{real}}(\widehat{}^{\mathrm{NR}},j)=\frac{2}{2\pi }_{r_{\mathrm{min}}}^{r_{\mathrm{max}}}𝑑r\sqrt{(r,\widehat{}^{\mathrm{NR}},j)}.$$
(36)
To evaluate the above integral we use the formula (3.9) of Ref. , derived by performing a complex contour integration. The result for the radial action variable $`_R^{\mathrm{real}}=\alpha i_r^{\mathrm{real}}`$ reads:
$`_R^{\mathrm{real}}(^{\mathrm{NR}},𝒥)`$ $`=`$ $`{\displaystyle \frac{\alpha \mu ^{1/2}}{\sqrt{2^{\mathrm{NR}}}}}\left[1{\displaystyle \frac{1}{4}}(\nu 3){\displaystyle \frac{^{\mathrm{NR}}}{\mu c^2}}{\displaystyle \frac{1}{32}}(56\nu 3\nu ^2)\left({\displaystyle \frac{^{\mathrm{NR}}}{\mu c^2}}\right)^2\right]`$ (38)
$`𝒥+{\displaystyle \frac{\alpha ^2}{c^2𝒥}}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{\nu }{2}}{\displaystyle \frac{^{\mathrm{NR}}}{\mu c^2}}\right)+{\displaystyle \frac{1}{8}}(16\nu ){\displaystyle \frac{\alpha ^4}{c^4𝒥^3}}.`$
Finally, to get the “energy-levels” we solve the above equation in terms of the relativistic energy $`^\mathrm{R}=^{\mathrm{NR}}+Mc^2`$. Introducing the Delaunay action variable $`𝒩=_R^{\mathrm{real}}+𝒥`$, we get:
$`^\mathrm{R}(𝒩,𝒥)`$ $`=`$ $`Mc^2{\displaystyle \frac{1}{2}}{\displaystyle \frac{\alpha ^2\mu }{𝒩^2}}+{\displaystyle \frac{1}{c^2}}\alpha ^4\mu [{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{𝒥𝒩^3}}+{\displaystyle \frac{1}{8}}(3\nu ){\displaystyle \frac{1}{𝒩^4}}]+{\displaystyle \frac{1}{c^4}}\alpha ^6\mu [{\displaystyle \frac{3}{8}}{\displaystyle \frac{1}{𝒥^2𝒩^4}}`$ (40)
$`+{\displaystyle \frac{1}{16}}(5+3\nu \nu ^2){\displaystyle \frac{1}{𝒩^6}}+{\displaystyle \frac{1}{4}}(32\nu ){\displaystyle \frac{1}{𝒥𝒩^5}}+{\displaystyle \frac{1}{8}}(6\nu 1){\displaystyle \frac{1}{𝒥^3𝒩^3}}].`$
At 0PC order we recover the well known result of the degeneracy of the energy-levels in the Coulomb problem. Let us observe that at 1PC order, identifying $`𝒩/\mathrm{}`$ with the principal quantum-number and $`𝒥/\mathrm{}`$ with the total angular-momentum quantum-number, we obtain that Eq. (40) gives, e.g., the correct bound-state energies of the singlet states of the positronium ($`e_1=e_2`$ and $`m_1=m_2`$) in the (classical) limit $`𝒥/\mathrm{}1`$. Moreover, within the approximation $`𝒥/\mathrm{}1`$, our method captures all the centrifugal barrier shifts that have to be added by hand in . However, we cannot recover from Eq. (40) the correct quantum energy-levels at 2PC level, because at this order radiation reaction effects should have been taken into account. Indeed, in electrodynamics they enter at 1.5PC order, with a dipole-type interaction. Only if we limit to systems with $`e_1/m_1=e_2/m_2`$, we can postpone radiation reaction effects at the quadrupole order, which means at 2.5 PC level. In the present work we are interested in the conservative part of the bound states dynamics, hence we do not make the restriction $`e_1/m_1=e_2/m_2`$. The radiative corrections which contribute to the main part of the Lamb shift have been evaluated in , using the quasi-potential approach, and are of the order $`\alpha ^5\mathrm{log}\alpha `$. Corrections of the order $`\alpha ^5`$, $`\alpha ^6`$, $`\alpha ^6\mathrm{log}\alpha `$ have also been partially obtained in the literature for some quantum bound states of positronium and muonium .
## III “Effective” one-body description
The basic idea of the present work is to map the “real” two-body dynamics, described in the previous section, to an “effective” dynamics of a test particle of mass $`m_0`$ and charge $`e_0`$, moving in an external electromagnetic field. The action for the test particle is given by:
$$S_{\mathrm{eff}}=\left(m_0cds_0+\frac{1}{c}e_0A_\mu ^{\mathrm{eff}}(z)dz^\mu \right),$$
(41)
where $`A_{\mathrm{eff}}^\mu =(\mathrm{\Phi }_{\mathrm{eff}},𝑨_{\mathrm{eff}})`$. It is straightforward to derive that the effective Hamiltonian satisfies the well known equation
$$\frac{(_{\mathrm{eff}}e_0\mathrm{\Phi }_{\mathrm{eff}})^2}{c^2}=m_0^2c^2+\left(𝒑\frac{e_0}{c}𝑨_{\mathrm{eff}}\right)^2.$$
(42)
The effective electromagnetic field $`A_{\mathrm{eff}}^\mu `$ will be constructed in the form of an expansion in the dimensionless parameter $`\alpha _0/(m_0c^2R)`$, where $`\alpha _0=e_0^2`$ is the coupling constant and $`\alpha _0/(m_0c^2)`$ is the classical charge radius of $`m_0`$. Hence, we pose:
$`\mathrm{\Phi }_{\mathrm{eff}}(R)={\displaystyle \frac{e_0\varphi _0}{R}}\left[1+\varphi _1{\displaystyle \frac{\alpha _0}{m_0c^2R}}+\varphi _2\left({\displaystyle \frac{\alpha _0}{m_0c^2R}}\right)^2+\mathrm{}\right],`$ (44)
$`𝑨_{\mathrm{eff}}(R)={\displaystyle \frac{e_0𝒂}{cR}}\left[a_0+a_1{\displaystyle \frac{\alpha _0}{m_0c^2R}}+\mathrm{}\right],`$
where $`\varphi _0,\varphi _1,\varphi _2`$ and $`a_0,a_1`$ are dimensionless parameters and $`𝒂`$ is a vector with the dimension of a velocity. All these unknown coefficients will be fixed by the matching between the “real” and the “effective” description. Note that, in the above equations the variable $`R`$ stands for the effective radial coordinate and differs from the real separation $`R`$ used in Sec. II. Moreover, in Eqs. (44), (44) we have indicated only the terms we shall need up to 2PC order.
The dynamics of the one-body problem can be described, in a coordinate-invariant manner, in the Hamilton-Jacobi framework, by considering the “energy-levels” of the bound states of the particle $`m_0`$ in the external electromagnetic field. The Hamilton-Jacobi equation can be obtained from Eq. (42) posing $`_{\mathrm{eff}}=_0`$ and introducing the Hamilton principal-function $`S_{\mathrm{eff}}/𝑹=𝒑`$. Limiting to the motion in the equatorial plane ($`\theta =\pi /2`$) we can separate the variables, writing
$$S_{\mathrm{eff}}=_0t+𝒥_0\phi +S_R^0(R,_0,𝒥_0),$$
(45)
where $`_0`$ and $`𝒥_0|𝓙_0|`$ are the conserved energy and angular momentum defined by Eq. (41). The effective radial action variable reads
$$_R^{\mathrm{eff}}=\frac{2}{2\pi }_{R_{\mathrm{min}}}^{R_{\mathrm{max}}}𝑑R\frac{dS_R^0}{dR}.$$
(46)
Like in the two-body description we can derive the “energy-levels” of the “effective” one-body problem. They can be written as Note that, if a vector potential is present, the energy-levels could also depend on the magnetic number $`𝒥_z^0`$. In the present paper when dealing with a vector potential (see Sec. III B) we shall assume that the source of the magnetic field is the angular momentum, hence the magnetic field will be perpendicular to the plane of motion. This choice implicitly assumes $`𝒥_0𝒥_z^0`$.:
$`_0(𝒩_0,𝒥_0)`$ $`=`$ $`m_0c^2{\displaystyle \frac{1}{2}}{\displaystyle \frac{m_0\alpha _0^2}{𝒩_0^2}}+{\displaystyle \frac{1}{c^2}}\alpha _0^4m_0\left({\displaystyle \frac{_{3,1}}{𝒥_0𝒩_0^3}}+{\displaystyle \frac{_{4,0}}{𝒩_0^4}}\right)`$ (47)
$`+`$ $`{\displaystyle \frac{1}{c^4}}\alpha _0^6m_0\left[{\displaystyle \frac{_{3,3}}{𝒥_0^3𝒩_0^3}}+{\displaystyle \frac{_{4,2}}{𝒥_0^2𝒩_0^4}}+{\displaystyle \frac{_{5,1}}{𝒥_0𝒩_0^5}}+{\displaystyle \frac{_{6,0}}{𝒩_0^6}}\right],`$ (48)
where $`𝒩_0=_R^{\mathrm{eff}}+𝒥_0`$ and $`_{i,j}`$ are combinations of the coefficients $`\varphi _0,\varphi _1,\varphi _2`$ and $`a_0,a_1`$ given in Eqs. (44), (44).
Let us now define the rules to match the “real” to the “effective” problem. Like in , we find very natural sticking with the following relations between the adiabatic invariants:
$$𝒩=𝒩_0,𝒥=𝒥_0.$$
(49)
However, the way the “energy-levels”, Eq. (40) and Eq. (47), are related is more subtle. If we simply identify $`_0(𝒥_0,𝒩_0)=^\mathrm{R}(𝒥,𝒩)+(m_0M)c^2`$, and impose that the mass of the effective test particle coincides with the reduced mass, i.e. $`m_0=\mu `$, we obtain that already at 1PC order it is impossible to reduce the two-body dynamics to a one-body description. Hence, following we assume that there is a one-to-one mapping between the “real” and the “effective” energy-levels of the general form:
$$\frac{_0^{\mathrm{NR}}}{m_0c^2}=\frac{^{\mathrm{NR}}}{\mu c^2}\left[1+\alpha _1\frac{^{\mathrm{NR}}}{\mu c^2}+\alpha _2\left(\frac{^{\mathrm{NR}}}{\mu c^2}\right)^2\right],$$
(50)
where $`\alpha _1`$ and $`\alpha _2`$ are unknown coefficients that will be fixed by the matching. Given the aforesaid “rules”, we shall investigate in the subsequent sections three feasible ways the mapping can be implemented. The diverse descriptions differ by the choice of the effective electromagnetic field and the spacetime metric.
### A Effective scalar potential depending on the energy
In this section we study the possibility of reducing the two-body dynamics to a one-body one introducing, in the “effective” description, the scalar potential $`\mathrm{\Phi }_{\mathrm{eff}}`$ displayed in Eq. (44), and assuming that the vector potential $`𝑨_{\mathrm{eff}}`$ is zero. In this case the derivative of the radial Hamilton principal-function is given by:
$$\frac{dS_R^0}{dR}=2m_0_0^{\mathrm{NR}}2m_0e_0\mathrm{\Phi }_{\mathrm{eff}}\frac{𝒥_0^2}{R^2}+\frac{(_0^{\mathrm{NR}})^2}{c^2}+\frac{e_0^2\mathrm{\Phi }_{\mathrm{eff}}^2}{c^2}\frac{2e_0_0^{\mathrm{NR}}\mathrm{\Phi }_{\mathrm{eff}}}{c^2},$$
(51)
where we have introduced the non-relativistic energy $`_0^{\mathrm{NR}}=_0^\mathrm{R}m_0c^2`$. Plugging the above expression in Eq. (46) we get:
$`_R^{\mathrm{eff}}(_0^{\mathrm{NR}},𝒥_0)={\displaystyle \frac{\alpha _0m_0^{1/2}}{\sqrt{2_0^{\mathrm{NR}}}}}\left[\varphi _0{\displaystyle \frac{3\varphi _0}{4}}{\displaystyle \frac{_0^{\mathrm{NR}}}{m_0c^2}}+{\displaystyle \frac{5\varphi _0}{32}}\left({\displaystyle \frac{_0^{\mathrm{NR}}}{m_0c^2}}\right)^2\right]𝒥_0`$ (52)
$`+{\displaystyle \frac{\alpha _0^2}{𝒥_0c^2}}\left[{\displaystyle \frac{\varphi _0^2}{2}}\varphi _0\varphi _1\varphi _0\varphi _1{\displaystyle \frac{_0^{\mathrm{NR}}}{m_0c^2}}\right]+{\displaystyle \frac{1}{8}}{\displaystyle \frac{\alpha _0^4}{𝒥_0^3c^4}}\left[\varphi _0^412\varphi _0^3\varphi _1+8\varphi _0^2\varphi _2+4\varphi _0^2\varphi _1^2\right].`$ (53)
Identifying Eq. (52) with Eq. (38), assuming $`m_0=\mu `$ and using Eqs. (49), (50) we obtain the equations for the unknowns $`\varphi _0,\varphi _1,\varphi _2`$, $`a_0,a_1`$ and $`\alpha _1`$ and $`\alpha _2`$. In particular, at 0PC order we have
$$\varphi _0\alpha _0=\alpha ,$$
(54)
and we find quite natural to pose $`\varphi _0=1`$, that is $`e_0^2=\alpha _0=\alpha =e_1e_2`$. The equations at 1PC level are:
$$\varphi _0\alpha _0(2\alpha _13)=\alpha (\nu 3),\alpha _0^2(\varphi _0^22\varphi _0\varphi _1)=\alpha ^2,$$
(55)
while at 2PC order they read:
$`\varphi _0\alpha _0(512\alpha _112\alpha _1^2+16\alpha _2)=\alpha (56\nu 3\nu ^2),`$ (56)
$`\alpha _0^4(\varphi _0^4+4\varphi _0^2\varphi _1^212\varphi _0^3\varphi _1+8\varphi _0^2\varphi _2)=\alpha ^4(16\nu ),`$ (57)
$`\varphi _0\varphi _1\alpha _0^2={\displaystyle \frac{\nu }{2}}\alpha ^2.`$ (58)
Let us notice that at 1PC order, Eq. (55) gives $`\alpha _1=\nu /2`$ and $`\varphi _1=0`$. Then at 2PC order one can solve Eqs. (56) and (57) in terms of $`\alpha _2`$ and $`\varphi _2`$, obtaining $`\alpha _2=0`$ and $`\varphi _2=3\nu /4`$, but Eq. (58) is inconsistent. To solve this incompatibility we are obliged to introduce another parameter in the “effective” description. A simple possibility is to suppose that the diverse coefficients that appear in the effective scalar potential depend on an external parameter $`E_{\mathrm{ext}}`$, having the dimension of an energy, that is:
$`\varphi _0(E_{\mathrm{ext}})`$ $`=`$ $`\varphi _0^{(0)}+\varphi _0^{(2)}{\displaystyle \frac{E_{\mathrm{ext}}}{m_0c^2}}+\varphi _0^{(4)}\left({\displaystyle \frac{E_{\mathrm{ext}}}{m_0c^2}}\right)^2,`$ (59)
$`\varphi _1(E_{\mathrm{ext}})`$ $`=`$ $`\varphi _1^{(0)}+\varphi _1^{(2)}{\displaystyle \frac{E_{\mathrm{ext}}}{m_0c^2}},`$ (60)
$`\varphi _2(E_{\mathrm{ext}})`$ $`=`$ $`\varphi _2^{(0)}.`$ (61)
We find that in order to implement the matching with the “real” description the parameter $`E_{\mathrm{ext}}`$ should be fixed equal to the “effective” non-relativistic energy, i.e. $`E_{\mathrm{ext}}_0^{\mathrm{NR}}`$. In more details, the introduction of an energy dependence in the coefficients $`\varphi _0,\varphi _1,\varphi _2`$ reshuffles the $`c^2`$ expansion of Eq. (52), modifying the Eqs. (55)–(58) and allowing to solve in many ways the constraint equations. The simplest solution is envisaged by requiring that the energy-dependence enters only at 2PC order in the coefficient $`\varphi _1`$. In this case, the solution reads:
$`\varphi _0^{(0)}=1,\varphi _0^{(2)}=0,\varphi _0^{(4)}=0,`$ (62)
$`\varphi _1^{(0)}=0,\varphi _1^{(2)}={\displaystyle \frac{\nu }{2}},\varphi _2^{(0)}={\displaystyle \frac{3}{4}}\nu ,`$ (63)
$`\alpha _1={\displaystyle \frac{\nu }{2}},\alpha _2=0.`$ (64)
To summarize, we have succeeded in mapping the two-body dynamics onto the one of a test particle of mass $`m_0=\mu `$ moving in the external scalar potential:
$$\mathrm{\Phi }_{\mathrm{eff}}(R,E_{\mathrm{ext}})=\frac{e_0}{R}\left[1\frac{\nu }{2}\left(\frac{E_{\mathrm{ext}}}{m_0c^2}\right)\left(\frac{\alpha _0}{m_0c^2R}\right)\frac{3\nu }{4}\left(\frac{\alpha _0}{m_0c^2R}\right)^2\right],$$
(65)
where $`E_{\mathrm{ext}}_0^{\mathrm{NR}}`$. We have found that the matching is implemented relating the “real” and “effective” energy-levels by the formula:
$$\frac{_0^{\mathrm{NR}}}{m_0c^2}=\frac{^{\mathrm{NR}}}{\mu c^2}\left[1+\frac{\nu }{2}\frac{^{\mathrm{NR}}}{\mu c^2}\right],$$
(66)
which, as noticed in , gives the following relation between the real total relativistic energy $``$ and the effective relativistic energy $`_0`$:
$$\frac{_0}{m_0c^2}\frac{^2m_1^2c^4m_2^2c^4}{2m_1m_2c^4}.$$
(67)
The above equation has a rather interesting property. In the limit $`m_1m_2`$ the effective energy of the effective particle equals the energy of the particle 1 in the rest frame of particle 2 (and reciprocally if $`m_2m_1`$). Moreover, the result (67) coincides with the one derived in Ref. in the context of quantum electrodynamics. We find quite remarkable that our way of relating the “real” and “effective” energy-levels agrees with the one introduced in . Nevertheless, we consider the dependence on the energy of the effective scalar potential, Eq. (65), quite unsatisfactory, though envisaged by Todorov et al. in the quasi-potential approach. Indeed, in our context the presence of an external parameter in the scalar potential obscures the nature of the mapping and complicates the possibility of incorporating radiation reaction effects. Certainly, this cannot be achieved straightforwardly in the way suggested in for the gravitational case.
As a final remark, let us note that if we were using the effective description introduced in the quasi-potential approach by Todorov et al. , we should have considered a test particle with effective mass, $`m_{\mathrm{eff}}`$, and effective energy, $`_{\mathrm{eff}}`$, given by:
$$m_{\mathrm{eff}}(_{\mathrm{real}})=\frac{m_1m_2c^2}{_{\mathrm{real}}},_{\mathrm{eff}}\frac{_{\mathrm{real}}^2m_1^2c^4m_2^2c^4}{2_{\mathrm{real}}}.$$
(68)
We have investigated the possibility of introducing an energy dependence in the effective mass of the test particle, but we found that, in this case, it is not possible to overcome the inconsistency in the matching equations that raised at 2PC order. A way out could be to introduce also an energy dependence in the effective coupling $`\alpha _{\mathrm{eff}}`$, but we find this possibility not very appealing.
### B Effective vector potential depending on the angular momentum
We have seen in the previous section that, at 2PC level, in order to cope with an inconsistency of the constraint equations, we were obliged to introduce an external parameter in the coefficients of the scalar potential. In this section we shall investigate the possibility of overcoming the above inconsistency by introducing, in the “effective” description, a scalar potential $`\mathrm{\Phi }_{\mathrm{eff}}`$, independent of any external parameter, and a vector potential $`𝑨_{\mathrm{eff}}`$ which will depend on an external vector $`𝑱_{\mathrm{ext}}`$. In order to implement the matching, we have found that it is sufficient to limit to the following form of the vector potential (see Eq. (44)):
$$𝑨_{\mathrm{eff}}=\frac{e_0(𝑱_{\mathrm{ext}}𝑹)}{m_0cR^3}\left[a_0+a_1\frac{\alpha _0}{m_0c^2R}+\mathrm{}\right],$$
(69)
where $`𝑱_{\mathrm{ext}}`$ is supposed to be perpendicular to the plane of motion Note that, with this choice of the vector potential the magnetic field will be perpendicular to the plane of motion.. In the Hamilton-Jacobi framework, restricting to $`\theta =\pi /2`$, we have
$$𝒑=\frac{S_{\mathrm{eff}}}{𝑹}=\widehat{𝒆}_R\frac{S_{\mathrm{eff}}}{R}+\widehat{𝒆}_\phi \frac{1}{R}\frac{S_{\mathrm{eff}}}{\phi },$$
(70)
where $`\widehat{𝒆}_R`$ and $`\widehat{𝒆}_\phi `$ are vectors of the orthonormal basis. Due to the particular choice of the vector $`𝑱_{\mathrm{ext}}`$ we made, the following equation holds:
$$𝑨_{\mathrm{eff}}=\frac{e_0J_{\mathrm{ext}}\widehat{𝒆}_\phi }{m_0cR^2}\left[a_0+a_1\frac{\alpha _0}{m_0c^2R}+\mathrm{}\right],$$
(71)
where $`J_{\mathrm{ext}}=|𝑱_{\mathrm{ext}}|`$. Finally, using $`S_{\mathrm{eff}}/\phi =𝒥_0`$ (see Eq. (45)), we get:
$$𝒑𝑨_{\mathrm{eff}}=\frac{e_0J_{\mathrm{ext}}𝒥_0}{m_0cR^3}\left[a_0+a_1\frac{\alpha _0}{m_0c^2R}+\mathrm{}\right],𝑨_{\mathrm{eff}}^2=\frac{e_0^2J_{\mathrm{ext}}^2a_0^2}{m_0^2c^2R^4}+\mathrm{}.$$
(72)
Note the crucial fact that, with the very special choice of the vector potential we made, $`𝒑𝑨_{\mathrm{eff}}`$ does not depend on $`𝒑_R`$. Plugging the above expressions in the Hamilton-Jacobi equation, Eq. (42), with $`_{\mathrm{eff}}=_0^{\mathrm{NR}}+m_0^2c^2`$ we obtain:
$`{\displaystyle \frac{dS_R^0}{dR}}`$ $`=`$ $`2m_0_0^{\mathrm{NR}}2m_0e_0\mathrm{\Phi }_{\mathrm{eff}}{\displaystyle \frac{𝒥_0^2}{R^2}}+{\displaystyle \frac{(_0^{\mathrm{NR}})^2}{c^2}}+{\displaystyle \frac{e_0^2\mathrm{\Phi }_{\mathrm{eff}}^2}{c^2}}{\displaystyle \frac{2e_0_0^{\mathrm{NR}}\mathrm{\Phi }_{\mathrm{eff}}}{c^2}}`$ (74)
$`+{\displaystyle \frac{2𝒥_0J_{\mathrm{ext}}}{R^2}}\left[a_0{\displaystyle \frac{\alpha _0}{m_0c^2R}}+a_1\left({\displaystyle \frac{\alpha _0}{m_0c^2R}}\right)^2\right]{\displaystyle \frac{J_{\mathrm{ext}}^2}{R^2}}a_0^2\left({\displaystyle \frac{\alpha _0}{m_0c^2R}}\right)^2,`$
where $`\mathrm{\Phi }_{\mathrm{eff}}`$ is given by Eq. (44). Evaluating the radial action variable (see Eq. (46)) we finally get:
$`_R^{\mathrm{eff}}(_0^{\mathrm{NR}},𝒥_0,J_{\mathrm{ext}})={\displaystyle \frac{\alpha _0m_0^{1/2}}{\sqrt{2_0^{\mathrm{NR}}}}}[\varphi _0{\displaystyle \frac{3\varphi _0}{4}}{\displaystyle \frac{_0^{\mathrm{NR}}}{m_0c^2}}+{\displaystyle \frac{5\varphi _0}{32}}\left({\displaystyle \frac{_0^{\mathrm{NR}}}{m_0c^2}}\right)^2]𝒥_0+{\displaystyle \frac{\alpha _0^2}{𝒥_0c^2}}[{\displaystyle \frac{\varphi _0^2}{2}}`$ (75)
$`\varphi _0\varphi _1\varphi _0a_0{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}}+{\displaystyle \frac{_0^{\mathrm{NR}}}{m_0c^2}}(\varphi _0\varphi _1\varphi _0a_0{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}}+a_1{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}}+a_0^2{\displaystyle \frac{J_{\mathrm{ext}}^2}{𝒥_0^2}})]+{\displaystyle \frac{1}{8}}{\displaystyle \frac{\alpha _0^4}{𝒥_0^3c^4}}[\varphi _0^4`$ (76)
$`12\varphi _0^3\varphi _1+8\varphi _0^2\varphi _2+4\varphi _0^2\varphi _1^2+24\varphi _0^2\varphi _1a_0{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}}12\varphi _0^3a_0{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}}+12\varphi _0^2a_1{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}}+24\varphi _0^2a_0^2{\displaystyle \frac{J_{\mathrm{ext}}^2}{𝒥_0^2}}].`$ (77)
Let us impose that the above equation coincides with the analogous expression for the “real” description, given by Eq. (38). Assuming $`m_0=\mu `$ and using Eqs. (49), (50) we derive the new constraint equations to be satisfied. At 0PC order we still have $`\varphi _0\alpha _0=\alpha `$, and we pose $`\varphi _0=1`$, while at 1PC level we get:
$$\varphi _0\alpha _0(2\alpha _13)=\alpha (\nu 3),\alpha _0^2\left(\varphi _0^22\varphi _0\varphi _12\varphi _0a_0\frac{J_{\mathrm{ext}}}{𝒥_0}\right)=\alpha ^2.$$
(79)
The first equation in (79) gives $`\alpha _1=\nu /2`$, while the second one is automatically satisfied if we make the rather natural requirement that either the Coulomb potential does not have any correction at 1PC order ($`\varphi _1=0`$) or the vector potential enters only at the next Coulombian order ($`a_0=0`$). Finally, the 2PC order constraints read:
$`\varphi _0\alpha _0(512\alpha _112\alpha _1^2+16\alpha _2)=\alpha (56\nu 3\nu ^2),`$ (80)
$`\alpha _0^4(\varphi _0^4+4\varphi _0^2\varphi _1^212\varphi _0^3\varphi _1+8\varphi _0^2\varphi _2+24\varphi _0^2\varphi _1a_0{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}}`$ (81)
$`12\varphi _0^3a_0{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}}+24\varphi _0^2a_0^2{\displaystyle \frac{J_{\mathrm{ext}}^2}{𝒥_0^2}}+12\varphi _0^2a_1{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}})=\alpha ^4(16\nu ),`$ (82)
$`\alpha _0^2\left(\varphi _0\varphi _1+\varphi _0a_0{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}}a_1{\displaystyle \frac{J_{\mathrm{ext}}}{𝒥_0}}a_0^2{\displaystyle \frac{J_{\mathrm{ext}}^2}{𝒥_0^2}}\right)=\nu {\displaystyle \frac{\alpha ^2}{2}}.`$ (83)
Plugging the results obtained at 0PC and 1PC order in Eqs. (80)–(83) and assuming that the external vector $`J_{\mathrm{ext}}`$ coincides with the constant of motion $`𝒥_0`$, we end up with the unique, rather simple solution:
$$\varphi _2=0,a_1=\frac{\nu }{2},\alpha _2=0.$$
(84)
In conclusion, in this Section we have obtained that at 2PC order it is possible to reduce the two-charge dynamics to the one of a test particle moving in an effective electromagnetic field described by a Coulomb potential $`\mathrm{\Phi }_{\mathrm{eff}}(R)=e_0/R`$ and a vector potential dependent on the external vector $`𝑱_{\mathrm{ext}}`$ ($`𝓙_0`$):
$$𝑨_{\mathrm{eff}}(R,J_{\mathrm{ext}})=\frac{\nu }{2}\frac{e_0\alpha _0}{m_0^2c^3}\frac{(𝑱_{\mathrm{ext}}𝑹)}{R^4}.$$
(85)
Moreover, quite remarkably, we have found, under rather natural assumptions, that the one-to-one mapping between the “real” and the “effective” energy-levels is still given by the formula (67). However, as already discussed at the end of the previous section, the fact that the electromagnetic field still has to depend on external parameters is not very desirable. In the next section we shall investigate a feasible way out.
### C Effective metric
So far we have seen that in order to succeed in reducing the two-body dynamics onto a one-body description we were obliged to introduce external parameters, which have been identified either with the energy or the angular momentum of the test particle $`m_0`$. This result is not very appealing, especially when we want to incorporate radiation reaction effects. A possible way out would be to relax the hypothesis that in the one-body description the test particle move in a flat spacetime. The effective spacetime metric should be viewed as an effective way of describing the global exchange of energy between the two charged particles in the “real” description.
The most general spherical symmetric metric written in Schwarzschild gauge has the form:
$$ds_{\mathrm{eff}}^2=A(R)c^2dt^2+B(R)dR^2+R^2(d\theta ^2+\mathrm{sin}\theta ^2d\phi ^2),$$
(86)
where the coefficients $`A(R)`$ and $`B(R)`$ are given as an expansion in the dimensionless parameter $`\alpha _0/(m_0c^2R)`$, that is:
$`A(R)=1+A_1{\displaystyle \frac{\alpha _0}{m_0c^2R}}+A_2\left({\displaystyle \frac{\alpha _0}{m_0c^2R}}\right)^2+A_3\left({\displaystyle \frac{\alpha _0}{m_0c^2R}}\right)^3+\mathrm{},`$ (87)
$`B(R)=1+B_1{\displaystyle \frac{\alpha _0}{m_0c^2R}}+B_2\left({\displaystyle \frac{\alpha _0}{m_0c^2R}}\right)^2+\mathrm{}.`$ (88)
The reduction, from the two-body problem to the one-body one, can simply be implemented assuming that in the “effective” description only the scalar potential $`\mathrm{\Phi }_{\mathrm{eff}}`$ is different from zero. In this case the derivative of the Hamilton principal-function reads:
$$\frac{dS_R^0}{dR}=\frac{B(R)}{c^2A(R)}(_0+m_0c^2e_0\mathrm{\Phi }_{\mathrm{eff}})^2\frac{B(R)}{R^2}𝒥_0^2B(R)m_0^2c^2,$$
(89)
and for the radial action variable we derive:
$`_R^{\mathrm{eff}}(_0^{\mathrm{NR}},𝒥_0)`$ $`=`$ $`{\displaystyle \frac{\alpha _0m_0^{1/2}}{\sqrt{2_0^{\mathrm{NR}}}}}\left[𝒜+{\displaystyle \frac{_0^{\mathrm{NR}}}{m_0c^2}}+𝒞\left({\displaystyle \frac{_0^{\mathrm{NR}}}{m_0c^2}}\right)^2\right]𝒥_0`$ (91)
$`+{\displaystyle \frac{\alpha _0^2}{𝒥_0c^2}}\left[𝒟+{\displaystyle \frac{_0^{\mathrm{NR}}}{m_0c^2}}\right]+{\displaystyle \frac{\alpha _0^4}{𝒥_0^3c^4}},`$
where the various coefficients can be written explicitly as:
$`𝒜`$ $`=`$ $`\varphi _0{\displaystyle \frac{1}{2}}A_1,`$ (92)
$``$ $`=`$ $`{\displaystyle \frac{3}{4}}\varphi _0+\left(B_1{\displaystyle \frac{7}{8}}A_1\right),`$ (93)
$`𝒞`$ $`=`$ $`{\displaystyle \frac{5}{32}}\varphi _0+\left({\displaystyle \frac{B_1}{4}}{\displaystyle \frac{19}{64}}A_1\right),`$ (94)
$`𝒟`$ $`=`$ $`\varphi _0\left(\varphi _1{\displaystyle \frac{B_1}{2}}+A_1\right)+{\displaystyle \frac{1}{2}}\varphi _0^2{\displaystyle \frac{1}{4}}A_1B_1+{\displaystyle \frac{A_1^2}{2}}{\displaystyle \frac{A_2}{2}},`$ (95)
$``$ $`=`$ $`\varphi _0\left(\varphi _1+A_1{\displaystyle \frac{B_1}{2}}\right)+A_1^2A_2{\displaystyle \frac{1}{2}}A_1B_1{\displaystyle \frac{B_1^2}{8}}+{\displaystyle \frac{B_2}{2}},`$ (96)
$``$ $`=`$ $`{\displaystyle \frac{1}{64}}\left(24A_1^448A_1^2A_2+8A_2^2+16A_1A_38A_1^3B_1+8A_1A_2B_1A_1^2B_1^2+4A_1^2B_2\right)`$ (100)
$`+{\displaystyle \frac{\varphi _0}{16}}(16\varphi _1A_1^2+24A_1^3+8\varphi _2A_1+8\varphi _1A_232A_1A_2+8A_3+4\varphi _1A_1B_1`$
$`8A_1^2B_1+4A_2B_1A_1B_1^2+4A_1B_2)+{\displaystyle \frac{\varphi _0^3}{4}}(6\varphi _1+4A_1B_1)+{\displaystyle \frac{\varphi _0^4}{8}}`$
$`+{\displaystyle \frac{\varphi _0^2}{16}}(8\varphi _1^240\varphi _1A_1+32A_1^2+16\varphi _220A_2+8\varphi _1B_110A_1B_1B_1^2+4B_2).`$
The above expressions coincide with the ones obtained in pure general relativity , once the limit $`\varphi _00`$ is considered and $`\alpha _0`$ is identified with the analogous quantity in the gravitational case, i.e. with $`Gm_1m_2`$ ($`G`$ is the Newton constant). Let us now equate the “real”, Eq. (38) and the “effective”, Eq. (91), radial action variables, assuming that the following relations hold: $`𝒥_0=𝒥`$, $`m_0=\mu `$ and Eq. (50). At 0PC order we get the constraint $`\alpha _0(\varphi _0A_1/2)=\alpha `$ which can be naturally fulfilled imposing that $`A_1=0`$ and posing $`\varphi _0=1`$, as above. At 1PC level we derive:
$`2\varphi _0\alpha _0(2\alpha _13)+\alpha _0(7A_18B_12A_1\alpha _1)=2\alpha (\nu 3),`$ (101)
$`2\alpha _0^2(\varphi _0^2+\varphi _0(2A_1B_12\varphi _1))+\alpha _0^2(2A_1^22A_2A_1B_1)=2\alpha ^2.`$ (102)
If we demand that at this order the scalar potential and the effective metric do not differ from the Coulomb potential and the flat spacetime metric, respectively, i.e. we pose $`\varphi _1=0,A_2=0,B_1=0`$, we find that Eq. (101) gives $`\alpha _1=\nu /2`$ while Eq. (102) is automatically satisfied. Inserting these values in the constraint equations at 2PC order and imposing that there are no corrections to the Coulomb potential at this order ($`\varphi _2=0`$) we obtain the unique simple solution:
$$\alpha _2=0,A_3=\nu ,B_2=\nu .$$
(103)
Hence, we have found that with the introduction of an effective metric we are not obliged to introduce in the electromagnetic field any dependence on external parameters, neither the energy nor the angular momentum. Moreover, up to 2PC order we find that there is no need of modifying the Coulomb scalar potential, i.e. $`\mathrm{\Phi }_{\mathrm{eff}}(R)=e_0/R`$ and the “energy-levels” of the real and “effective” description are still related by Eq. (67). Finally, the external spacetime metric is simply given by:
$$A(R)=1+\nu \left(\frac{\alpha _0}{m_0c^2R}\right)^3,B(R)=1\nu \left(\frac{\alpha _0}{m_0c^2R}\right)^2.$$
(104)
## IV Conclusions
In this paper we have analysed the application of a new approach to studying the relativistic dynamics of the bound states of two classical charged particles, with comparable masses, interacting electromagnetically. The key idea, originally introduced investigating the two-body problem in general relativity , has been to map the “real” two-body problem onto the one of a test particle moving in an external electromagnetic field.
We have found that the matching can be implemented imposing the following rather natural “rules”: i) the adiabatic invariants $`𝒩`$ and $`𝒥`$ in the two descriptions have to be identified; ii) the reduced mass of the “real” system, $`\mu `$, has to coincide with the mass of the effective particle, $`m_0`$, and iii) the energy axis between the two problems has to be transformed. Let us note immediately that, a bottom-line of our results has been that, in all the three cases considered (see Sec. III A, III B and III C), we have found quite naturally that the energy axis, between the two descriptions, has to change in such a way that the effective energy of the effective particle coincides with the energy of the particle 1 in the rest frame of particle 2 in the limit $`m_1m_2`$ (and vice versa) (see Eq. (67)).
Nevertheless, contrary to the results obtained in general-relativity , the requirements i), ii) and iii) envisaged above, do not fix uniquely the external electromagnetic field, with which the effective test particle $`m_0`$ interacts. In fact, we have found that, in order to overcome an inconsistency in the constraint equations which define the matching, we had to introduce an external parameter either in the scalar potential, Eq. (65), or in the vector potential, Eq. (85). These parameters have to be identified with the non-relativistic energy and the angular momentum of the effective test particle $`m_0`$, respectively. As pointed out above and in Ref. , the dependence of the effective electromagnetic field on some external parameter makes the mapping between the two descriptions quite awkward and complicates the inclusion of radiation reaction effects. A possible solution of this issue is to relax the hypothesis that the test particle moves in a flat spacetime. Indeed, in this case we have found that the conditions i), ii) and iii) fix rather naturally the external scalar potential and the effective metric. They provide, up to 2PC order, an effective Coulomb potential and a rather simple $`\nu `$-deformed flat metric (see Eq. (104)).
Once the matching has been successfully defined, to have a complete knowledge of the “real” dynamics through the auxiliary “effective” one, we can construct, like in , the canonical transformation which relates the variables of the relative motion in the “real” description, to the coordinates and momenta of the test particle in the “effective” problem. However, this calculation goes beyond the scope of the present paper.
Finally, a last remark. In Sec. III B we have introduced a vector potential in the effective description in such a way that the source of the magnetic field is the angular momentum of the system. This study suggests the investigation, in the general relativity context , of relaxing the hypothesis of mapping the “real” two-body dynamics onto the one of a test particle moving in a deformed Schwarzschild spacetime. Indeed, it could well be possible to match the two problems appealing to an effective deformed Kerr spacetime.
## Acknowledgements
It is a pleasure to thank Thibault Damour, Scott Hughes, Gerhard Schäfer and Kip Thorne for useful discussions and/or for comments on this manuscript.
This research is supported by the Richard C. Tolman Fellowship and by NSF Grant AST-9731698 and NASA Grant NAG5-6840.
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# Can a frustrated spin-cluster model describe the low-temperature physics of NaV2O5 ?
## Abstract
Recent experimental evidence suggest the existence of three distinct V-valence states (V<sup>+4</sup>, V<sup>+4.5</sup> and V<sup>+5</sup>) in the low-temperature phase of NaV<sub>2</sub>O<sub>5</sub> in apparent discrepancy with the observed spin-gap. We investigate a novel spin cluster model, consisting of weakly coupled, frustrated four-spin clusters aligned along the crystallographic b-axis that was recently proposed to reconcile these experimental observations. We have studied the phase diagram and the magnon dispersion relation of this model using DMRG, exact diagonalization and a novel cluster-operator theory. We find a spin-gap for all parameter values and two distinct phases, a cluster phase and a Haldane phase. We evaluate the size of the gap and the magnon dispersion and find no parameter regime which would reproduce the experimental results. We conclude that this model is inappropriate for the low-temperature regime of NaV<sub>2</sub>O<sub>5</sub>.
PACS numbers: 75.10.d, 72.15.N, 71.20.B, 75.10.J
Introduction Recent investigations of electronically quasi one-dimensional (1D) transition metal compounds probe the limits of our understanding of the interplay between structural and electronic effects in such low dimensional materials. In NaV<sub>2</sub>O<sub>5</sub>, a prototypical example for this class of materials, V-ions are arranged in ladders along the crystallographic b-direction. Measurements of the magnetic susceptibility in the high-temperature phase indicate the presence of only one equivalent V-site with valence V<sup>+4.5</sup>, consistent with a model where the electrons in bonding V-O-V orbitals along the rungs of the ladder form a 1D Heisenberg chain.
At $`T_C=34\text{K}`$ the unit-cell doubles along the a- and b- and quadruples along the c-axis in a phase transition of as-of-yet unknown origin. At the same time a spin-gap of $`\mathrm{\Delta }_{min}=10\text{meV}`$ opens and charge ordering $`2\text{V}^{+4.5}\text{V}^{+4}+\text{V}^{+5}`$ sets in . The observed charge ordering is inconsistent with a generic spin-Peierls scenario and raises the question about the driving force (lattice, magnetic or Coulomb) responsible for this transition. Since NaV<sub>2</sub>O<sub>5</sub> is an insulator, the discussion of the material is simplified by the introduction of pseudospins for the charge degrees of freedom that couple to the spin-degrees of freedom . The effective spin-Hamiltonian depends, consequently, on the pattern of charge order and may differ in the high- and the low-temperature phase.
The occurrence of two well defined magnon-branches for $`T<T_C`$ in NaV<sub>2</sub>O<sub>5</sub> along the $`a`$ direction (perpendicular to the chains), as measured by neutron scattering , had been explained tentatively by a model, where the charge orders in a ‘zig-zag’ pattern in the low-temperature phase . This proposal has been questioned by recent analysis of the low-temperature crystal structure . Based on bond-charge models, the existence of three different V-valence states (V<sup>+4</sup>, V<sup>+4.5</sup> and V<sup>+5</sup>) has been proposed , as illustrated in Fig. 1. In this analysis, pairs of V<sup>+4.5</sup> form dimerized spin-chains on every other ladder, which alone could explain the observed spin-gap . A puzzle is posed however, by the presence of free isolated moments on the V<sup>+4</sup> ions on the remaining ladders, which is inconsistent with the existence of such a gap.
As one possible reconciliation, Boer et al. recently proposed that the V<sup>+4</sup> moments are quenched by their interaction with the neighboring V<sup>+4.5</sup> sites of the adjacent dimerized V-O-V ladder. Within this model, clusters of six Vanadiums each (and with four spins) would be weakly coupled and the observed spin-gap would arise not from the dimerization but locally from the gap of the isolated clusters.
To distinguish between these fundamentally different mechanisms we study this model by a series of complementary approaches, using DMRG , exact-diagonalization and a novel bond-cluster theory, to map all physically relevant regions of its phase diagram. We find that the ground state varies continuously from a cluster-phase for large $`J^{}`$ to a Haldane-phase for small $`J^{}`$ (see Fig. 1). We evaluate the gap and the dispersion and find that there is no parameter regime that would explain the neutron-scattering data .
The spin-cluster model We denote by $`𝐒_{n,i}`$ ($`i=1,\mathrm{},4`$) the four spins of the $`n`$’th cluster, compare Fig. 1. The Hamiltonian is then
$`H=J_1{\displaystyle \underset{n}{}}𝐒_{n,1}𝐒_{n,2}+J_2{\displaystyle \underset{n}{}}𝐒_{n,1}𝐒_{n+1,2}`$
$$+J^{}\underset{n}{}\left(𝐒_{n,1}+𝐒_{n,2}\right)\left(𝐒_{n,3}+𝐒_{n,4}\right),$$
(1)
where $`J_1=J(1+\delta )`$ and $`J_2=J(1\delta )`$ (with $`J_1,J_2,J^{}>0`$). $`\delta `$ is the degree of dimerization. For $`J^{}=0`$ the $`𝐒_{n,1/2}`$ form a dimerized chain with an in-chain gap $`J\delta ^{2/3}`$. A particular property of Eq. (1) is the local coupling to the total spin $`𝐒_{n,3}+𝐒_{n,4}`$, which is consequently a (locally) conserved quantity, $`(𝐒_{n,3}+𝐒_{n,4})^2=S_n(S_n+1)`$ for any $`n`$. In the ground-state $`S_n1`$.
We consider first an isolated cluster and denote by $`s_{ij}`$ and $`t_{ij}^\alpha `$ the wavefunctions of the singlet and of the triplets ($`\alpha =1,0,+1`$) of the spins $`i`$ and $`j`$ ($`i,j=1,\mathrm{}4,`$). The low-energy states are
$$\psi _1=\frac{1}{\sqrt{3}}\left[t_{12}^0t_{34}^0t_{12}^+t_{34}^{}t_{12}^{}t_{34}^+\right],$$
(2)
$$\psi _2=s_{12}^{}s_{34}^{},\psi _3^\alpha =s_{12}^{}t_{34}^\alpha ,$$
(3)
$$\psi _4^0=\frac{1}{\sqrt{2}}\left[t_{12}^+t_{34}^{}t_{12}^{}t_{34}^+\right],$$
(4)
where $`\psi _4^0`$ is the $`S^z=0`$ component of the triplet $`\psi _4^\alpha `$. The corresponding energies are $`E_1=2J^{}+J_1/4`$, $`E_2=E_3=3J_1/4`$ and $`E_4=J^{}+J_1/4`$ .
For $`J^{}/J_1>0.5`$ the singlet $`\psi _1`$ is the ground-state (we denote this region the ‘cluster phase’). For $`J^{}/J_1<0.5`$ the ground-state of the isolated cluster is four-fold degenerate, the singlet $`\psi _2`$ and the triplet $`\psi _3`$ have the same energy. Note that the intercluster coupling $`J_2`$ will not mix $`\psi _2`$ and $`\psi _3`$, since the local spin $`𝐒_{n,3}+𝐒_{n,4}`$ is conserved. Intercluster coupling $`J_2`$ will lead to an antiferromagnetic interaction $`J_H(J^{}J_2)^2/J_1^3`$ between the moments of the $`\psi _3`$ states, as can be evaluated easily in second-order perturbation in $`J_2`$ (using the complete set of eigenstates of the cluster). The total energy is therefore lowered by $`J_2`$ when all cluster states are $`\psi _3`$. The $`S=1`$ moments of the $`\psi _3`$ states thus form an effective spin-1 chain with a Haldane gap $`\mathrm{\Delta }_H=0.41050J_H`$ . We denote this region therefore the ‘Haldane phase’.
We have evaluated the energy gap of the spin-cluster model by DMRG , using the finite-size algorithm with open boundary conditions for systems with $`L=32`$ and $`L=64`$ spins. The ground-state has $`N_{}=L/2`$ up-spins and $`N_{}=L/2`$ down-spins. We retained typically 60 states of the density matrix, checking the convergence by additional calculations with 40 and 90 states respectively. We evaluated the gap by two complementary methods, namely (i) by targeting two states in the sector with $`N_{}=L/2=N_{}`$ and (ii) by targeting the ground states in (a) the sector with $`N_{}=L/2=N_{}`$ and (b) $`N_{}=L/2+1`$ and $`N_{}=L/21`$. We find complete consistency and present the results in Fig. 2 for some selected values for the dimerization $`\delta `$. The finite-size corrections are smaller than the symbol sizes. We find a rapidly decreasing gap as a function of decreasing $`J^{}/J`$ and a smooth crossover between the cluster- and the Haldane phase. As the symmetry of these two phases is the same, we do not expect a phase transition in the thermodynamic limit.
Cluster-operator theory In the cluster-phase two low-lying triplet modes, $`\psi _3^\alpha `$ and $`\psi _4^\alpha `$, are relevant. In order to take the effect of the intercluster coupling $`J_2`$ into account we describe the seven degrees of freedom of cluster $`n`$ by bosonic degrees of freedom: $`s_n^{}`$ for the singlet ($`\psi _1`$) $`b_{n,3,\alpha }^{}`$ and $`b_{n,4,\alpha }^{}`$ for the triplets ($`\psi _3`$ and $`\psi _4`$). The low-lying singlet $`\psi _2`$ does not couple and may be disregarded here. This approach generalizes the bond-operator theory for dimerized spin-chains to the case of spin-clusters. The constraint $`s_n^{}s_n^{}+_{\tau ,\alpha }b_{n,\tau ,\alpha }^{}b_{n,\tau ,\alpha }^{}=1`$ ($`\tau =3,4`$) restricts the bosonic Hilbert space to the physical one. The spin-operators take the form
$$S_{n,1/2}^z=\pm \frac{b_{n,3,0}^{}s_n^{}+s_n^{}b_{n,3,0}^{}}{\sqrt{12}}\frac{b_{n,4,0}^{}s_n^{}+s_n^{}b_{n,4,0}^{}}{\sqrt{6}}.$$
(5)
Note, that there are no terms $`b_{n,\tau ,\alpha }^{}b_{n,\tau ^{},\alpha ^{}}^{}`$ corresponding to triplet-triplet interactions. In linearized Holstein-Primakov approximation (LHP), we substitute $`s_n^{}1`$ and $`s_n^{}1`$ in Eq. (5) and in similar expressions for $`S_{n,1/2}^{x/y}`$. This approximation retains spin-rotational invariance and we may disregard the index $`\alpha =1,0,1`$ for the triplet operators. We obtain for the LHP Hamilton-operator in momentum space $`H^{(LHP)}=H_0+H_2^{(1)}+H_2^{(2)}`$ with $`H_0=_{k,\tau }\mathrm{\Delta }_\tau b_{k,\tau }^{}b_{k,\tau }^{}`$ ($`\mathrm{\Delta }_\tau =E_\tau E_1`$). The intercluster coupling is given by
$`H_2^{(1)}={\displaystyle \frac{J_2}{12}}{\displaystyle \underset{k}{}}[2\mathrm{cos}(2bk)(2b_{k,4}^{}b_{k,4}^{}b_{k,3}^{}b_{k,3}^{})`$
$$+i2\sqrt{2}\mathrm{sin}(2bk)(b_{k,3}^{}b_{k,4}^{}b_{k,4}^{}b_{k,3}^{})]$$
(6)
and
$`H_2^{(2)}={\displaystyle \frac{J_2}{12}}{\displaystyle \underset{k}{}}[\mathrm{cos}(2bk)(2b_{k,4}^{}b_{k,4}^{}b_{k,3}^{}b_{k,3}^{})`$
$$i2\sqrt{2}\mathrm{sin}(2bk)b_{k,4}^{}b_{k,3}^{}+\text{h.c.}].$$
(7)
Here $`b=3.611\text{Å}`$ is the lattice constant of the high-temperature phase. Note the opposite sign in the dispersion of two triplets. It is straightforward to diagonalize $`H^{(LHP)}`$. We define $`c=(J_2/6)\mathrm{cos}(2bk)`$, $`2t=\mathrm{\Delta }_4^2+\mathrm{\Delta }_3^2+2c(2\mathrm{\Delta }_4\mathrm{\Delta }_3)`$ and $`s=\mathrm{\Delta }_3^2\mathrm{\Delta }_4^2+2c\mathrm{\Delta }_3\mathrm{\Delta }_4(2\mathrm{\Delta }_3\mathrm{\Delta }_4)2\mathrm{\Delta }_3\mathrm{\Delta }_4J_2^2/9`$. The dispersion $`\omega _\pm =\omega _\pm (k)`$ of the two magnon branches (each branch is three-fold degenerate) in LHP-approximation is then
$$\omega _\pm ^2=t\pm \sqrt{t^2s}.$$
(8)
We have included the results for the magnon gap in Fig. 2. For large ratios $`J^{}/J`$ the LHP-result becomes asymptotically exact, in this limit it is equivalent to perturbation theory in $`J_2`$. In the LHP-approximation the transition to the Haldane phase is signaled by a vanishing of the energy gap, the crossover cannot be described by the cluster-operator theory.
In Fig. 3 we present the magnon-dispersion Eq. (8) for $`J^{}=2J`$ and compare the LHP-results (lines) with an exact-diagonalization study of a system with 16 sites (filled circles) . The agreement is very good, due to the large gap and (correspondingly) small correlation length. Note that the low-lying magnon, which corresponds to $`\psi _4`$ (see inset of Fig. 2), has its minimum at $`k=\pi /(2b)`$.
In Fig. 4 we present the magnon-dispersion Eq. (8) for $`J^{}=0.8J`$ which is closer to the transition to the Haldane phase. The agreement with the exact diagonalization and the DMRG data is not good, since the precursors to the Haldane phase are not included in the cluster-operator theory. The low-lying magnon, which corresponds to $`\psi _3`$ (see inset of Fig. 2), has its minimum now at $`k=0`$ and $`k=\pi /b`$ and a maximum at $`k=\pi /(2b)`$, as measured by neutron scattering . The cluster-operator theory substantially overestimates the size of the magnon dispersion relative to the exact-diagonalization result near to the Haldane phase. The physical reason for this discrepancy can be understood: The lattice constant of the effective spin-1 chain in the Haldane phase is $`2b`$ and the minimum of the magnon dispersion is therefore at $`\pi /(2b)`$ in the Haldane phase . It changes therefore at the crossover from the cluster phase and the Haldane phase. This change in the location of the gap is not included in the cluster-operator theory.
Discussion The exchange constant along $`b`$ is $`J529560\text{K}`$ in the high-temperature phase of NaV<sub>2</sub>O<sub>5</sub> and the inter-ladder coupling is probably very small, a $`J^{}/J1/45`$ has been found in an analysis of the magnon-dispersion for $`T<T_C`$ in a model with zig-zag charge order . This small ratio is consistent with the very small coupling along $`a`$ found in a LDA-study . There are, however, two reasons why $`J^{}`$ might be larger in the low-temperature phase. (a) As noted by Mack and Horsch , there is a near cancellation for $`T>T_C`$ in between paths with intermediate singlet and triplet states and energies $`E_{s/t}`$: $`J^{}=t_{xy}^2(1/E_s1/E_t)`$, where $`t_{xy}`$ is the V-V hopping matrix element in $`a`$ direction. A corresponding calculation for $`T<T_C`$ in the phase shown in Fig. 1 yields $`J^{}=2t_{xy}^2(1/E_s+1/U1/E_t)`$ ($`U`$ is the onsite Hubbard-$`U`$). (b) $`t_{xy}`$ might be substantially larger in the low-temperature phase, since the smallness of $`t_{xy}`$ for $`T>T_C`$ is a subtle band-structure effect . We have therefore scanned the complete phase diagram of the spin-cluster Hamiltonian in order to determine whether there exists a parameter range able to fit the neutron scattering data.
We have tried to reproduce, within the spin-cluster model, four known properties of NaV<sub>2</sub>O<sub>5</sub>: (i) The gap (averaged over $`k_a`$) is $`\mathrm{\Delta }_{min}=10\text{meV}`$. (ii) The maximum of the dispersion of the lowest magnon branch is at $`\pi /(2b)`$, the minimum at $`0`$ and $`\pi /b`$. (iii) The value of the maximum of the dispersion of the lowest magnon branch is $`\mathrm{\Delta }_{max}40\text{meV}`$ , i.e. the ratio is $`\mathrm{\Delta }_{max}/\mathrm{\Delta }_{min}4`$. (iv) The value of the coupling along $`b`$ is $`J441\text{K}=38\text{meV}`$ for $`T<T_{SP}`$ .
Condition (ii) implies that only the cluster-phase of Hamiltonian Eq. (1) with $`J^{}<J_1`$ is a candidate for the low-temperature phase of NaV<sub>2</sub>O<sub>5</sub>. This implies $`J_1/2<J^{}<J_1`$. Within the cluster-operator theory one obtains $`\mathrm{\Delta }_{max}/\mathrm{\Delta }_{min}=4`$ for values of $`J^{}`$ near to the gap closing. One needs consequently large coupling constants $`J`$ (see inset of Fig. 3). in order to reproduce $`\mathrm{\Delta }_{min}=10\text{meV}`$. We have evaluated the values of $`J^{}`$ and $`J`$ needed to reproduce the gap-ratio as a function of dimerization $`\delta `$ and find a minimum in $`J`$ for $`\delta =0.2`$ (see inset of Fig. 3). This minimum is $`J126\text{meV}`$, substantially larger than the experimental value $`J38\text{meV}`$. Note, that the cluster-operator theory overestimates the dispersion in this phase and underestimates the value of $`J`$ needed. We therefore conclude safely, that the model is not able to reproduce the measured magnon-dispersion of NaV<sub>2</sub>O<sub>5</sub> and that Eq. (1) is unlikely to be the appropriate model for the low-temperature phase of NaV<sub>2</sub>O<sub>5</sub>, at least in its one-dimensional version. It might be possible, in principle, that two-dimensional couplings change the scenario obtained in the present study, though we note, that an increase in dimensionality does, in general, reduce the size of a spin-gap.
We would like to acknowledge discussions with P. Lemmens and the support of the German Science Foundation, the BMBF and the Fonds der chemischen Industrie.
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# On the symmetry and uniqueness of solutions of the Ginzburg-Landau equations for small domains
## 1 Introduction and main results
In this paper, we study the properties of a superconducting cylinder submitted to an exterior magnetic field H<sub>0</sub> parallel to the axis of the cylinder. According to the Ginzburg-Landau theory of superconductivity, the sample is in a state that minimizes the following energy:
$$E_\kappa (\mathrm{\Psi },\text{A})=_\mathrm{\Omega }|(\frac{1}{\kappa }i\text{A})\mathrm{\Psi }|^2+\frac{1}{2}(|\mathrm{\Psi }|^21)^2+(\text{curl }\text{A}\text{H}\text{0})^2d\mathrm{\Omega }.$$
(1.1)
Here $`\mathrm{\Omega }`$ is a simply connected domain in $`\mathrm{IR}^2`$ with characteristic length $`d`$. As usual, $`\kappa `$ is the Ginzburg-Landau parameter that characterizes the type of superconductor, $`\text{A}(x,y)`$ is the vector potential, so that curl A is the magnetic field and $`\mathrm{\Psi }(x,y)`$ is the order parameter. Because $`\text{H}\text{0}=H_0\text{e}\text{z}`$ is along the $`z`$ axis, we can assume without loss of generality that $`\text{A}=(A_1,A_2,0)`$. For a detailed description of the model, one may refer to or .
The minimization process yields the classical Ginzburg-Landau equations (see for more details):
$$\{\begin{array}{cc}(\frac{1}{\kappa }i\text{A})^2\mathrm{\Psi }=\mathrm{\Psi }(|\mathrm{\Psi }|^21)\text{in}\mathrm{\Omega },\hfill & \\ \text{curl }\text{curl }\text{A}=\frac{i}{2\kappa }(\mathrm{\Psi }^{}\mathrm{\Psi }\mathrm{\Psi }\mathrm{\Psi }^{})+\text{A}|\mathrm{\Psi }|^2\text{in}\mathrm{\Omega },\hfill & \\ \frac{\mathrm{\Psi }}{n}i\mathrm{\Psi }\text{A}\text{n}=0\text{on}\mathrm{\Omega },\hfill & \\ \text{curl }\text{A}\times \text{n}=\text{H}_0\times \text{n}\text{on}\mathrm{\Omega }.\hfill & \end{array}$$
$`(GL)`$
It is common to use the gauge where $`\text{ div }\text{A}=0`$ and $`\text{A}\text{n}=0`$ on $`\mathrm{\Omega }`$. Notice that this system has a special set of solutions with $`\mathrm{\Psi }0`$ and $`\text{curl }\text{A}\text{H}_0`$, called normal solutions. They correspond to a case where superconductivity is destroyed. In the following, we will be concerned with solutions which are not in this set.
In this problem, there are 3 variable parameters: $`\kappa `$, $`H_0`$ and $`d`$ the characteristic size of $`\mathrm{\Omega }`$. According to the values of these parameters, the properties of solutions (existence of non normal solutions, number, symmetry) change. It has been the purpose of various authors to study the properties of solutions in some asymptotic limits: for instance, Giorgi and Phillips have proved that for $`\mathrm{\Omega }`$ and $`\kappa `$ fixed and $`H_0`$ large enough, then the only solution is the normal solution. In the case where $`d`$ is fixed and $`\kappa `$ is large, Sandier and Serfaty have proved that the solutions can have a lot of vortices for a certain range of magnetic field. There have been many other contributions in other asymptotic regimes.
In this paper, we are interested in the case where $`d`$, the characteristic length of $`\mathrm{\Omega }`$ is small, in particular compared to $`1/\kappa `$. In this setting, it is expected that $`\mathrm{\Psi }`$ is almost constant and has no zero, that is no vortices, since a vortex is of size $`1/\kappa `$. This result is often used in the physics literature, but we did not know any rigorous proof. This is what we show:
###### Theorem 1
Assume that $`D`$ is a fixed simply connected bounded domain and let $`\mathrm{\Omega }=dD`$. For any $`d_0>0`$, there exists $`d_1>0`$, such that if $`d<\mathrm{min}(d_0,d_1/\kappa )`$, for any $`H_0=H_0(d,\kappa )`$, then any solution of $`(GL)`$ which is not a normal solution is such that $`\mathrm{\Psi }`$ has no zero. Moreover, $`|\mathrm{\Psi }/|\mathrm{\Psi }|_{\mathrm{}}1|C\kappa d`$.
We let $`H_0`$ be a function of $`d`$ and $`\kappa `$ because it is expected that a superconducting solution exists up to fields $`H_0`$ of the order of $`C/d`$ as we will explain in more details later. Additionally, we prove a symmetry and uniqueness result when the domain is a small disc:
###### Theorem 2
Let $`\mathrm{\Omega }`$ be a disc of radius $`d`$. There exist constants $`d_0`$ and $`d_1`$, such that if $`d<\mathrm{min}(d_0,d_1/\kappa )`$, for any $`H_0=H_0(d,\kappa )`$, then any solution of $`(GL)`$ is radially symmetric, that is $`\mathrm{\Psi }(x,y)=\mathrm{\Psi }(r)`$ and $`\text{A}(x,y)=A(r)\text{e}_\theta `$, where $`r=\sqrt{x^2+y^2}`$.
###### Theorem 3
Under the hypotheses of Theorem 2, if there exists a non normal solution $`(\mathrm{\Psi },\text{A})`$ of $`(GL)`$, then this solution is unique up to multiplication of $`\mathrm{\Psi }`$ by a constant of modulus 1.
###### Corollary 4
Under the hypotheses of Theorem 2, there exists $`H_{}=H_{}(d,\kappa )`$, with $`lim_{d0}dH_{}(d,\kappa )=2\sqrt{2}`$ and for $`H_0<H_{}`$, there exists a unique non normal solution of $`(GL)`$, while for $`H_0H_{}`$, the only solutions are the normal solutions.
We make a change of variables $`x^{}=x/d`$, $`y^{}=y/d`$ so that the new variable lies in a domain of unit size $`D`$. We also define $`\text{B}=\kappa d\text{A}`$ and $`h_0=\kappa d^2H_0`$. In the following, we will assume additionnally that
$$\text{ div }\text{B}=0\text{in}\mathrm{\Omega }\text{and}\text{B}\text{n}=0\text{on}\mathrm{\Omega }.$$
(1.2)
The equations then become
$$\{\begin{array}{cc}(i\text{B})^2\psi =\kappa ^2d^2\psi (|\psi |^21)\text{in}D,\hfill & \\ \mathrm{\Delta }\text{B}=d^2\left(\frac{i}{2}(\psi ^{}\psi \psi \psi ^{})+\text{B}|\psi |^2\right)\text{in}D,\hfill & \\ \frac{\psi }{n}=0\text{on}D,\hfill & \\ \text{curl }\text{B}\times \text{n}=\text{h}_0\times \text{n}\text{on}D.\hfill & \end{array}$$
$`(GL_d)`$
Note that another way of writing the equation for B is
$$\mathrm{\Delta }\text{B}=d^2(i\psi ,\psi i\text{B}\psi ),$$
(1.3)
where $`(.,.)`$ is the real part of the scalar product in $`\text{ }\mathrm{C}`$.
We allow $`h_0`$ to vary with $`d`$ and $`\kappa `$ but we will prove that if there exists a solution, then in fact $`h_0`$ is bounded, that is $`H_0`$ is bounded by $`C/\kappa d^2`$.
The proof consists in obtaining a priori estimates for the solutions $`(\psi ,\text{B})`$. This is done in Section 2. In Section 3, we use these a priori estimates and the Poincaré inequality to derive that $`\psi `$ is nearly constant, hence has no zero. This will prove Theorem 1. Then, in Section 4, we define the functions
$$\stackrel{~}{\psi }(x,y)=\psi (x,y),\stackrel{~}{\text{B}}(x,y)=\left(\begin{array}{cc}B_1(x,y)\hfill & \\ B_2(x,y)\hfill & \end{array}\right),$$
$$w(x,y)=\frac{1}{|\psi |_{\mathrm{}}}(\psi (x,y)\stackrel{~}{\psi }(x,y)),\text{z}(x)=\text{B}(x,y)\stackrel{~}{\text{B}}(x,y),$$
(1.4)
which satisfy elliptic PDE’s with small right hand side terms. Then we use that $`\psi `$ is nearly constant and hence we get that $`w`$ and z are identically zero. This will prove Theorem 2. In Section 5, we obtain the uniqueness result proving that any solution is necessary a local minimizer of the energy. Finally in Section 6, we show how our proof can be adapted to the one dimensional case to provide a symmetry result in this setting and in Section 7 we summarize our results in terms of bifurcation curves and prove Corollary 4.
Remark: If $`\mathrm{\Omega }`$ is simply connected and symmetric in the $`y`$ direction, then our proof also gives that $`\psi `$ and B are symmetric, in the sense that $`\stackrel{~}{\psi }=\psi `$ and $`\stackrel{~}{\text{B}}=\text{B}`$.
We note that in the case with no magnetic field, Bethuel, Brezis and Helein have proved the symmetry of solutions when the domain is small using the Poincaré inequality. The radial symmetry of local minimizers has also been studied in for a different system with no magnetic field. Here the proof is more involved since we have to deal with the system with applied magnetic field.
Let us point out that the uniqueness proof of Section 5 is inspired by where it is obtained that any radial solution with a giant vortex of degree $`N`$ at the origin is a local minimizer of the energy.
## 2 A priori estimates
###### Proposition 5
Fix $`p>1`$. Assume that $`(psi,\text{B})`$ is a non normal solution of $`(GL_d)`$ and (1.2). For all constants $`d_0`$ and $`d_1`$, if $`d<\mathrm{min}(d_0,`$ $`d_1/\kappa )`$, then $`\psi `$ and B are bounded in $`W^{2,p}(D)`$ by constants independent of $`d`$ and $`\kappa `$. Moreover, for fixed $`\kappa `$,
$$\underset{d0}{lim}\text{ h}_0=0\text{and}\underset{d0}{lim}\text{B}_{W^{2,p}(D)}=0.$$
(2.1)
The goal of this section consists in proving Proposition 5.
### 2.1 First estimates for $`\psi `$
First recall from Du-Gunzburger-Peterson that
$$|\psi |1\text{a.e.}$$
(2.2)
Next we define
$$u(x,y)=|\psi (x,y)|^2\text{and}𝚷=i\text{B}.$$
(2.3)
One can compute easily that
$$\frac{1}{2}\mathrm{\Delta }u+\kappa ^2d^2u(1u)=|𝚷\psi |^2.$$
(2.4)
Integrate (2.4) in $`D`$, use the boundary condition $`\psi /n=0`$ and get the key estimate for $`\psi `$:
$$𝚷\psi _{L^2(D)}\kappa d\psi _{L^2(D)}.$$
(2.5)
### 2.2 Estimates for B
We decompose $`\text{B}=h_0\text{b}_0+\text{b}_r`$ where $`\text{b}_0`$ is chosen such that
$$\text{curl }\text{curl }\text{b}_0=0\text{and}\text{ div }\text{b}_0=0\text{in}D,\text{curl }\text{b}_0=\text{e}_z\text{and}\text{b}_0\text{n}=0\text{on}D,$$
and $`\text{b}_r=\text{B}h_0\text{b}_0`$.
A particular choice for $`\text{b}_0`$ when $`\mathrm{\Omega }`$ is a disc is $`\text{b}_0=\frac{1}{2}(y,x,0)`$. Note that in any case $`\text{b}_0`$ and $`\text{b}_r`$ are vectors in the $`xy`$ plane and their $`curl`$ are in the $`z`$ direction. Making the difference between the equations for B and $`\text{b}_0`$, we obtain the equation for $`\text{b}_r`$
$`\text{curl }\text{curl }\text{b}_r=d^2(i\psi ,𝚷\psi )\text{and}\text{ div }\text{b}_r=0\text{in}D,`$ (2.6)
$`\text{curl }\text{b}_r=0\text{and}\text{b}_r\text{n}=0\text{on}D.`$ (2.7)
Since $`\text{ div }\text{b}_r=0`$ and $`\text{b}_r\text{n}=0`$ on $`D`$, then $`\text{curl }\text{b}_r_{L^p(D)}`$ is a norm in $`W^{1,p}(D)`$, that is
$$\text{b}_r_{W^{1,p}(D)}C\text{curl }\text{b}_r_{L^p(D)}.$$
(2.8)
Now Cauchy-Schwarz inequality and (2.6) imply that
$$|\text{curl }\text{curl }\text{b}_r|_2d^2|𝚷\psi |_2.$$
(2.9)
Note that $`|.|_2`$ is the 2-norm for vectors in $`\mathrm{IR}^2`$. From (2.9), we get the $`W^{2,2}`$ bound for $`\text{b}_r`$. Thus, because of (2.5), (2.9) and the boundary condition (2.7), it follows from the Poincaré inequality that
$$\text{curl }\text{b}_r_{W^{1,2}(D)}C\kappa d^3.$$
(2.10)
We gather (2.8) and (2.10) to obtain the key estimate for $`\text{b}_r`$:
$$\text{b}_r_{W^{2,2}(D)}C\kappa d^3,$$
(2.11)
where $`C`$ is independent of $`d`$ and $`\kappa `$. In order to get the bounds for B, we only need to prove that $`lim_{d0}h_0=0`$. Let us assume that $`h_0`$ is bounded below by some constant $`m`$. Then we use an estimate by Giorgi and Phillips (Lemma 2.8 p.349): there exists $`C_2`$ (which depends on $`m`$) such that
$$C_2h_0_D|\psi |^2_D|\psi ih_0\text{b}_0\psi |^2.$$
(2.12)
In order to bound the right hand side of (2.12), we write $`h_0\text{b}_0=\text{B}\text{b}_r`$ and use (2.5) and (2.11) to get
$$(C_2h_0)^{1/2}(_D|\psi |^2)^{1/2}\kappa d(_D|\psi |^2)^{1/2}(1+C_1d^2).$$
(2.13)
So we obtain a contradiction with the lower bound for $`h_0`$ when $`d`$ is small enough. Hence B is bounded in $`W^{2,2}`$, independently of $`d`$ and $`\kappa `$.
Once we have bounded B, the equation for $`\psi `$ can be thought of as a linear elliptic equation with bounded coefficients, so that the classical elliptic estimates imply bounds for $`\psi `$.
Finally, in order to get the $`W^{2,p}`$ bound for B, that is for $`\text{b}_r`$, we bootstrap equation (2.6) (see ). Thus we obtain that
$$\text{curl }\text{b}_r_{W^{2,p}(D)}Cd^2.$$
(2.14)
and (2.1) is true for B.
## 3 $`\psi `$ is nearly constant
###### Proposition 6
For all $`d_0>0`$, there exists $`d_1>0`$, such that if $`d<\mathrm{min}(d_0,d_1/\kappa )`$, then $`\psi /|\psi |_{\mathrm{}}_{\mathrm{}}`$ and $`|\psi /|\psi |_{\mathrm{}}1|`$ are small, thus $`\psi `$ has no zero.
We define $`\varphi =\psi /|\psi |_{\mathrm{}}`$, then the equation for $`\varphi `$ can be written as
$$\mathrm{\Delta }\varphi 2i\text{B}\varphi =\kappa ^2d^2\varphi (|\psi |_{\mathrm{}}^2|\varphi |^21)+|\text{B}|^2\varphi $$
and for $`p<2`$
$$|\text{B}\varphi |_p|\text{B}𝚷\varphi |_p+|\text{B}^2\varphi |_p|\frac{1}{|\psi |_{\mathrm{}}}𝚷\psi |_2|\text{B}|_{\frac{2p}{2p}}+|\text{B}^2|_p.$$
(3.1)
Since $`|\varphi |1`$, and because of (2.5) and the previous estimates for B, we see that the right hand side of (3.1) is small when $`d`$ is small. Hence we go back to the equation for $`\varphi `$ and the Agmon Douglis Nirenberg estimates imply a bound for $`\varphi `$ in $`W^{2,p}`$ for any $`p<2`$. Hence by Sobolev embedding, $`\varphi `$ is bounded in $`L^p`$ for any $`p`$ and going back to the equation for $`\varphi `$, we derive a bound for $`\varphi `$ in $`W^{2,p}`$ for any $`p`$, hence in $`C^1`$. In particular, $`\varphi `$ is equicontinuous. Now, when $`d`$ tends to 0, we know that $`\text{B}_{\mathrm{}}`$ tends to 0, so if we multiply the equation for $`\varphi `$ by $`\varphi ^{}`$ and integrate, we find that $`\varphi _{L^2(D)}`$ is small. Thus, since $`\varphi `$ is equicontinuous, we see that $`\varphi _{\mathrm{}}`$ tends to 0 with $`d`$, which proves the first part of Proposition 6.
In fact we can get a precise estimate for the smallness of $`|\psi /|\psi |_{\mathrm{}}1|`$. Take the equation (2.4), multiply by $`u`$ and integrate to get
$$_D|u|^2\kappa ^2d^2_Du^2(1u).$$
(3.2)
Next we define $`v=u/|\psi |_{\mathrm{}}^2`$. Then (3.2) implies that
$$_D|v|^2\kappa ^2d^2|D|.$$
We call $`v_{mean}`$ the mean value of $`v`$ in $`D`$. The previous estimate and the Poincaré inequality yield
$$vv_{mean}_{L^2(D)}C_0\kappa d.$$
(3.3)
Let us call $`r`$ an upper bound for $`v`$, independent of $`d`$ and $`\kappa `$. We have seen in the first part of the proof that it exists.
Let us now prove that
$$|v(x)v_{mean}|\alpha xD,\text{where}\alpha ^2=16C_0r\kappa d/\sqrt{\pi },$$
(3.4)
where $`C_0`$ comes from (3.3). Assume that (3.4) is not true. Then there exists a point $`x_0`$ where this does not hold. Set $`\eta =\alpha /4r`$. Note that when $`\kappa d`$ is small, then $`\eta `$ is small. If $`\mathrm{dist}(x_0,D)2\eta `$, then for $`\kappa d`$ small enough, there always exists a point $`x_1`$ in $`D`$ with $`\mathrm{dist}(x_1,D)2\eta `$ and $`\mathrm{dist}(x_1,x_0)2\eta `$. If $`\mathrm{dist}(x_0,D)2\eta `$ then we set $`x_1=x_0`$. Now we have $`|v(x)v_{mean}|>\alpha /4`$ for $`|xx_1|\eta `$. Define $`D_1=D(x_1,\eta )`$, then
$$vv_{mean}_{L^2(D_1)}>\eta \sqrt{\pi }\alpha /4.$$
This provides a contradiction with the definition of $`\alpha `$ and (3.3).
Now recall that $`v`$ is necessarily 1 for at least one point in $`D`$. So we use (3.4) to get
$$|\psi (x)/|\psi |_{\mathrm{}}1|C\kappa dxD.$$
This implies in particular that for small $`\kappa d`$, $`\psi `$ is never equal to zero.
## 4 Radial symmetry of the solutions in the case of a disc
In this Section, we assume that $`\mathrm{\Omega }`$ is a disc of radius $`d`$. We are going to use the previous a priori estimates to obtain radial symmetry. We define the functions $`w`$ and z as in (1.4). They satisfy the following equations
$`\mathrm{\Delta }w2i\text{B}w{\displaystyle \frac{2i}{|\psi |_{\mathrm{}}}}\text{z}\stackrel{~}{\psi }=\kappa ^2d^2w(|\psi |^2+\stackrel{~}{\psi }(\psi +\stackrel{~}{\psi })1)`$ (4.3)
$`+w|\text{B}|^2+{\displaystyle \frac{\stackrel{~}{\psi }}{|\psi |_{\mathrm{}}}}\text{z}(\text{B}+\stackrel{~}{\text{B}})\text{in}D,`$
$`{\displaystyle \frac{w}{n}}=0\text{on}D,`$
$`\text{curl }\text{curl }\text{z}=d^2\left[(i|\psi |_{\mathrm{}}w,𝚷\psi )+(i\stackrel{~}{\psi },|\psi |_{\mathrm{}}wi\text{B}|\psi |_{\mathrm{}}wi\stackrel{~}{\psi }\text{z})\right]`$
$`\text{and}\text{ div }\text{z}=0\text{in}D,`$
$`\text{z}\text{n}=0\text{and}\text{curl }\text{z}\times \text{n}=0\text{on}D.`$ (4.4)
###### Proposition 7
There are constants $`d_0`$ and $`d_1`$ such that, if $`d<\mathrm{min}(d_0,`$ $`d_1/\kappa )`$, then $`w\text{z}0`$.
Let us multiply (4.3) by $`w`$ and integrate. The first term only gives an integral of the gradient of $`w`$ because of the boundary condition (4.3). We recall that $`|\psi |1`$ and $`|w|2`$ to obtain
$`{\displaystyle _D}|w|^2`$ $``$ $`4\text{B}_{\mathrm{}}{\displaystyle _D}|w|+(4\kappa ^2d^2+\text{B}_{\mathrm{}}^2){\displaystyle _D}w^2`$
$`+4({\displaystyle \frac{\psi }{\psi _{\mathrm{}}}}_{\mathrm{}}+\text{B}_{\mathrm{}}){\displaystyle _D}|\text{z}|.`$
Then we use our previous estimates for $`\psi `$ and B in Proposition 5 and 6. In particular we recall that $`\text{B}_{\mathrm{}}`$ and $`\frac{1}{|\psi |_{\mathrm{}}}\psi _{\mathrm{}}`$ are small when $`d`$ is small to get: for all $`\epsilon >0`$, if $`d`$ is small enough, then
$$w_2\epsilon w_2+\epsilon \text{z}_2.$$
Because of the definition of $`w`$, one can see easily that $`w`$ has mean value zero. We can use the Poincaré inequality to get
$$\mu w_2\epsilon w_2+\epsilon \text{z}_2,$$
(4.5)
where $`\mu >0`$ is independent of $`d`$ and $`\kappa `$. Thus
$$w_2\frac{\epsilon }{\mu \epsilon }\text{z}_2\text{z}_2,$$
(4.6)
provided $`\epsilon /(\mu \epsilon )`$ is less then 1.
Similarly, we multiply equation (4.3) by z, integrate and use that the boundary condition implies $`\text{curl }\text{z}=0`$ on $`D`$ to get
$$_D|\text{curl }\text{z}|_2^2d^2_D|\text{z}|_2|\mathrm{\Pi }\psi +w+\text{B}+\text{z}|_2].$$
We can apply the previous estimates which give bounds on $`\psi `$ and B to derive
$$\text{z}_2\epsilon \text{z}_2.$$
(4.7)
Then we write $`\text{z}=\text{z}_m+\text{z}_t`$ where $`\text{z}_m`$ is the mean value of z (so $`\text{z}_m`$ is a vector) and $`\text{z}_t`$ has mean value zero. Again we use the Poincaré inequality for $`\text{z}_t`$ and equation (4.7) to get
$$\text{z}_t_2\frac{\epsilon }{\mu \epsilon }\text{z}_m_2\text{and}w_2\frac{\epsilon \mu }{(\mu \epsilon )^2}\text{z}_m_2.$$
(4.8)
Now we integrate (4.3) over $`D`$, use the boundary conditions which imply that the curl curl term contribution vanishes and we get
$$_D\frac{|\stackrel{~}{\psi }|^2}{|\psi |_{\mathrm{}}^2}\text{z}=_D(iw,𝚷\psi )+(i\stackrel{~}{\psi },wi\text{B}w).$$
(4.9)
This yields
$$|\text{z}_m|_D\frac{|\psi |^2}{|\psi |_{\mathrm{}}^2}w_2𝚷\psi _2+\psi _2wi\text{B}w_2+_D|\text{z}_t|C\epsilon |\text{z}_m|.$$
The last inequality comes from (4.8) and the bounds for $`\psi `$ and B in Section 3. If $`\text{z}_m`$ is not zero, the boundedness of $`𝚷\psi `$, B and (4.8) imply that
$$_D|\psi |^2/|\psi |_{\mathrm{}}^2$$
is small. This provides a contradiction with Proposition 6, where we have proved that $`|\psi |/|\psi |_{\mathrm{}}`$ is nearly 1 for small $`d`$.
Let us complete the proof of Theorem 2. Indeed, once we know that $`w`$ and z are identically zero, then it implies as in that $`\psi `$ and $`|\text{B}|`$ are radial, where $`|.|`$ is the modulus. Thus, we see that either B vanishes everywhere on the circle $`|(x,y)|=r`$ and there is nothing to prove, or nowhere on this circle. In this case, we can argue that $`\text{B}(r,\theta )`$ is in the direction of $`𝐞_\theta =(\mathrm{sin}\theta ,\mathrm{cos}\theta )`$. Indeed, since $`\text{z}=0`$, it implies in particular that $`B_1(r,0)=0`$, hence $`\text{B}(r,0)`$ is in the direction of $`(0,1)`$. Then we use a reflection in the axis $`T`$ in the direction of $`(\mathrm{cos}\theta ,\mathrm{sin}\theta )`$ and derive that the corresponding $`w`$ and z are zero. This implies that $`\text{B}(r,\theta )𝐞_r=0`$. Since $`B_1(r,0)=0`$, a continuity argument implies that $`\text{B}(r,\theta )=B(r)𝐞_\theta `$, as we had claimed.
Let us show that one can choose $`\psi `$ to be real. Since $`\psi `$ is radial and has no zero, we can write $`\psi =f(r)e^{i\varphi (r)}`$, where $`f`$ is real. The estimate of Proposition 6 yields that $`\varphi ^{}(r)`$ is small for small $`d`$. Thus, if we change gauge from $`(\psi ,\text{B})`$ to $`(f,\text{Q})`$ where $`\text{Q}=\text{B}\varphi ^{}(r)\text{e}_r`$, then $`(f,\text{Q})`$ is a solution of $`(GL_d)`$, $`\text{Q}\text{n}=0`$ on $`\mathrm{\Omega }`$ and $`|\text{Q}|_{\mathrm{}}`$ is small. Notice that in this section, we have not explicitely used that $`\text{ div }\text{B}=0`$. We have only used that $`|\text{B}|_{\mathrm{}}`$ is small, which was proved in Proposition 5 under the hypothesis $`\text{ div }\text{B}=0`$. Now, since $`|\text{Q}|_{\mathrm{}}`$ is small, this is enough to do the proof of this section with $`(f,\text{Q})`$ instead of $`(\psi ,\text{B})`$. It yields that Q is along $`\text{e}_\theta `$ and that $`\varphi ^{}=0`$. Thus one can choose $`\psi `$ up to multiplication by a constant of modulus 1. This finishes the proof of Theorem 2.
## 5 Uniqueness for the ball
In this section, we assume again that $`D`$ is a ball. We have proved in the previous section that we can assume that $`\psi =f(r)`$ is real and $`\text{B}=\text{Q}=Q(r)\text{e}_\theta `$. Then the Ginzburg-Landau energy in this case is
$$E(f,Q)=_Df_{}^{}{}_{}{}^{2}+f^2Q^2+\frac{\kappa ^2d^2}{2}(f^21)^2+\frac{1}{d^2}(\frac{1}{r}(rQ)^{}h_0)^2.$$
(5.10)
We also have $`\text{ div }\text{Q}=0`$ then $`\frac{1}{r}(rQ)^{}_{L^2}`$ is an $`H^1`$ norm. We are going to prove that any radial solution of $`(GL_d)`$ is a local minimizer of the energy (5.10). If $`(f,Q)`$ is such a solution, then for any $`g`$ and $`\text{P}=P(r)\text{e}_\theta `$ in $`H^1`$,
$`0`$ $`=`$ $`E^{}(f,Q)(g,P)={\displaystyle _D}f^{}g^{}+QPf^2+Q^2fg`$ (5.11)
$`+`$ $`{\displaystyle \frac{1}{d^2}}({\displaystyle \frac{1}{r}}(rP)^{})({\displaystyle \frac{1}{r}}(rQ)^{}h_0)+\kappa ^2d^2(f^21)fg.`$
Moreover
$`E^{\prime \prime }(f,Q)(g,P)`$ $`=`$ $`{\displaystyle _D}g_{}^{}{}_{}{}^{2}+g^2Q^2+4fgQP+P^2f^2`$ (5.12)
$`+`$ $`{\displaystyle \frac{1}{d^2}}({\displaystyle \frac{1}{r}}(rP)^{})^2+\kappa ^2d^2g^2(3f^21)fg`$
###### Proposition 8
Let $`(f,Q)`$ be a radial solution of $`(GL_d)`$ with $`d`$ and $`d\kappa `$ small. Then $`E^{\prime \prime }(f,Q)(g,P)C(g/f_{H^1}+P_{H^1})`$, where $`C`$ depends on $`d`$, $`\kappa d`$ and $`f_{\mathrm{}}`$.
The proof follows ideas from . First compute $`E^{}(f,Q)(g^2/f,0)`$. This yields
$$0=_Df^{}(\frac{g^2}{f})^{}+g^2Q^2+\kappa ^2d^2(f^21)g^2.$$
Hence
$$E^{\prime \prime }(f,Q)(g,P)=_D((\frac{g}{f})^{})^2f^2+4fgQP++P^2f^2+\frac{1}{d^2}(\frac{1}{r}(rP)^{})^2+2\kappa ^2d^2g^2f^2.$$
We write
$$4fgQP+\kappa ^2d^2g^2f^2=(\kappa dfg+\frac{2}{\kappa d}PQ)^2\frac{4}{\kappa ^2d^2}P^2Q^2.$$
Hence
$$E^{\prime \prime }(f,Q)(g,P)_D((\frac{g}{f})^{})^2f^2+\kappa ^2d^2g^2f^2+\frac{1}{d^2}(\frac{1}{r}(rP)^{})^2\frac{4}{\kappa ^2d^2}P^2Q^2.$$
Now we use that for $`d`$ small enough, $`Q_{\mathrm{}}`$ is small hence $`8_DP^2Q^2\kappa ^2_D(\frac{1}{r}(rP)^{})^2`$, thus the estimate of the Proposition holds.
This Proposition yields that any solution of $`(GL_d)`$ is non degenerate. We are going to prove uniqueness using the bifurcation curve. More precisely, any radial solution of $`(GL_d)`$ solves
$$\{\begin{array}{cc}f^{\prime \prime }+\frac{1}{r}f^{}=f(f^2+Q^21)\hfill & \\ q^{\prime \prime }+\frac{1}{r}q^{}\frac{1}{r^2}q=f^2q,\hfill & \end{array}$$
(5.13)
with the boundary conditions
$$f^{}(0)=f^{}(1)=0\text{and}q(0)=0,(q^{}+\frac{1}{r}q)(1)=h_0.$$
(5.14)
We use shooting techniques, that is solve (5.13) with the initial conditions
$$f(0)=\beta ,f^{}(0)=0,q(0)=0,q^{}(0)=\alpha .$$
(5.15)
Then the same kind of proof as in Kwong’s paper which applies the Sturm comparison principle to equation (5.13) shows that $`f(r;\beta ,\alpha )`$, $`q(r;\beta ,\alpha )`$ and $`q^{}(r;\beta ,\alpha )`$ are increasing functions of $`\beta `$ and $`\alpha `$. Moreover, for each $`\beta `$ in $`(0,1)`$, there exists a unique $`\alpha `$ such that $`f^{}(1;\beta ,\alpha )=0`$, that is the boundary condition is satisfied by $`f`$. Then
$$h(\beta )=(q^{}+\frac{1}{r}q)(1;\beta ,\alpha (\beta ))$$
is a continuous function of $`\beta `$ with limit 0 as $`\beta `$ tends to 1 and a finite limit as $`\beta `$ tends to 0. The nondegeneracy result of Proposition 8 implies that this curve is decreasing and it provides uniqueness.
Note that uniqueness could also be proved by combining the compactness of solutions and Proposition 8 with Crandall Rabinowitz bifurcation Theorem. The Crandall Rabinowitz bifurcation Theorem ensures that only one branch comes off the normal solutions. Then the non degeneracy result of Proposition 8 and the implicit function Theorem guarantee that the value $`\beta `$ continues locally uniquely as $`h_0`$ is varied. This ensures uniqueness.
## 6 The 1-dimensional case
When the superconducting material is an infinite slab of thickness $`2d`$ between the planes $`x=d`$ and $`x=d`$, it is usual to assume that both $`\mathrm{\Psi }`$ and A are uniform in the $`y`$ and $`z`$ direction, and that the exterior magnetic field is tangential to the slab, that is $`\text{H}\text{0}=`$(0,0,$`H_0`$). A suitable gauge can then be chosen so that $`\mathrm{\Psi }=f(x)`$ is a real function, and $`\text{A}=q(x)\text{e}_y`$, where e<sub>y</sub> is the unit vector along the $`y`$ direction (see for more details). The model can then be simplified to a system of 2 coupled ODE’s $`(f,q)`$ satisfy
$$\{\begin{array}{cc}\frac{1}{\kappa ^2}f^{\prime \prime }=f(f^2+q^21)\text{in}(d,d),\hfill & \\ f^{}(\pm d)=0,\hfill & \\ q^{\prime \prime }=qf^2\text{in}(d,d),\hfill & \\ q^{}(\pm d)=H_0.\hfill & \end{array}$$
$`(gl_d)`$
In this setting, there are two types of solutions: symmetric solutions where $`f`$ is even and $`q`$ is odd, and asymmetric solutions. A complete numerical study of the number and symmetry of these solutions has been done in .
Using the same techniques as in the previous sections, we can prove
###### Theorem 9
There exist constants $`d_0`$ and $`d_1`$, such that if $`d<\mathrm{min}(d_0,`$ $`d_1/\kappa )`$, then any solution of $`(gl_d)`$ is symmetric, that is $`f`$ is even and $`q`$ is odd.
This theorem together with the result of uniqueness for symmetric solutions proved in give a global uniqueness result for the solutions of $`(gl_d)`$ with small $`d`$.
The proof, as in the previous sections consists first in deriving a priori estimates for the functions $`f`$ and $`q`$. We recall from that a solution is such that $`f`$ has a unique maximum, which we call $`x_0`$, with $`\beta =f(x_0)=f_{\mathrm{}}`$, and $`q`$ is increasing with a unique zero $`x_1`$. A similar proof to what we did in Section 2 and 3 yields:
###### Proposition 10
There exist constants $`d_0`$ and $`d_1`$, such that if $`d<\mathrm{min}`$ $`(d_0,d_1/\kappa )`$, then
$$|\frac{f(x)}{\beta }1|C\kappa d,\text{and}|q(x)H_0x|Cd^2.$$
(6.1)
Using similar techniques as in , one can get more precise estimates, which in particular give a relation between $`\beta `$ and $`H_0`$.
###### Proposition 11
Let $`0<d<min(1,1/2\kappa )`$, then
$$\beta (1\frac{\kappa ^2}{2}(xx_0)^2)f(x)\beta x(d,d),$$
(6.2)
and there exists $`\alpha `$ such that
$`\alpha (xx_1)q(x)\alpha (xx_1)(1+2\beta d)x(x_1,d),`$ (6.3)
$`\alpha (xx_1)(12\beta d)q(x)\alpha (xx_1)x(d,x_1),`$ (6.4)
$`|x_1|2\kappa ^2d^3+2d^2,`$ (6.5)
$`{\displaystyle \frac{\sqrt{3(1\beta ^2)}}{\sqrt{1+3\beta d}}}\alpha d{\displaystyle \frac{\sqrt{3(1\beta ^2)+2d^2\kappa ^2}}{\sqrt{1d^2\kappa ^2/2}}}`$ (6.6)
These estimates mean that $`f`$ is nearly constant and that $`q`$ is nearly equal to $`\alpha x`$, because its zero $`x_1`$ is very small and recall from Proposition 10 that $`\alpha `$ and $`H_0`$ are very close. We will not give the proof of Proposition 11 since Proposition 10 is enough for our purposes.
The second step of the proof consists in defining the functions $`w`$ and $`z`$ as before (note that the distances have been non dimensionnalized by $`d`$ as in the previous sections:
$$w(x)=\frac{1}{f_{\mathrm{}}}(f(x)f(x))\text{and}z(x)=\frac{q(x)+q(x)}{d}.$$
(6.7)
We derive the equations satisfied by $`w`$ and $`z`$, multiply them respectively by $`w`$ and $`z`$ and use the estimates of Proposition 10 to obtain
$$|w^{\prime \prime }|\epsilon |w|+\epsilon |z|\text{and}|z^{\prime \prime }|\epsilon |w|+\epsilon |z|.$$
Then we use the Poincaré inequality to get the equivalent of (4.6) and (4.7).
Finally, we see that since $`_1^1f^2q=0`$, then
$$_1^1z\frac{f^2}{\beta ^2}=_1^1qw\frac{(\stackrel{~}{f}+f)}{\beta }.$$
(6.8)
As previously, we derive that if $`z_m`$ is different from zero, (6.8) provides a contradiction with (6.1), which means that $`f(x)/\beta `$ is nearly constant.
## 7 Remarks on the bifurcation curve
In order to describe the solutions of the Ginzburg-Landau equations, it is common to draw the bifurcation diagram. The natural bifurcation diagram, which appears in the physics literature, is to plot the infinity norm of $`\mathrm{\Psi }`$ (that we call $`\beta `$) against $`H_0`$. In our setting, recall that $`\mathrm{\Psi }`$ is nearly constant.
For the 1-dimensionnal model, Kwong has proved that for each $`\beta `$, there exists a unique $`H_0(\beta )`$ such that $`(gl_d)`$ has a symmetric solution. Recall that in our setting, we have proved that solutions are indeed symmetric. Aftalion and Troy have established a complete picture of the bifurcation curves in the different regimes of $`\kappa `$ and $`d`$. In the regime where $`d`$ is small compared to $`1/\kappa `$, then the curve $`\beta `$ against $`H_0`$ is decreasing from 1 to 0 as $`H_0`$ increases from 0 to $`C/d`$. In , it is proved that $`lim_{d0}dH_0(\beta )=\sqrt{3(1\beta ^2)}`$ uniformly with respect to $`\beta `$ in $`(0,1)`$. Thus, for $`H_0`$ larger than $`\sqrt{3}/d`$, the only solution is the normal solution. For $`H_0`$ smaller, there exists a unique solution which is not normal; it is symmetric and has $`\mathrm{\Psi }`$ or $`f`$ almost constant to $`\beta `$ between 0 and 1. The value of $`H_0`$ for which the normal solution loses its stability has been studied in detail by Bolley-Helffer .
In the 2 dimensionnal case, the same kind of behaviour holds. A natural question is to wonder when the superconducting solution exists. Using the result of Giorgi-Phillips, we have proved that if there exists a non normal solution then $`d^2H_0`$ tends to 0 with $`d`$. Hence, if $`H_0`$ is not small enough, the only solution is the normal solution. On the other hand, it is known that for $`H_0`$ small, there exists a superconducting solution close to 1.
When there is a superconducting solution, we have also proved that these solutions are such that $`\mathrm{\Psi }`$ is nearly constant. Let us call this constant $`\beta `$. Then there is a relation between the value of $`\beta `$ and $`H_0`$. This relation is obtained by integrating the equation for $`\mathrm{\Psi }`$ using the boundary condition. One gets
$$\mathrm{\Psi }(\mathrm{\Psi }^2+|A|^21)=0.$$
In the 1-d case, since $`q(x)`$ is equivalent to $`H_0x`$, we get
$$\underset{d0}{lim}d^2H_0^2=3(1\beta ^2),$$
as proved in . In the radial case, one can prove with a similar estimate to Section 2 that $`A(r)`$ is asymptotically $`H_0r/2`$, hence
$$\underset{d0}{lim}d^2H_0^2=8(1\beta ^2).$$
This gives an asymptotic limit of the bifurcation curve for small $`d`$ and finishes the proof of Corollary 4.
Ackowledgements: The authors would like to thank the referee for helpful suggestions and the second author would like to thank the Ecole Normale Supérieure for their support and hospitality.
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# Non-abelian quantum Hall states – exclusion statistics, 𝐾-matrices and duality –
## 1. Introduction
The fractional quantum Hall effect has led to the identification of new states of matter, which can be characterized as incompressible quantum fluids with off-diagonal long-range order (‘topological order’). After the initial discovery of the ‘principal Laughlin series’ of quantum Hall fluids at filling factor $`\nu =1/m`$, a large class of so-called abelian quantum Hall fluids has been identified, accounting for the rich spectrum of fractional quantum Hall plateaus that have been observed in the lowest Landau level.
The observation of a quantum Hall plateau at filling factor $`\nu =5/2`$ (see for a recent experiment) has made it clear that the traditional set of abelian quantum Hall states (which all share the property of having an odd-denominator filling factor) will not suffice for explaining the phenomena observed in the second Landau level. Prompted by this development, new categories of incompressible quantum fluids have been proposed. Among them are various ‘paired’ or ‘clustered’ states, such as the Pfaffian states first proposed by Moore and Read . The quasiparticles over these states satisfy what is called non-abelian braid statistics, and by abuse of language one speaks of ‘non-abelian quantum Hall states’.
The characteristic order of the abelian quantum Hall fluids should be viewed as ‘topological’ and it can be characterized by a collection of integer numbers, which together constitute a so-called $`K`$-matrix. Many of the low energy characteristics of the Hall fluid are encoded in this $`K`$-matrix and the quantum numbers of the elementary electron-type excitations. They include the filling factor $`\nu `$ and the spin Hall conductance $`\sigma `$. In addition the (fractional) quantum numbers of the various quasihole type excitations are determined by using (the inverse of) the $`K`$-matrix.
A systematic framework for the physical implications of the topological order embodied in the $`K`$-matrix is provided by effective Chern-Simons and conformal field theories for bulk and edge excitations, respectively. In a systematic treatment of the low energy dynamics, these theories arise as special limits of a unifying field theory for the low energy behaviour of quantum Hall systems .
It is well-known that bulk-excitations over fractional quantum Hall fluids satisfy fractional (anyonic) braid statistics. Closely related to this are the fractional exclusion statistics of both bulk and edge excitations . It has been observed that for abelian quantum Hall fluids, the (edge) statistics matrix $`𝐊`$ (in the sense of Haldane’s definition of exclusion statistics ) is closely related to the $`K`$-matrix.
The main purpose of the present paper is to present a $`K`$-matrix structure associated to specific series of non-abelian quantum Hall states. To this end, we study the exclusion statistics of edge excitations over these quantum Hall states, and identify from that analysis statistics matrices $`𝐊`$. We shall then argue that these same matrices can be viewed as $`K`$-matrices for these non-abelian quantum Hall states.
Our analysis here builds on earlier results published in . In we presented our present results in brief form, and elaborated on the physical meaning of the newly obtained $`K`$-matrices.
This paper is organized as follows. In Sect. 2 we briefly review the $`K`$-matrix theory for abelian quantum Hall states, and make a generalization in order to be able to treat spin singlet states. We continue in Sect. 3 by making the link to abelian exclusion statistics. We argue that the statistics matrix is (basically) given by the $`K`$-matrix. Also, we introduce an important notion of duality. In Sect. 4 we generalize this concept to the non-abelian case, where composites and pseudoparticles play a vital role. It is argued that the well known formulas for the physical quantities such as the filling factor, derived from the abelian $`K`$-matrix structure, still hold for the non-abelian $`K`$-matrices describing various clustered non-abelian quantum Hall states. Sect. 5 deals with the relation between the universal chiral partition function (UCPF) and exclusion statistics. In Sect. 6 and 7 the $`K`$-matrices for two classes of non-abelian clustered states are identified. Sect. 8 is reserved for discussions, while some of the more mathematical results, for instance on character formulas, are discussed in the appendices.
## 2. $`K`$-matrices for abelian quantum Hall states
In this section, we briefly review the $`K`$-matrix structure for abelian quantum Hall states. We do not derive, but merely state the results we need in this paper. For a more detailed review, see for instance .
The information needed to describe an abelian quantum Hall state can be encoded in the following way, henceforth referred to as the fqH data. The four important ‘objects’ which will do the job are the $`K`$-matrix (which will also play the role as statistics matrix), the so called charge and spin vectors, $`𝐭`$ and $`𝐬`$, respectively, and the angular momentum vector $`𝐣`$. A few remarks with respect to the notation of spin vectors need to be made at this point. In , the concept of a ‘spin vector’ was introduced. This ‘spin vector’ is in fact related to the angular momentum of the electrons on (for instance) the sphere and is needed to calculate the so-called shift. In our case we need to distinguish between this angular momentum vector and the vector containing the real spin of the particles. Therefore, we have denoted the angular momentum vector by $`𝐣`$, and the vector containing the spin quantum numbers by $`𝐬`$.
In order to have the possibility to connect the $`K`$-matrix with the statistics matrix (as we will do in the following sections), we will distinguish between the $`K`$-matrix for the ‘electron part’ and the ‘quasihole part’ of the theory. These will be denoted by $`𝐊_e`$ and $`𝐊_\varphi `$, respectively. The corresponding charge, spin and angular momentum vectors are $`𝐭_e`$, $`𝐭_\varphi `$, $`𝐬_e`$, $`𝐬_\varphi `$, $`𝐣_e`$, and $`𝐣_\varphi `$ in an obvious notation. In all the cases we considered, it is possible to choose a basis in which the $`K`$-matrices are just each others inverse, $`𝐊_e=𝐊_\varphi ^1`$.
As stated above, the $`K`$-matrices will play several roles in the theory. First of all, they couple the different Chern-Simons gauge fields which play a central role in a Lagrangian description of the quantum Hall states. In the abelian case, the Chern-Simons part of the Lagrangian for a system on a surface of genus $`g`$ reads as follows
$$_{\text{CS}}=\frac{1}{4\pi }ϵ^{\mu \nu \lambda }\left(𝐊_e^{ij}a_\mu ^i_\nu a_\lambda ^j+2𝐭_e^iA_\mu _\nu a_\lambda ^i+2𝐣_e^i\omega _\mu _\nu a_\lambda ^i+2𝐬_e^i\beta _\mu _\nu a_\lambda ^i\right),$$
(2.1)
where the fields $`a`$ are the Chern-Simons gauge fields. The Greek indices run over $`\{0,1,2\}`$, and the Roman indices over the number of channels. The first three terms in Eq. (2.1) are rather standard and described in, for instance, . The first term is the famous Chern-Simons term, the other three describe the couplings to various fields. The gauge field $`A_\mu `$ describes the electromagnetic field and $`\omega _\mu `$ is the ‘spin connection’ which gives rise to the curvature of the space on which the system is defined. To explain the last term, we briefly discuss the concept of the spin Hall conductance and the related spin filling factor $`\sigma `$ (see and references therein).
In general, one would define the spin conductance in the same way as the charge conductance, namely as a response to a certain field. In the case of a quantum Hall system, the role of the electric field is taken over by a gradient in the Zeeman energy. The gauge field describing this is denoted by $`\beta _\mu `$ in Eq. (2.1). The spin Hall conductance is then related to the ‘spin-current’ induced perpendicular to the direction of the gradient of the Zeeman energy.
Let us now briefly recall the results obtained from this formulation for the filling factors and the shift corresponding to a surface of genus $`g`$. The filling factors can be calculated by means of simple inner products<sup>1</sup><sup>1</sup>1Throughout this paper the transpose in equations like (2) is implicitly understood in order to simplify the notation.
$`\nu =𝐭_e𝐊_e^1𝐭_e=𝐭_\varphi 𝐊_\varphi ^1𝐭_\varphi ,`$
$`\sigma =𝐬_e𝐊_e^1𝐬_e=𝐬_\varphi 𝐊_\varphi ^1𝐬_\varphi .`$ (2.2)
The relation between the charge (and spin) vectors of the electron and quasihole parts are given by
$$𝐭_\varphi =𝐊_e^1𝐭_e,𝐬_\varphi =𝐊_e^1𝐬_e.$$
(2.3)
The last important property we will discuss is the so called ‘shift’ in the flux on surfaces of general genus $`g`$. The relation between the number of electrons $`N_e`$ and the corresponding number of flux quanta $`N_\mathrm{\Phi }`$ is given by
$$N_\mathrm{\Phi }=\frac{1}{\nu }N_e𝒮,$$
(2.4)
where the shift $`𝒮`$ is given by
$$𝒮=\frac{2(1g)}{\nu }(𝐭_e𝐊_e^1𝐣_e).$$
(2.5)
Although $`𝐣_e`$ plays a somewhat different role than $`𝐭_e`$ and $`𝐬_e`$, we define $`𝐣_\varphi `$ by analogy to (2.3)
$$𝐣_\varphi =𝐊_e^1𝐣_e.$$
(2.6)
In the present paper, we shall establish that the various relations given above are not just valid for the abelian case. They also apply in the non-abelian case, under the condition that a formulation is used in which the pseudoparticles do not carry charge or spin (see Sect. 4). We shall see that in all the cases we consider, such a formulation can indeed be given.
The other important role the $`K`$-matrices play will be described in the next section, namely the role as statistics matrix in the sense of the Haldane exclusion statistics of the (quasi) particles. Also, we will explain a notion of ‘duality’ which is important in this context, and rederive some of the relations given above.
## 3. Abelian exclusion statistics
An important consequence of the concept of an ‘ideal gas of fractional statistics particles’ is the notion of 1-particle distribution functions which generalize the familiar Fermi-Dirac and Bose-Einstein distributions. These distributions can be derived from ‘1-particle grand canonical partition functions’. These quantities, which we denote by $`\lambda _i`$, satisfy the following set of equations, which were independently derived by Isakov, Dasnières de Veigy-Ouvry and Wu (IOW)
$$\left(\frac{\lambda _i1}{\lambda _i}\right)\underset{j}{}\lambda _j^{𝐊_{ij}}=z_i,$$
(3.1)
where $`\lambda _i=\lambda _i(z_1,\mathrm{},z_n)`$, with $`z_i=e^{\beta (\mu _iϵ)}`$ the generalized fugacity of species $`i`$. Note that the energy $`ϵ`$ may also include contributions from the coupling of the charge and spin of the quasiparticles to external electric and magnetic fields. Hence the information about charge and spin of the quasiparticles is also encoded in these generalized fugacities. The fugacities of the particles will be important for the distinction between abelian and non-abelian statistics, as we will point out later. The matrix $`𝐊`$ is the so-called ‘statistics matrix’ and describes, at least in the original situation in which Haldane introduced his new notion of statistics, the statistical interaction of particles of different species.
From the solutions $`\lambda _i`$ of the IOW-equations (3.1) the one-particle distribution functions $`n_i(ϵ)`$ are obtained as
$$n_i(ϵ)=z_i\frac{}{z_i}\mathrm{log}\underset{j}{}\lambda _j{}_{|z_i=e^{\beta (\mu _iϵ)}}{}^{}=\underset{j}{}z_j\frac{}{z_j}\mathrm{log}\lambda _i{}_{|z_i=e^{\beta (\mu _iϵ)}}{}^{},$$
(3.2)
where we have assumed that the matrix $`𝐊`$ is symmetric.
The relation between, on the one hand, the $`K`$-matrix of an abelian quantum Hall fluid and, on the other hand, the exclusion statistics of its charged edge excitations, can be described as follows. The charged edge excitations are described by a specific Conformal Field Theory (CFT), also known as a chiral Luttinger liquid. Following a procedure first proposed in , one may associate a notion of fractional exclusion statistics to a set of fundamental excitations in this CFT. Selecting a particular set of negatively charged ‘electron type’ excitations together with a ‘dual’ set of positively charged quasihole excitations, one precisely finds fractional exclusion statistics in the sense of Haldane, with statistics matrix $`𝐊`$ given by
$$𝐊=𝐊_e𝐊_\varphi ,$$
(3.3)
with $`𝐊_e`$ and $`𝐊_\varphi `$ the $`K`$-matrices for the abelian quantum Hall state. For the principal Laughlin series at filling fraction $`\nu =1/m`$, this result was obtained in , in its general form it first appeared in our paper . The relation of the identification (3.3) with character identities involving so called Universal Chiral Partition Functions will be discussed in Sect. 5.
In , a slightly different identification between the $`K`$-matrix and a statistics matrix, amounting to $`𝐊=𝐊_e`$, was proposed. The two proposals can be reconciled by realizing that we, in our analysis of edge excitations, restrict ourselves to quanta of positive energy only. From the duality relations that we discuss below, one learns that, in a precise sense, quasihole quanta of positive energy can be traded for holes in a ‘Fermi sea’ of electron-type quanta at negative energy, and in this way one arrives at a complete description in terms of the matrix $`𝐊_e`$ alone.
One of the main themes in this paper will be the identification of statistics matrices $`𝐊`$ for excitations over non-abelian quantum Hall states. Extending the identification (3.3) to the non-abelian case, we shall propose $`K`$-matrices for the non-abelian quantum Hall states. We would like to stress that, although many of the formulas from the well known abelian $`K`$-matrix description still hold for the generalized $`K`$-matrices we find here, the description for the non-abelian states is on an entirely different footing. The abelian $`K`$-matrices were introduced to describe quantum Hall states in the ‘most general’ way, i.e. by trying to implement the hierarchical schemes in a general way. In the non-abelian case, we need the $`K`$-matrix structure to keep track of the non-abelian statistics. So although we use a matrix structure, we are not describing a hierarchical situation.
We continue this section with a discussion of the fundamental ‘particle-hole’ duality between the electron and the quasihole sectors of the theory. To show how this duality works, we assume that we have $`n`$ quasiholes $`\varphi `$ and $`n`$ electron-like particles $`\mathrm{\Psi }`$ described by the matrices $`𝐊_\varphi `$ and $`𝐊_e`$, respectively. We assume that (i) $`𝐊_\varphi =𝐊_e^1`$, and (ii) there is no mutual exclusion statistics between the two sectors (meaning that the statistics matrix is given by the direct sum (3.3)). These two conditions in fact constitute what we mean by duality in this context. In the context of low-energy effective actions for abelian fqH systems, a similar notion of duality has been considered (see, e.g. and references therein).
With the matrices $`𝐊_\varphi `$ and $`𝐊_e`$, two independent systems of IOW-equations can be written down, and these systems are related by the duality (for clarity, we will denote the single level partition function for the quasiholes and electron-like particles by $`\lambda _i`$ and $`\mu _i`$ respectively; the corresponding fugacities will be denoted by $`x_i`$ and $`y_i`$)
$$\lambda _i=\frac{\mu _i}{\mu _i1},x_i=\underset{j}{}y_j^{(𝐊_e)_{ij}^1},$$
(3.4)
as can be verified easily.
As an illustration of the duality, we calculate the central charge of the conformal field theory that describes the edge excitations. We focus on the abelian case. In the non-abelian case, which we discuss in the next section, there will be a subtraction term due to the presence of pseudoparticles.
In general, for abelian quantum Hall states, the central charge $`c_{\text{CFT}}`$ is given by
$$c_{\text{CFT}}=\frac{6}{\pi ^2}_0^1\frac{dz}{z}\mathrm{log}\lambda _{\text{tot}}(z),$$
(3.5)
where $`\lambda _{\text{tot}}(z)`$ denotes the product $`_j\lambda _j`$ evaluated at $`z_j=z`$ for all $`j`$. It has been shown (see and references therein), that this can be rewritten in the following form
$$c_{\text{CFT}}=\frac{6}{\pi ^2}\underset{i}{}L(\xi _i),$$
(3.6)
where $`L(z)`$ is Rogers’ dilogarithm
$$L(z)=\frac{1}{2}_0^z𝑑y\left(\frac{\mathrm{log}y}{1y}+\frac{\mathrm{log}(1y)}{y}\right).$$
(3.7)
The quantities $`\xi _i`$ which appear in Eq. (3.6) are solutions to the central charge equations
$$\xi _i=\underset{j}{}(1\xi _j)^{𝐊_{ij}}.$$
(3.8)
For the abelian quantum Hall case, we have two matrices $`𝐊_\varphi `$ and $`𝐊_e`$ and we need the solutions $`\xi _i`$ and $`\eta _i`$ of the equations
$$\xi _i=\underset{j=1}{\overset{n}{}}(1\xi _j)^{(𝐊_\varphi )_{ij}},\eta _i=\underset{j=1}{\overset{n}{}}(1\eta _j)^{(𝐊_e)_{ij}}.$$
(3.9)
By virtue of the duality, these solutions are related by a simple equation: $`\eta _i=1\xi _i`$. This leads to
$$\underset{i}{}L(\xi _i)+\underset{i}{}L(\eta _i)=\underset{i}{}\left(L(\xi _i)+L(1\xi _i)\right)=nL(1)=n\frac{\pi ^2}{6}.$$
(3.10)
So in the abelian case, we correctly find that the central charge is just given by the number of species in the theory, $`c_{\text{CFT}}=n`$.
## 4. Non-abelian exclusion statistics
In this section, we focus on $`K`$-matrices and statistics matrices for non-abelian quantum Hall states. We shall first introduce new types of particles, pseudoparticles and composite particles, and explain the role they play in the non-abelian case. We also extend the notion of duality to the non-abelian case. After that we discuss various aspects (filling factors and shift map) of the quantum Hall data $`𝐊`$, $`𝐭`$, $`𝐬`$ and $`𝐣`$ in the non-abelian case.
Among the new particles that appear in non-abelian theories are so called ‘composite’ particles in the electron sector. These will show up as particles which have multiple electron charges. We introduce an integer label $`l_i`$ for an order-$`l_i`$ composite particle of charge $`(𝐭_e)_i=l_i`$.
In the quasihole sector, we encounter so called pseudoparticles, which do not carry any energy, but rather act as a book-keeping device that keep track of ‘internal degrees of freedom’ of the physical quasiholes. Pseudoparticles were first introduced in the TBA analysis of integrable systems with non-diagonal particle scattering (see, e.g. ); in the context of exclusion statistics they have been discussed in . We assign the label $`l_i=0`$ to all pseudoparticles.
An important observation, first made in , is that the duality between the electron and quasihole sectors naturally links the presence of composite particles in one sector to the presence of pseudoparticles in the other. Physically, this is a link between the pairing physics of the non-abelian quantum Hall states and the non-abelian statistics of their fundamental excitations.
### 4.1. Composites, pseudoparticles and null-particles
The presence of pseudoparticles and composite particles calls for a slight generalization of the discussion of the previous section. When focusing on the dependence of the $`\lambda _i`$ on the energy $`ϵ`$, the natural specialization of the generalized fugacities $`z_i`$ is given by $`z_i=z^{l_i}`$, with $`z=e^{\beta ϵ}`$. In the presence of $`l_i1`$, the 1-particle distribution functions take the form \[note that a composite particle labeled by $`ϵ`$ carries energy $`l_iϵ`$\]
$$n_i(ϵ)=z_i\frac{}{z_i}\mathrm{log}\underset{j}{}[\lambda _j]^{l_j}{}_{|z_i=e^{\beta (\mu _il_iϵ)}}{}^{}=\underset{j}{}l_jz_j\frac{}{z_j}\mathrm{log}\lambda _i{}_{|z_i=e^{\beta (\mu _il_iϵ)}}{}^{}.$$
(4.1)
With the following definition of $`\lambda _{\text{tot}}(z)`$
$$\lambda _{\text{tot}}(z)=\underset{i}{}[\lambda _i(z_j=z^{l_j})]^{l_i},$$
(4.2)
the central charge $`c_{\text{CFT}}`$ is again given by the expression (3.5). We note that in the specialized IOW equations, with $`z_i=z^{l_i}`$, the right hand side of the equations for pseudoparticles is equal to 1. When focusing on quantum numbers other than energy, such as spin, we will consider slightly more general versions of the quantity $`\lambda _{\text{tot}}`$.
In all examples (abelian and non-abelian) that are explicitly discussed in this paper, we assume a choice of particle basis such that $`𝐥_e=𝐭_e`$. For the abelian quantum Hall states we further assume that $`(𝐭_e)_i=1`$ for all $`i`$. In the quasihole sector we specify $`(𝐥_\varphi )_i=\frac{1}{q_{\text{qp}}}(𝐊_\varphi )_{ij}(𝐥_e)_j`$, where $`q_{\text{qp}}`$ is the smallest (elementary) charge in the quasihole sector. \[This implies that, even in the abelian case, we may treat some of the quasiholes as composites of the most fundamental ones, thereby generalizing the discussion of the previous section.\]
Under these assumptions, we find that under duality $`\lambda _{\text{tot}}(x)`$ and $`\mu _{\text{tot}}(y)`$ are related in the following way
$$\lambda _{\text{tot}}(x)=x^\gamma \mu _{\text{tot}}^\alpha (y),y=x^\beta ,$$
(4.3)
with
$$\alpha =\beta =\frac{1}{q_{\text{qp}}},\gamma =\frac{\nu }{q_{\text{qp}}^2}.$$
(4.4)
A clear sign of non-abelian statistics is found in the way the quantity $`\lambda _i`$ for physical particles depends on the fugacity $`z_i`$. Putting $`z_l=1`$ for all pseudoparticles, and focusing on the small $`z`$ behaviour of $`\lambda _i`$, one finds
$$\lambda _i=1+\alpha _iz_i+o(z^2).$$
(4.5)
In the abelian case, $`\alpha _i=1`$, whereas in the non-abelian case $`\alpha _i>1`$. The factors $`\alpha _i`$ lead to multiplicative factors in the Boltzmann tails of the one-particle distribution functions for physical particles. The quantities $`\alpha _i`$ are in fact the largest eigenvalues of the fusion matrix , i.e., the quantum dimensions of the conformal field theory associated to the quantum Hall state, and can easily be calculated for the cases we deal with (see Sects. 6 and 7.2).
In , we presented a generalized $`K`$-matrix structure for some recently proposed quantum Hall states. The proposed $`K`$-matrices were identified via their role as statistics matrices for the fundamental charged edge excitations. In the quasihole sector, the non-abelian statistics leads to a specific set of pseudoparticles and an associated statistics matrix $`𝐊_\varphi `$ . The matrix $`𝐊_e`$, related to $`𝐊_\varphi `$ by the duality $`𝐊_e=𝐊_\varphi ^1`$, refers to particles which are identified as composites of the fundamental electron-like excitation. From the point of view of the wave functions for the non-abelian quantum Hall states , the presence of composite excitations is very natural. This is because the non-abelian states show a behaviour which is called clustering (of order $`k`$, where $`k`$ is a label of the states ). This order-$`k`$ clustering means that up to $`k`$ particles can come to the same position, without making the wave function zero, whereas, as soon as $`k+1`$ particles are located at the same positions, the wave function becomes identically zero. In it was argued that the wave functions which show pairing (at $`k=2`$), are related (in the non-magnetic limit, i.e. in the limit of $`\nu \mathrm{}`$) to BCS superconductivity.
Composite particles are identified as particles whose generalized fugacities are specific combinations of the generalized fugacities of other particles, i.e., all quantum numbers of composite particles are completely determined in terms of the quantum numbers of their constituents. It has been shown in that particular kinds of composite particles, so-called null-particles, accounting for the null-states in the quasiparticle Fock spaces, are often needed to interpret the system in terms of Haldane’s exclusion statistics or, equivalently, to write the partition function in UCPF form (see also Sect. 5.2).
We now turn to the computation of the central charge $`c_{\text{CFT}}`$ the non-abelian case. It was shown in , that the presence of pseudoparticles leads to a simple correction term that is subtracted from the abelian result $`c_{\text{CFT}}=n`$. For the pseudoparticles, a system of equations like Eq. (3.9) can be written down
$$\xi _i^{}=\underset{j}{^{}}(1\xi _j^{})^{𝐊_{ij}},$$
(4.6)
where the prime indicates that the product is restricted to pseudoparticles. The correction term is given by a sum over the dilogarithm of the solutions of (4.6), leading to
$$c_{\text{CFT}}=n\frac{6}{\pi ^2}\underset{i}{^{}}L(\xi _i^{}).$$
(4.7)
### 4.2. On filling factors
Up to now, we merely asserted that the statistics matrices $`𝐊`$ can also serve as (generalized) $`K`$-matrices for non-abelian quantum Hall states. To make this statement more clear, we will now investigate how some of the ‘$`K`$-matrix results’ for abelian quantum Hall states generalize to the non-abelian case. In this derivation, we make the assumption that the pseudoparticles do not carry charge or spin. In all cases that are explicitly considered in Sects. 6 and 7 this assumption holds in the simplest formulation. If pseudoparticles do carry spin or charge, the formulas we obtain below need to be modified.
Let us start with the filling factor corresponding to state which is described by the IOW-equations, for a statistics matrix $`𝐊_e`$, charge vector $`𝐭_e`$, and labels $`𝐥_e=𝐭_e`$. We couple the system to an electric field by taking $`y_i=y^{(𝐭_e)_i}`$. \[This is when the orientation of the electric field is such that the response is carried by the negatively charged excitations.\] The large $`y`$ (i.e. low temperature) behaviour of the IOW-equations (3.1) is then given by the following set of relations
$$\underset{j}{}\mu _j^{(𝐊_e)_{ij}}y^{(𝐭_e)_i},$$
(4.8)
which imply, when $`𝐊`$ is symmetric (which is assumed throughout the paper) and invertible
$$\mu _{\text{tot}}=\underset{i}{}\mu _i^{(𝐭_e)_i}y^{𝐭_e𝐊_e^1𝐭_e}.$$
(4.9)
Because the left hand side of Eq. (4.9) in the $`T0`$ limit determines the filling factor $`\nu `$ through $`\mu _{\text{tot}}y^\nu `$, we find the well-known formula
$$\nu =𝐭_e𝐊_e^1𝐭_e.$$
(4.10)
For the opposite orientation of the electric field, a similar expression is obtained by starting from the $`K`$-matrix for the (positively charged) quasiholes
$$\nu =𝐭_\varphi 𝐊_\varphi ^1𝐭_\varphi .$$
(4.11)
This result could also have been obtained by using Eq. (4.10) and the transformation properties of $`𝐊_e`$ and $`𝐭_e`$ under duality. We remark that the above derivations explicitly assume that only the physical particles respond to the electric field, i.e., that all pseudoparticles are neutral.
Let us now turn to the spin Hall conductance, and the corresponding spin filling factor. The derivation of the corresponding spin filling factor
$$\sigma =𝐬_e𝐊_e^1𝐬_e,$$
(4.12)
goes along the same lines as the derivation of the electron filling factor. As an extra step, one needs to relate the fugacities of the spin up and down particles by $`y_{}=1/y_{}=z`$. This results in
$$\underset{i}{}\mu _i^{(𝐬_e)_i}z^{𝐬_e𝐊_e^1𝐬_e},$$
(4.13)
leading to Eq. (4.12). It is important to note that this formula only holds in the cases where the pseudoparticles in the $`\varphi `$-sector do not carry spin. As a check on this formula, one would like to have a procedure to obtain the spin filling factor directly from the wave functions, as is possible for the electron filling factor. To do this, one has to count the zeros of the wave function with respect to one reference particle (of a given spin, say, up). The total number of zeros gives the total flux needed on the sphere as a linear function of the total number of electrons $`N_e`$. By using the relation between $`N_e`$ and $`N_\mathrm{\Phi }`$ given in (2.4) one obtains the electron filling factor and the shift. To obtain the spin filling factor, one has to keep track of two different types of zeros, namely those with respect to a particle of the same spin, and the ones with respect to particles of the other spin. We will denote the number of these zeros by $`N_\mathrm{\Phi }^{}`$ and $`N_\mathrm{\Phi }^{}`$ respectively. The electron and spin filling factors are obtained from
$`N_\mathrm{\Phi }=N_\mathrm{\Phi }^{}+N_\mathrm{\Phi }^{}`$ $`=`$ $`{\displaystyle \frac{1}{\nu }}N_e𝒮,`$
$`N_\mathrm{\Phi }^{}N_\mathrm{\Phi }^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sigma }}N_e𝒮.`$ (4.14)
We applied this procedure to the non-abelian spin singlet states of (the explicit form of the wave functions will be given elsewhere ), and indeed found the same results for the electron and spin filling factor as obtained from the $`K`$-matrix formalism, Eq. (7.1). Also the electron filling factor for the Read-Rezayi states is reproduced correctly, see Eq. (6.1). In addition, for both types of states we found that the shift on the sphere is in agreement with (2.5) for $`g=0`$.
Summarizing, we have presented evidence that duality relations
$$𝐊_\varphi =𝐊_e^1,𝐭_\varphi =𝐊_e^1𝐭_e,𝐬_\varphi =𝐊_e^1𝐬_e,𝐣_\varphi =𝐊_e^1𝐣_e.$$
(4.15)
are applicable to both abelian and non-abelian quantum Hall states, and that the expressions (2) for the filling factors $`\nu `$ and $`\sigma `$ apply to the non-abelian case, in a formulation where pseudoparticles do not carry spin or charge.
### 4.3. Shift map
Suppose we have a fractional quantum Hall system described by the data $`(𝐊_e,𝐭_e,𝐬_e,𝐣_e)`$. We can then construct a family of fractional quantum Hall systems, parametrized by $`M\text{}_+`$, by applying the ‘shift map’ $`𝒮_M`$ introduced in . In the cases we consider, $`M`$ odd (even) corresponds to a fermionic (bosonic) state respectively. At the level of trial wave functions $`\mathrm{\Psi }(z)`$, $`𝒮_M`$ simply acts as a multiplicative Laughlin factor $`_{i<j}(z_iz_j)^M`$. Thus, $`𝒮_M`$ increases the number of flux quanta by
$$N_\mathrm{\Phi }N_\mathrm{\Phi }+M(N_e1)=(\frac{1}{\nu }+M)N_e(𝒮+M),$$
(4.16)
i.e.,
$$\nu ^1\nu ^1+M,\sigma \sigma ,𝒮𝒮+M.$$
(4.17)
In fact, $`𝒮_M`$ acts on the fqH data $`(𝐊_e,𝐭_e,𝐬_e,𝐣_e)`$ as
$`𝒮_M𝐊_e`$ $`=`$ $`𝐊_e+M𝐭_e𝐭_e,`$
$`𝒮_M𝐭_e`$ $`=`$ $`𝐭_e,`$
$`𝒮_M𝐬_e`$ $`=`$ $`𝐬_e,`$
$`𝒮_M𝐣_e`$ $`=`$ $`𝐣_e+\frac{M}{2}𝐭_e.`$ (4.18)
One easily checks that (4.3), together with (4.10), leads to the shift in $`\nu ^1`$ as given in (4.17). By duality (4.15) one obtains
$`𝒮_M𝐊_\varphi `$ $`=`$ $`𝐊_\varphi \frac{M}{1+\nu M}𝐭_\varphi 𝐭_\varphi ,`$
$`𝒮_M𝐭_\varphi `$ $`=`$ $`\frac{1}{1+\nu M}𝐭_\varphi ,`$
$`𝒮_M𝐬_\varphi `$ $`=`$ $`𝐬_\varphi ,`$
$`𝒮_M𝐣_\varphi `$ $`=`$ $`𝐣_\varphi \frac{M}{2}\left(\frac{\nu 𝒮1}{1+\nu M}\right)𝐭_\varphi .`$ (4.19)
A few remarks should be made. By using the duality (4.15), one actually finds for the action of the shift map on $`𝐬_\varphi `$: $`𝒮_M𝐬_\varphi =𝐬_\varphi +\frac{M(𝐭_\varphi 𝐬_e)}{1+\nu M}𝐭_\varphi `$. However, the shift map is only supposed to act on the charge component of the particles, thus we would like to demand that $`𝒮_M𝐬_\varphi =𝐬_\varphi `$. Therefore, for consistency, we require
$$𝐭_\varphi 𝐬_e=𝐭_e𝐊_e^1𝐬_e=0,$$
(4.20)
leading to (4.3). Of course, relation (4.20) is just the statement that for spin singlet states there should be a $`\text{}_2`$ symmetry $`(𝐭_e,𝐬_e)(𝐭_e,𝐬_e)`$. Eq. (4.20) is fulfilled for all our examples (if we take $`𝐬_e=0`$ for the spin polarized states). Although, in general, $`𝐣_e`$ has to be treated as an independent variable, for the examples discussed in Sects. 6 and 7 all formulas are consistent with the relation $`𝐣_e=𝐬_e+(𝒮/2(1g))𝐭_e`$.
In this paper we will be mainly concerned with fractional quantum Hall systems corresponding to conformal field theories $`\widehat{𝔤}_{k,M}`$ which are deformations of the conformal field theory based on the affine Lie algebra $`\widehat{𝔤}_k`$ at level $`k`$. The $`\widehat{𝔤}`$-symmetry greatly simplifies the determination of the fqH data $`(𝐊_e,𝐭_e,𝐬_e,𝐣_e)`$ for $`\widehat{𝔤}_k`$. The fqH data for $`(\widehat{𝔤})_{k,M}`$ are then simply obtained by applying the shift operator $`𝒮_M`$ as in (4.3). The action of the shift map can be visualized as follows. Charge is usually identified with a particular direction in the weight lattice of $`𝔤`$. The degrees of freedom associated to this direction can be represented by a chiral boson compactified on a circle of some radius $`R`$. The shift map $`𝒮_M`$ has the effect of rescaling the radius $`R`$ while keeping all other directions in the weight diagram fixed.
### 4.4. Composites
The description of a physical system in terms of a set of $`n`$ quasiparticles with mutual exclusion statistics given by a matrix $`(𝐊_{ij})_{1i,jn}`$ is not unique. In particular one may extend the number of quasiparticles by introducing composites as we will now explain.
Consider the IOW-equations (3.1) with
$$𝐊=\left(\begin{array}{ccc}a_{11}& \mathrm{}& a_{1n}\\ \mathrm{}& & \mathrm{}\\ a_{n1}& \mathrm{}& a_{nn}\end{array}\right),𝐳=\left(\begin{array}{c}z_1\\ \mathrm{}\\ z_n\end{array}\right).$$
(4.21)
If we define the operation $`𝒞_{ij}`$, corresponding to adding a composite of the quasiparticles $`i`$ and $`j`$ to the system, by
$$𝒞_{ij}𝐊=\left(\begin{array}{cccccc}a_{11}& \mathrm{}& & a_{1n}& \mathrm{}& a_{1i}+a_{1j}\\ \mathrm{}& & & \mathrm{}& \mathrm{}& \mathrm{}\\ & & a_{ij}+1& & \mathrm{}& \\ & & & & \mathrm{}& \\ & a_{ji}+1& & & \mathrm{}& \\ & & & & \mathrm{}& \\ a_{n1}& \mathrm{}& & a_{nn}& \mathrm{}& a_{ni}+a_{nj}\\ [1.5]4& \mathrm{}& [1.5]1\\ a_{i1}+a_{j1}& \mathrm{}& & a_{in}+a_{jn}& \mathrm{}& a_{ii}+2a_{ij}+a_{jj}\end{array}\right),$$
(4.22)
and
$$𝒞_{ij}𝐳=(𝐳_1,\mathrm{},𝐳_n;𝐳_i𝐳_j),$$
(4.23)
such that, in particular,
$`𝒞_{ij}𝐭`$ $`=`$ $`(𝐭_1,\mathrm{},𝐭_n;𝐭_i+𝐭_j),`$
$`𝒞_{ij}𝐬`$ $`=`$ $`(𝐬_1,\mathrm{},𝐬_n;𝐬_i+𝐬_j),`$ (4.24)
then the two systems are equivalent, at least at the level of thermodynamics. The solutions $`\{\lambda _i\}`$ to the IOW-equations defined by $`(𝐊,𝐳)`$ and $`\{\lambda _i^{}\}`$ defined by $`(𝐊^{},𝐳^{})=(𝒞_{ij}𝐊,𝒞_{ij}𝐳)`$ are simply related by
$`\lambda _i^{}`$ $`={\displaystyle \frac{\lambda _i+\lambda _j1}{\lambda _j}},`$ $`\lambda _j^{}`$ $`={\displaystyle \frac{\lambda _i+\lambda _j1}{\lambda _i}},`$
$`\lambda _{n+1}^{}`$ $`={\displaystyle \frac{\lambda _i\lambda _j}{\lambda _i+\lambda _j1}},`$ $`\lambda _k^{}`$ $`=\lambda _k,(ki,j,n+1).`$ (4.25)
Note that, in particular, it follows $`\lambda _i=\lambda _i^{}\lambda _{n+1}^{}`$ and $`\lambda _j=\lambda _j^{}\lambda _{n+1}^{}`$ such that $`\lambda _{\text{tot}}=\lambda _{\text{tot}}^{}`$. Also, from $`\lambda _i=\lambda _i^{}\lambda _{n+1}^{}`$ and $`\lambda _j=\lambda _j^{}\lambda _{n+1}^{}`$ one sees that the original one-particle partition functions for $`i`$ and $`j`$, receive contributions from the new particles $`i`$ and $`j`$, respectively, as well as from the composite particle $`n+1`$. The operation $`𝒞_{ij}`$ has the effect that states in the spectrum containing both particles $`i`$ and $`j`$ get less dense (their mutual exclusion statistics is bumped up by $`1`$), while the resulting ‘gaps’ are now filled by the new composite particle.
A consistency check on the equivalence of the systems described by $`(𝐊,𝐳)`$ and $`(𝐊^{},𝐳^{})`$ is the fact that both lead to the same central charge as a consequence of the five-term identity for Rogers’ dilogarithm (see ).
Finally, note that the shift map $`𝒮_M`$ of Eq. (4.3) and composite operation $`𝒞_{ij}`$ of Eqs. (4.22) and (4.4) commute, i.e.
$$𝒮_M𝒞_{ij}=𝒞_{ij}𝒮_M,$$
(4.26)
as one would expect.
## 5. The UCPF and exclusion statistics
### 5.1. Quasiparticle basis and truncated partition function
Quasiparticles in two dimensional conformal field theories are represented by so-called chiral vertex operators $`\varphi ^{(i)}(z)`$ that intertwine between the irreducible representations of the chiral algebra. Given a set of quasiparticles $`\varphi ^{(i)}(z)`$, $`i=1,\mathrm{},n`$, one has to determine a basis for the Fock space created by the modes $`\varphi _s^{(i)}`$, i.e., a maximal, linearly independent set of vectors
$$\varphi _{s_N}^{(i_N)}\mathrm{}\varphi _{s_2}^{(i_2)}\varphi _{s_1}^{(i_1)}|\omega ,$$
(5.1)
with suitable restrictions on the mode sequences $`(s_1,\mathrm{},s_N)`$ (which may depend on the ‘fusion paths’ $`(i_1,\mathrm{},i_N)`$), as well as a set of vacua $`|\omega `$ (see for more details). The partition function $`Z(𝐳;q)`$ is then defined by
$$Z(𝐳;q)=\mathrm{Tr}\left((\underset{i}{}z_i^{N_i})q^{L_0}\right),$$
(5.2)
where the trace is taken over the basis (5.1) and $`N_i`$ denotes the number operator for quasiparticles of type $`i`$ while $`L_0=_is_i`$ for a state of type (5.1). During this discussion on the UCPF, we use the following, in the literature standard notation $`q=e^{\beta ϵ_0}`$, where $`ϵ_0`$ is some fixed energy scale, and $`z_i=e^{\beta \mu _i}`$.
Exclusion statistics in conformal field theory can be studied by means of recursion relations for truncated partition functions . Truncated partition functions $`P_𝐋(𝐳;q)`$, for $`𝐋=(L_1,\mathrm{},L_n)`$, are defined by taking the partition function of those states (5.1) where all the modes $`s`$ for quasiparticles of species $`i`$ satisfy $`sL_i`$. By definition, for large $`𝐋`$, we will have (see for more details)
$$P_{𝐋+𝐞_i}(z;q)/P_𝐋(z;q)\lambda _i(z_iq^{L_i}),$$
(5.3)
where $`𝐞_i`$ denotes the unit vector in the $`i`$-direction. In particular, if the generalized fugacities $`z_i`$ are given by $`z_i=z^{l_i}`$, for some fixed $`z`$, and the quasiparticle modes are truncated by $`L_i=l_iL`$, then we find, using (4.2)
$$P_{L+1}(z;q)/P_L(z;q)\lambda _{\text{tot}}(zq^L),$$
(5.4)
where $`P_L(z;q)=P_{l_1L,l_2L,\mathrm{},l_nL}(z_i=z^{l_i};q)`$. Thus, given a set of recursion relations for the truncated partition functions $`P_𝐋(z;q)`$, one derives algebraic equations for the one-particle partition functions $`\lambda _i(z)`$ by taking the large $`𝐋`$ limit. In particular one can find an equation for $`\lambda _{\text{tot}}(z)`$ from $`P_L(z;q)`$ by using (5.4). For all conformal field theories that have been studied this way it turns out that one finds agreement between these $`\lambda `$-equations and the IOW-equations (3.1) corresponding to a specific statistics matrix $`𝐊`$ (see, in particular, ).
### 5.2. The universal chiral partition function
Based on many examples, it has become clear that the characters of the representations of all conformal field theories can be written in the form of, what is now known as, a universal chiral partition function (UCPF) (see in particular, Ref. and references therein)
$$Z(𝐊;𝐐,𝐮|𝐳;q)=\underset{𝐦}{^{}}\left(\underset{i}{}z_i^{m_i}\right)q^{\frac{1}{2}𝐦𝐊𝐦+𝐐𝐦}\underset{i}{}\left[\genfrac{}{}{0.0pt}{}{((\mathbb{𝟙}𝐊)𝐦+𝐮)_i}{m_i}\right],$$
(5.5)
where $`𝐊`$ is a (rational) $`n\times n`$ matrix, $`𝐐`$ and $`𝐮`$ are certain $`n`$-vectors and the sum over $`m_1,\mathrm{},m_n`$, is over the nonnegative integers subject to some restrictions (which, throughout this paper, are taken to be such that the coefficients in the $`q`$-binomials are integer). The $`q`$-binomial (Gaussian polynomial) is defined by
$$\left[\genfrac{}{}{0.0pt}{}{M}{m}\right]=\frac{(q)_M}{(q)_m(q)_{Mm}},(q)_m=\underset{k=1}{\overset{m}{}}(1q^k).$$
(5.6)
The vectors $`𝐐`$ and $`𝐮`$ as well as the restrictions on the summation variables, will in general depend on the particular representation of the conformal field theory, while $`𝐊`$ is independent of the representation. To write the conformal characters in the form (5.5) may require introducing null-quasiparticles which account for null-states in the quasiparticle Fock space . The null-quasiparticles are certain composites, hence their fugacities $`z_i`$ in (5.5) are specific combinations of the fugacities of their constituents.
It has been conjectured that the UCPF (5.5) is precisely the partition function (5.2) of a set of quasiparticles with exclusion statistics given by the same matrix $`𝐊`$, where $`u_i=\mathrm{}`$ corresponds to a physical quasiparticle and $`u_i<\mathrm{}`$ to a pseudoparticle . This conjecture has been verified in numerous examples (see for references). A convincing piece of evidence in support of this conjecture is the fact that the asymptotics of the character (5.5) (in the thermodynamic limit $`q1^{}`$) is given by exactly the same formula as the one for the IOW-equations (see also for $`z_i=1`$). In the next section we establish the correspondence in a more direct way.
For future convenience let us introduce the limiting form of the UCPF (5.10) when all $`u_i\mathrm{}`$, i.e. the case that all quasiparticles are physical and the exclusion statistics is abelian
$$Z_{\mathrm{}}(𝐊;𝐐)=\underset{𝐦}{^{}}\left(\underset{i}{}z_i^{m_i}\right)\frac{q^{\frac{1}{2}𝐦𝐊𝐦+𝐐𝐦}}{_i(q)_{m_i}}.$$
(5.7)
Note that the limiting UCPFs (5.7) are not all independent, but satisfy (see )
$$Z_{\mathrm{}}(𝐊;𝐐)=Z_{\mathrm{}}(𝐊;𝐐+𝐞_i)+z_iq^{\frac{1}{2}𝐊_{ii}+𝐐_i}Z_{\mathrm{}}(𝐊;𝐐+𝐊𝐞_i),$$
(5.8)
as a consequence of
$$\frac{1}{(q)_m}=\frac{q^m}{(q)_m}+\frac{1}{(q)_{m1}}.$$
(5.9)
### 5.3. Relation to exclusion statistics
The relation between the UCPF and exclusion statistics can be made more explicit as follows. Suppose the truncated partition functions $`P_𝐋(𝐳;q)`$ are given by ‘finitized UCPFs’ of the form
$$P_𝐋(𝐳;q)=\underset{𝐦}{^{}}\left(\underset{i}{}z_i^{m_i}\right)q^{\frac{1}{2}𝐦𝐊𝐦+𝐐𝐦}\underset{i}{}\left[\genfrac{}{}{0.0pt}{}{(𝐋+(\mathbb{𝟙}𝐊)𝐦+𝐮)_i}{m_i}\right],$$
(5.10)
for some vectors $`(𝐐,𝐮)`$. Of course, the number of parameters in this expression is overdetermined. Usually we think of $`𝐮`$ as being fixed while the meaning of the parameters $`𝐋`$ are determined by the cut-off scale. We can of course absorb the $`𝐮`$ by shifts in $`𝐋`$ (in fact, in practice we often make shifts in the definition of $`𝐋`$ to simplify the recursion relations). We also remark that we have introduced finitization parameters $`L_i`$ also for the pseudoparticles in (5.10) to facilitate deriving recursion relations. In making the identification with the truncated partition functions these parameters are kept at a fixed (usually ‘small’ or even zero) value.
Using
$$\left[\genfrac{}{}{0.0pt}{}{M}{m}\right]=\left[\genfrac{}{}{0.0pt}{}{M1}{m}\right]+q^{Mm}\left[\genfrac{}{}{0.0pt}{}{M1}{m1}\right],$$
(5.11)
we find that $`P_𝐋(𝐳;q)`$ satisfies the system of recursion relations
$$P_𝐋(𝐳;q)=P_{𝐋𝐞_i}(𝐳;q)+z_iq^{\frac{1}{2}𝐊_{ii}+𝐐_i+𝐮_i+𝐋_i}P_{𝐋𝐊𝐞_i}(𝐳;q).$$
(5.12)
Upon dividing by $`P_𝐋(𝐳;q)`$, setting $`q=1`$, taking the large $`𝐋`$ limit, and using (5.3), we obtain
$$1=\lambda _i^1+z_i\underset{j}{}\lambda _j^{𝐊_{ji}},$$
(5.13)
which are equivalent to the IOW-equations (3.1) with statistics matrix $`𝐊`$.
Moreover, for any polynomial $`P_𝐋(𝐳;q)`$ satisfying the recursion relation (5.12), the polynomial
$$Q_𝐋(𝐳;q)=\left(\underset{i}{}z_i^{L_i}\right)q^{\frac{1}{2}𝐋𝐊𝐋+(𝐐+𝐮)𝐋}P_{𝐊𝐋}(𝐳;q^1),$$
(5.14)
satisfies the recursion relations (5.12) with dual data $`(𝐊^{};𝐐^{},𝐮^{},𝐳^{})`$, given by (cf. (3.4))
$$𝐊^{}=𝐊^1,𝐐^{}+𝐮^{}=𝐊^1(𝐐+𝐮),z_i^{}=\underset{j}{}z_j^{𝐊_{ij}^1}.$$
(5.15)
Thus, under the assumption that the set of finitized UCPFs (5.10), for fixed $`𝐐+𝐮`$, form a complete set of solutions to (5.12), the dual polynomial $`Q_𝐋(𝐳^{},q)`$ of (5.14) can again be written as a (finite) linear sum of finitized UCPFs with dual data (5.15). Moreover, by taking the large $`𝐋`$ limit of (5.14), using Eqs. (5.3) and (5.13), one recovers the duality relations (3.4) and (4.3).
The above calculation shows that, for quasiparticles whose truncated partition function is given by an expression of the form (5.10), the thermodynamics of these quasiparticles is described by Haldane’s exclusion statistics with statistics matrix $`𝐊`$. Even though many truncated characters are indeed of the form (5.10) (we will encounter various examples in the remainder of this paper) this is not the general situation. However, in examples it turns out that for all recursion relations for truncated characters there is an associated recursion relation, leading to the same $`\lambda `$-equation, which does admit a solution of the form (5.10). The true solution to this recursion relation will in general differ from (5.10) by terms of order $`q^L`$. In a sense we can talk about the universality class of recursion relations as those recursion relations that give rise to the same $`\lambda `$-equations and hence the same exclusion statistics.
### 5.4. Composites, revisited
In Sect. 4.4 we have seen, at the level of thermodynamics (i.e. the IOW-equations), how to introduce composite particles into the system in such a way that the resulting system is equivalent to the original system. Due to the intimate relation of exclusion statistics with the UCPF, explained in Sect. 5.3, one would expect that a similar construction is possible at the level of the UCPF. Indeed, upon substituting the following polynomial $`q`$-identity (see App. A for a proof)
$$\begin{array}{c}\left[\genfrac{}{}{0.0pt}{}{M_1}{m_1}\right]\left[\genfrac{}{}{0.0pt}{}{M_2}{m_2}\right]=\hfill \\ \hfill \underset{m0}{}q^{(m_1m)(m_2m)}\left[\genfrac{}{}{0.0pt}{}{M_1m_2}{m_1m}\right]\left[\genfrac{}{}{0.0pt}{}{M_2m_1}{m_2m}\right]\left[\genfrac{}{}{0.0pt}{}{M_1+M_2(m_1+m_2)+m}{m}\right],\end{array}$$
(5.16)
into the UCPF (5.10) at the $`(i,j)`$-th entry, and subsequently shifting the summation variables $`m_im_i+m`$, $`m_jm_j+m`$, yields an equivalent UCPF, based on $`n+1`$ quasiparticles with data $`(𝒞_{ij}𝐊;𝒞_{ij}𝐐,𝒞_{ij}𝐮)`$ and $`𝒞_{ij}𝐳`$, where
$`𝒞_{ij}𝐐`$ $`=`$ $`(𝐐_1,\mathrm{},𝐐_n;𝐐_i+𝐐_j),`$
$`𝒞_{ij}𝐮`$ $`=`$ $`(𝐮_1,\mathrm{},𝐮_n;𝐮_i+𝐮_j),`$ (5.17)
while $`𝒞_{ij}𝐊`$ and $`𝒞_{ij}𝐳`$ are defined in Eqs. (4.22) and (4.23), respectively. Various limiting forms of (5.16), relevant to introducing a composite of two physical particles or one physical particle and one pseudoparticle, are given in App. A as well.
## 6. $`𝔰𝔩_2`$: $`K`$-matrices for non-abelian spin polarized states
In this section we discuss a family of non-abelian spin polarized fractional quantum Hall systems with underlying conformal field theory $`(\widehat{𝔰𝔩_2})_{k,M}`$ and filling factor
$$\nu _{k,M}=\frac{k}{kM+2}.$$
(6.1)
For $`k=2`$ these systems, the so-called $`q`$-Pfaffians (where now $`q=1/\nu =M+1`$), were introduced in while the generalizations to $`k>2`$ were introduced in . The system contains a single quasihole $`\varphi `$, with charge $`1/(kM+2)`$ and an electron operator $`\mathrm{\Psi }`$ with charge $`1`$. At the $`(\widehat{𝔰𝔩_2})_k`$-point (i.e. $`M=0`$) the quasihole operator $`\varphi `$ has $`𝔰𝔩_2`$-weight $`\alpha /2`$, where $`\alpha `$ is the (positive) root of $`𝔰𝔩_2`$ and corresponds to one component of the chiral vertex operator transforming in the spin-$`1/2`$ representation (‘spinon’, see ), while the electron operator has weight $`\alpha `$ and corresponds to the current $`J_\alpha `$. For general $`M`$ the charge lattice has to be stretched.
The fqH data $`(𝐊_e,𝐭_e)`$ and their duals $`(𝐊_\varphi ,𝐭_\varphi )`$ for $`k=1`$ (corresponding to the abelian spin polarized Laughlin states with $`\nu =1/(M+2)`$ ) were discussed in and for $`k=2`$ (the $`q`$-Pfaffian) in . Here we discuss the generalization (see also ) to arbitrary $`k`$, corresponding to the Read-Rezayi states .
As indicated before, we analyze the conformal field theory $`(\widehat{𝔰𝔩_2})_{k,M}`$ by first analyzing the affine Lie algebra point $`M=0`$ and subsequently applying the shift map to obtain the result for general $`M`$.
The exclusion statistics and UCPF for the doublet of spinon operators in $`(\widehat{𝔰𝔩_2})_k`$ were studied in . It turns out that in this case we need $`k1`$ additional charge- and spin neutral pseudoparticles. Omitting the negative isospin spinon, we find (see, in particular, )
$$𝐊_\varphi =\left(\begin{array}{ccccccc}1& \frac{1}{2}& & & & \mathrm{}& \\ \frac{1}{2}& 1& \frac{1}{2}& & & \mathrm{}& \\ & \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}& \\ & & \frac{1}{2}& 1& \frac{1}{2}& \mathrm{}& \\ & & & \frac{1}{2}& 1& \mathrm{}& \frac{1}{2}\\ [1.5]7\\ & & & & \frac{1}{2}& \mathrm{}& \frac{1}{2}\end{array}\right),𝐭_\varphi =\left(\begin{array}{c}0\\ \mathrm{}\\ 0\\ \frac{1}{2}\end{array}\right),$$
(6.2)
leading, with (4.11), to a filling factor of $`\nu =k/2`$ in accordance with (6.1).
The data for arbitrary $`M`$ now follow by applying the shift map $`𝒮_M`$ of (4.3), i.e.
$$𝐊_\varphi ^M=𝒮_M𝐊_\varphi =\left(\begin{array}{ccccc}& & & \mathrm{}& \\ & \frac{1}{2}𝐀_{k1}& & \mathrm{}& \\ & & & \mathrm{}& \frac{1}{2}\\ [1.5]5\\ & & \frac{1}{2}& \mathrm{}& \frac{(k1)M+2}{2(kM+2)}\end{array}\right),𝐭_\varphi ^M=\left(\begin{array}{c}0\\ \mathrm{}\\ 0\\ \frac{1}{kM+2}\end{array}\right),$$
(6.3)
where, in order to simplify the notation, we have introduced the Cartan matrix $`𝐀_{k1}`$ of $`𝔰𝔩_k`$ (cf. (B.3)). One verifies that (4.11) is satisfied. The IOW-equations, determining the exclusion statistics of the quasiholes, can now be explicitly written down. E.g., for the $`q`$-Pfaffian ($`k=2`$) the following equation for $`\lambda _{\text{tot}}`$ easily follows from (3.1), in agreement with
$$(\lambda _{\text{tot}}1)(\lambda _{\text{tot}}^{\frac{1}{2}}1)=x^2\lambda _{\text{tot}}^{\frac{3M+2}{2(M+1)}}.$$
(6.4)
The small $`x`$ behaviour of $`\lambda _{\text{tot}}`$ for general $`k`$ was obtained from the IOW-equations in , with the result
$$\lambda _{\text{tot}}(x)=1+\alpha _kx+o(x^2),\alpha _k=2\mathrm{cos}\left(\frac{\pi }{k+2}\right).$$
(6.5)
It was argued that the factors $`\alpha `$ can also be obtained as quantum dimension of the appropriate CFT. It is easily checked that the small $`x`$ behaviour of $`\lambda _{\text{tot}}`$ in (6.4) indeed satisfies (6.5) for $`k=2`$. Similar equations for $`\lambda _{\text{tot}}`$ with $`k=3,4`$ were given in .
To determine the fqH data $`(𝐊_e,𝐭_e)`$ in the electron sector we observe that the electron operator $`\mathrm{\Psi }(z)`$ is identified with $`J_\alpha (z)`$. By acting with the negative modes of $`J_\alpha (z)`$ on the lowest weight vector in the lowest energy sector of some integrable highest weight module $`L(\mathrm{\Lambda })`$ at level $`k`$, one obtains what is known as the principal subspace $`W(\mathrm{\Lambda })`$ of $`L(\mathrm{\Lambda })`$ (or, rather, the reflected principal subspace). It is known that the character of the principal subspace can be written in the UCPF form (see App. B for a brief summary of the results for $`(\widehat{𝔰𝔩_n})_k`$). For $`(\widehat{𝔰𝔩_2})_k`$ this requires, besides the electron operator $`\mathrm{\Psi }`$ itself, clusters of up to $`k`$ electron operators. The corresponding $`K`$-matrix is given by the $`k\times k`$ matrix $`𝐊_e=2𝐁_k`$ where $`(𝐁_k)_{ij}=\mathrm{min}(i,j)`$ (see (B.4)), while $`𝐭_e=(1,2,\mathrm{},k)`$. Applying the shift map (4.3) thus gives
$$𝐊_e^M=\left(\begin{array}{cccc}M+2& 2M+2& \mathrm{}& kM+2\\ 2M+2& 2(2M+2)& \mathrm{}& 2(kM+2)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ kM+2& 2(kM+2)& \mathrm{}& k(kM+2)\end{array}\right),𝐭_e^M=\left(\begin{array}{c}1\\ 2\\ \mathrm{}\\ k\end{array}\right).$$
(6.6)
One easily verifies that the data $`(𝐊_\varphi ,𝐭_\varphi )`$ and $`(𝐊_e,𝐭_e)`$ are indeed related by the duality relations (4.15), and that Eqs. (4.10) and (4.11) are satisfied.
Moreover, the resulting IOW-equations for $`\mu _{\text{tot}}=\mu _1\mu _2^2`$ in case of the $`q`$-Pfaffian are given by
$$(\mu _{\text{tot}}^{2(M+1)}y^2)(\mu _{\text{tot}}^{M+1}y)=\mu _{\text{tot}}^{3M+2},$$
(6.7)
which are indeed related to (6.4) by the duality relations (4.3). Explicitly,
$$\lambda _{\text{tot}}(x)=y^2\mu _{\text{tot}}^{2(M+1)}(y),y=x^{2(M+1)}.$$
(6.8)
Finally, in order to show that the quasihole-electron system based on $`𝐊=𝐊_\varphi ^M𝐊_e^M`$ gives a complete description of the $`(\widehat{𝔰𝔩_2})_{k,M}`$ conformal field theory, we have to show that the chiral character of the latter can be written in terms of a (finite) combination of UCPF characters based on $`𝐊_\varphi ^M𝐊_e^M`$. This is indeed possible and discussed in App. C. Here we suffice to remark that the central charge, related to the asymptotic behaviour of the characters, works out correctly. Indeed, using standard dilogarithm identities one finds with (4.7)
$$c_\varphi +c_e=\frac{3k}{k+2},$$
(6.9)
which equals the central charge of $`(\widehat{𝔰𝔩_2})_{k,M}`$.
The above description of the Read-Rezayi states has an interesting application, namely the identification of a particle which acts as a supercurrent in the non-magnetic limit. This identification was made in , to which we refer for a more detailed discussion. We use the variable $`q=1/\nu =M+k/2`$, in terms of which the non-magnetic limit corresponds to $`q0`$. In this limit, all the statistics parameters of the largest composite (with charge $`k`$), go to zero, while the statistics parameters of the quasihole diverge. This is easily seen when one writes the statistic matrices (6.6) and (6.3) in terms of $`q`$. For these quantum Hall states the fundamental flux quantum is $`h/ke`$, because of the order-$`k`$ clustering. Upon piercing a quantum Hall state with this amount of flux, a quasihole with charge $`e/kq`$ is formed. This follows from the fact that the filling factor is $`e^2/qh`$ in physical units. For $`q1/k`$ this is the lowest charge possible and the electron like excitations correspond to multiple insertions of the flux quantum. This situation changes when we take the limit $`q0`$. Following , we take $`q=1/N`$, with $`N`$ a large integer. The largest composite is formed by inserting an amount of flux $`qkh/e=kh/Ne`$, thus a fraction of the flux quantum. The maximal occupation with this particle (in absence of other particles) is $`n_{\text{max}}=1/k^2q=N/k^2`$. Thus the maximal amount of flux that can be screened by this type of composites is $`(kh/Ne)(N/k^2)=h/ke`$, which is precisely the flux quantum. In conclusion we find that in the non-magnetic limit, the largest composite has bosonic statistics, and can screen an amount of flux up to the flux quantum. This clearly resembles the behaviour of the supercurrent in BCS superconductors.
## 7. $`𝔰𝔩_3`$: $`K`$-matrices for non-abelian spin singlet states
In a family of non-abelian spin singlet (NASS) states $`\mathrm{\Psi }_{k,M}`$ trial wave functions with filling factors
$$\nu _{k,M}=\frac{2k}{2kM+3},\sigma _{k,M}=2k,$$
(7.1)
was constructed. The system has two quasihole excitations $`\{\varphi _{},\varphi _{}\}`$ with one unit of up/down spin and charge $`1/(2kM+3)`$, while the electron operators $`\{\mathrm{\Psi }_{},\mathrm{\Psi }_{}\}`$ have charge $`1`$. The underlying conformal field theory is $`(\widehat{𝔰𝔩_3})_{k,M}`$. In terms of $`𝔰𝔩_3`$-weights the spin and charge assignment in the $`M=0`$ case is as follows. Denote the positive simple roots of $`𝔰𝔩_3`$ by $`\alpha _i`$, $`i=1,2`$ and the remaining positive non-simple root by $`\alpha _3=\alpha _1+\alpha _2`$. Let $`ϵ_i`$, $`i=1,2,3`$, denote the weights of the fundamental three dimensional irreducible representation $`\mathrm{𝟑}`$ of $`𝔰𝔩_3`$ such that $`ϵ_iϵ_j=\delta _{ij}1/3`$ and $`\alpha _i=ϵ_iϵ_{i+1}`$, $`i=1,2`$, then $`\{\varphi _{},\varphi _{}\}=\{\varphi ^{ϵ_1},\varphi ^{ϵ_2}\}`$ while $`\{\mathrm{\Psi }_{},\mathrm{\Psi }_{}\}=\{J_{\alpha _2},J_{\alpha _3}\}`$ (see Fig. 7.1). The charge and spin direction are identified in the $`𝔰𝔩_3`$ weight diagram as indicated in the figure. For other $`M`$ the analogous picture is obtained by ‘stretching’ the charge axis.
In the following sections we analyze the fqH data for the conformal field theory $`(\widehat{𝔰𝔩_3})_{k,M}`$. We first discuss the case $`k=1`$ (which corresponds to the abelian spin singlet Halperin state with parameters $`(M+2,M+2,M+1)`$ ) in some detail and then generalize to the non-abelian case $`k>1`$.
### 7.1. $`(\widehat{𝔰𝔩_3})_{k=1,M}`$
The exclusion statistics and UCPF character for the $`(\widehat{𝔰𝔩_3})_{k=1,M=0}`$ conformal field theory, in terms of the quasiparticles $`\{\varphi ^{ϵ_1},\varphi ^{ϵ_2},\varphi ^{ϵ_3}\}`$, were worked out in . Specializing to the subset $`\{\varphi _{},\varphi _{}\}=\{\varphi ^{ϵ_1},\varphi ^{ϵ_2}\}`$ we have
$$𝐊_\varphi =\frac{1}{3}\left(\begin{array}{cc}2& 1\\ 1& 2\end{array}\right),𝐭_\varphi =\left(\begin{array}{c}\frac{1}{3}\\ \frac{1}{3}\end{array}\right),𝐬_\varphi =\left(\begin{array}{c}1\\ 1\end{array}\right).$$
(7.2)
With (4.11) this leads to $`\nu =2/3`$ in agreement with (7.1). Applying the shift map (4.3), the fqH data for $`(\widehat{𝔰𝔩_3})_{k=1,M}`$ are thus given by
$$𝐊_\varphi ^M=𝒮_M𝐊_\varphi =\frac{1}{2M+3}\left(\begin{array}{cc}M+2& (M+1)\\ (M+1)& M+2\end{array}\right),$$
(7.3)
while
$$𝐭_\varphi ^M=\left(\begin{array}{c}\frac{1}{2M+3}\\ \frac{1}{2M+3}\end{array}\right),𝐬_\varphi ^M=\left(\begin{array}{c}1\\ 1\end{array}\right).$$
(7.4)
The IOW-equation for the total one-particle partition function $`\lambda _{\text{tot}}=\lambda _{}\lambda _{}`$, resulting from (7.3), is given by
$$\lambda _{\text{tot}}x_{}x_{}\lambda _{\text{tot}}^{\frac{2M+2}{2M+3}}(x_{}+x_{})\lambda _{\text{tot}}^{\frac{M+1}{2M+3}}1=0.$$
(7.5)
The $`K`$-matrix in the electron sector is determined as follows. First of all, the principal subspace of the $`(\widehat{𝔰𝔩_3})_{k=1,M=0}`$ integrable highest weight modules is generated by $`\{J_{\alpha _1},J_{\alpha _2}\}`$ and has a $`K`$-matrix given by (see App. B)
$$𝐊=\left(\begin{array}{cc}2& 1\\ 1& 2\end{array}\right).$$
(7.6)
The electron operators $`\{\mathrm{\Psi }_{},\mathrm{\Psi }_{}\}`$, however, are identified with $`\{J_{\alpha _2},J_{\alpha _3}\}`$. Interpreting $`J_{\alpha _3}`$ as the composite $`(J_{\alpha _1}J_{\alpha _2})`$, we can apply the construction of Sect. 4.4 and find an equivalent $`K`$-matrix for the combined $`\{J_{\alpha _1},J_{\alpha _2},J_{\alpha _3}\}`$ system
$$𝐊^{}=𝒞_{12}𝐊=\left(\begin{array}{ccc}2& 0& 1\\ 0& 2& 1\\ 1& 1& 2\end{array}\right).$$
(7.7)
Thus, we conclude that the electron fqH data are given by
$$𝐊_e=\left(\begin{array}{cc}2& 1\\ 1& 2\end{array}\right),𝐭_e=\left(\begin{array}{c}1\\ 1\end{array}\right),𝐬_e=\left(\begin{array}{c}1\\ 1\end{array}\right).$$
(7.8)
And thus, by applying the shift map
$$𝐊_e^M=𝒮_M𝐊_e=\left(\begin{array}{cc}M+2& M+1\\ M+1& M+2\end{array}\right),𝐭_e^M=\left(\begin{array}{c}1\\ 1\end{array}\right).$$
(7.9)
Note again that the fqH data in the electron and quasihole sectors, given in Eqs. (7.3), (7.4) and (7.9), are related by the duality (4.15).
The IOW-equation for $`\mu _{\text{tot}}=\mu _{}\mu _{}`$, resulting from (7.9), is given by
$$\mu _{\text{tot}}^{2M+3}\mu _{\text{tot}}^{2M+2}(y_{}+y_{})\mu _{\text{tot}}^{M+1}y_{}y_{}=0,$$
(7.10)
and is dual to (7.5) in the sense of (4.3). Explicitly,
$$\lambda _{\text{tot}}(x_{},x_{})=(y_{}y_{})^1\mu _{\text{tot}}(y_{},y_{})^{2M+3},$$
(7.11)
where
$$y_{}=x_{}^{(M+2)}x_{}^{(M+1)},y_{}=x_{}^{(M+1)}x_{}^{(M+2)}.$$
(7.12)
It remains to show that the fqH data $`(𝐊_\varphi ,𝐭_\varphi ,𝐬_\varphi )`$ and their duals $`(𝐊_e,𝐭_e,𝐬_e)`$ give a complete description of the chiral spectrum of the $`(\widehat{𝔰𝔩_3})_{k=1,M}`$ conformal field theory by constructing the $`(\widehat{𝔰𝔩_3})_{k=1,M}`$ characters in terms of (finite) linear combinations of UCPFs based on $`𝐊_e𝐊_\varphi `$. This is delegated to App. D. Here we only observe that, since there are no pseudoparticles, Eq. (3.10) immediately gives $`c_e+c_\varphi =2`$ which is the correct value of the central charge for $`(\widehat{𝔰𝔩_3})_{k=1,M}`$. Note also that $`c_\varphi `$ and $`c_e`$ separately depend on $`M`$ and are, in general, not simple rational numbers, e.g., for $`M=0`$ we have numerically $`c_e=0.6887`$ and $`c_\varphi =1.3113`$ while for $`M\mathrm{}`$ all the central charge is concentrated in the $`\varphi `$ sector.
Upon generalizing to higher levels $`k>1`$, it turns out we need an equivalent description of the system described above in terms of three quasihole operators, namely by adding a quasihole operator $`\varphi ^{ϵ_3}`$ of $`𝔰𝔩_3`$ weight $`ϵ_3`$, i.e., of charge $`2/3`$ (for $`M=0`$) and spinless. The $`K`$-matrix for this system can be obtained as a submatrix of the $`K`$-matrix describing quasiparticles in the $`\mathrm{𝟑}\mathrm{𝟑}^{}`$ of $`𝔰𝔩_3`$ or, equivalently, by using that $`\varphi ^{ϵ_3}`$ is the composite $`(\varphi ^{ϵ_1}\varphi ^{ϵ_2})`$ and using (4.22). We find
$$𝐊_\varphi ^{}_{}{}^{}M=𝒞_{12}𝐊_\varphi ^M=\frac{1}{2M+3}\left(\begin{array}{ccc}M+2& M+2& 1\\ M+2& M+2& 1\\ 1& 1& 2\end{array}\right),𝐭_\varphi ^{}_{}{}^{}M=\left(\begin{array}{c}\frac{1}{2M+3}\\ \frac{1}{2M+3}\\ \frac{2}{2M+3}\end{array}\right).$$
(7.13)
In the electron sector we can similarly introduce the composite $`(J_{\alpha _2}J_{\alpha _3})`$ and obtain
$$𝐊_e^{}_{}{}^{}M=𝒞_{12}𝐊_e^M=\left(\begin{array}{ccc}M+2& M+2& 2M+3\\ M+2& M+2& 2M+3\\ 2M+3& 2M+3& 4M+6\end{array}\right),𝐭_e^{}=\left(\begin{array}{c}1\\ 1\\ 2\end{array}\right).$$
(7.14)
Now we observe a curiosity; while obviously the fqH data (7.13) and (7.14) are dual, since they are equivalent to the dual systems given in (7.3) and (7.9), they are not related by the duality transformation (4.15) because both $`𝐊_\varphi `$ and $`𝐊_e`$ are not invertible. The equivalence can also be observed at the level of the resulting IOW-equations which are now given by
$`(\lambda _{\text{tot}}^{\frac{1}{2M+3}}x_{}x_{})(\lambda _{\text{tot}}x_{}x_{}\lambda _{\text{tot}}^{\frac{2M+2}{2M+3}}(x_{}+x_{})\lambda _{\text{tot}}^{\frac{M+1}{2M+3}}1)`$ $`=`$ $`0,`$
$`(\mu _{\text{tot}}^{2M+3}y_{}y_{})(\mu _{\text{tot}}^{2M+3}\mu _{\text{tot}}^{2M+2}(y_{}+y_{})\mu _{\text{tot}}^{M+1}y_{}y_{})`$ $`=`$ $`0.`$ (7.15)
Because of the first factor the equations (7.1) do not transform into eachother under (7.11). However, the physical solutions, which are determined by the second factor, do! Summarizing, we conclude that it is obvious that the notion of duality should have an extension that incorporates non-invertible $`K`$-matrices. We leave this for future investigation.
### 7.2. $`(\widehat{𝔰𝔩_3})_{k,M}`$
As argued in , the generalization of the results of the previous section to levels $`k>1`$ requires the addition of $`2(k1)`$ pseudoparticles incorporating the non-abelian statistics of the quasihole operators $`\{\varphi _{},\varphi _{}\}`$. Since these pseudoparticles couple differently to $`\{\varphi _{},\varphi _{}\}`$ than to the composite particle $`\varphi _{}=(\varphi _{}\varphi _{})`$ (i.e., different than the naive coupling given by the composite construction), it appears that the first construction in Sect. 7.1 does not generalize to higher levels.
It is known that for $`(\widehat{𝔰𝔩_n})_{k,M=0}`$ the pseudoparticles couple to the physical particles by means of the matrix $`𝐀_{n1}^1𝐀_k`$. Here we have used the result for the restricted Kostka polynomials as given in, e.g., (see the discussion in for details). Then, by applying the shift map (4.3), we obtain
$$𝐊_\varphi ^M=\left(\begin{array}{ccccc}& & & & \mathrm{}\\ & & & & \mathrm{}& \\ & & 𝐀_2^1𝐀_{k1}& & \mathrm{}& {\scriptscriptstyle \frac{2}{3}}& {\scriptscriptstyle \frac{2}{3}}& {\scriptscriptstyle \frac{1}{3}}\\ & & & & \mathrm{}& {\scriptscriptstyle \frac{1}{3}}& {\scriptscriptstyle \frac{1}{3}}& {\scriptscriptstyle \frac{2}{3}}\\ [1.5]8\\ & & {\scriptscriptstyle \frac{2}{3}}& {\scriptscriptstyle \frac{1}{3}}& \mathrm{}& {\scriptscriptstyle \frac{(4k1)M+6}{3(2kM+3)}}& {\scriptscriptstyle \frac{(4k1)M+6}{3(2kM+3)}}& {\scriptscriptstyle \frac{(2k2)M+3}{3(2kM+3)}}\\ & & {\scriptscriptstyle \frac{2}{3}}& {\scriptscriptstyle \frac{1}{3}}& \mathrm{}& {\scriptscriptstyle \frac{(4k1)M+6}{3(2kM+3)}}& {\scriptscriptstyle \frac{(4k1)M+6}{3(2kM+3)}}& {\scriptscriptstyle \frac{(2k2)M+3}{3(2kM+3)}}\\ & & {\scriptscriptstyle \frac{1}{3}}& {\scriptscriptstyle \frac{2}{3}}& \mathrm{}& {\scriptscriptstyle \frac{(2k2)M+3}{3(2kM+3)}}& {\scriptscriptstyle \frac{(2k2)M+3}{3(2kM+3)}}& {\scriptscriptstyle \frac{(4k4)M+6}{3(2kM+3)}}\end{array}\right),$$
(7.16)
where the components of $`𝐀_2`$ refer to the quasiholes in the $`\mathrm{𝟑}`$ and $`\mathrm{𝟑}^{}`$, respectively, and
$$𝐭_\varphi =(\underset{2(k1)}{\underset{}{0,0,\mathrm{},0}}|\frac{1}{2kM+3},\frac{1}{2kM+3},\frac{2}{2kM+3}).$$
(7.17)
For instance, for level $`k=2`$ we have
$$𝐊_\varphi ^M=\left(\begin{array}{cccccc}\frac{4}{3}& \frac{2}{3}& \mathrm{}& \frac{2}{3}& \frac{2}{3}& \frac{1}{3}\\ \frac{2}{3}& \frac{4}{3}& \mathrm{}& \frac{1}{3}& \frac{1}{3}& \frac{2}{3}\\ [1.5]6\\ \frac{2}{3}& \frac{1}{3}& \mathrm{}& \frac{7M+6}{12M+9}& \frac{7M+6}{12M+9}& \frac{2M+3}{12M+9}\\ \frac{2}{3}& \frac{1}{3}& \mathrm{}& \frac{7M+6}{12M+9}& \frac{7M+6}{12M+9}& \frac{2M+3}{12M+9}\\ \frac{1}{3}& \frac{2}{3}& \mathrm{}& \frac{2M+3}{12M+9}& \frac{2M+3}{12M+9}& \frac{4M+6}{12M+9}\end{array}\right).$$
(7.18)
Note that the matrix $`𝐊_\varphi ^M`$ of (7.16) is not invertible, as was observed for $`k=1`$ in Sect. 7.1. Thus, we cannot simply identify the dual sector by performing the transformation (4.15).
To obtain the dual sector we proceed as in Sect. 7.1. We start with the $`K`$-matrix of the principal subspace spanned by $`\{J_{\alpha _1},J_{\alpha _2}\}`$. As discussed in App. B, for $`(\widehat{𝔰𝔩_3})_k`$, this $`K`$-matrix is given by $`𝐊=𝐀_2𝐁_k`$ and requires, besides the currents $`\{J_{\alpha _1},J_{\alpha _2}\}`$ a set of $`2(k1)`$ composites
$$(\underset{l}{\underset{}{J_{\alpha _i}\mathrm{}J_{\alpha _i}}}),2lk,i=1,2.$$
(7.19)
Starting with this matrix we introduce additional composites according to the procedure of Sect. 4.4, beginning with the electron operator $`\mathrm{\Psi }_{}=(J_{\alpha _1}J_{\alpha _2})`$ (recall that $`\mathrm{\Psi }_{}=J_{\alpha _2}`$), and continuing until all composites
$$(\underset{n_{}}{\underset{}{\mathrm{\Psi }_{}\mathrm{}\mathrm{\Psi }_{}}}\underset{n_{}}{\underset{}{\mathrm{\Psi }_{}\mathrm{}\mathrm{\Psi }_{}}}),n_{}+n_{}k,$$
(7.20)
have been introduced. Note that the set of composites (7.20), for fixed $`n_{}+n_{}`$, span a $`(n_{}+n_{}+1)`$-dimensional irreducible representation of spin $`SU(2)`$. The electron $`K`$-matrix is then the $`\frac{1}{2}k(k+3)\times \frac{1}{2}k(k+3)`$ submatrix of the resulting $`𝐊`$ obtained by omitting the composites which cannot be written in terms of electron operators only. Let us be illustrate this procedure the case of $`k=2`$. Starting with the principal subspace $`K`$-matrix
$$𝐊=\left(\begin{array}{ccccc}2& 1& \mathrm{}& 2& 1\\ 1& 2& \mathrm{}& 1& 2\\ [1.5]5\\ 2& 1& \mathrm{}& 4& 2\\ 1& 2& \mathrm{}& 2& 4\end{array}\right),$$
(7.21)
we introduce, respectively, the composites $`\mathrm{\Psi }_{}=(J_{\alpha _1}J_{\alpha _2})`$, $`(J_{\alpha _2}(J_{\alpha _1}J_{\alpha _1}))`$, $`(J_{\alpha _2}(J_{\alpha _1}J_{\alpha _2}))`$, and $`(J_{\alpha _2}((J_{\alpha _2}(J_{\alpha _1}J_{\alpha _1})))`$. Then, after removing the rows and columns corresponding to $`J_{\alpha _1}`$, $`(J_{\alpha _1}J_{\alpha _1})`$ and $`(J_{\alpha _2}(J_{\alpha _1}J_{\alpha _1}))`$, we obtain
$$𝐊_e^{}=\left(\begin{array}{cccccc}2& 1& \mathrm{}& 2& 2& 1\\ 1& 2& \mathrm{}& 1& 2& 2\\ [1.5]6\\ 2& 1& \mathrm{}& 4& 3& 2\\ 2& 2& \mathrm{}& 3& 4& 3\\ 1& 2& \mathrm{}& 2& 3& 4\end{array}\right),𝐭_e=\left(\begin{array}{c}1\\ 1\\ 2\\ 2\\ 2\end{array}\right),𝐬_e=\left(\begin{array}{c}1\\ 1\\ 2\\ 0\\ 2\end{array}\right).$$
(7.22)
Similarly, one obtains the electron $`K`$-matrix for $`(\widehat{𝔰𝔩_3})_{k,M=0}`$ at higher levels, and the generalization to arbitrary $`M`$ follows, as before, by applying the shift map (4.3). Unfortunately, the procedure described above is ambiguous. The resulting $`K`$-matrix depends on the order in which the composites are taken as well as the precise identification of the clusters (7.20) with the original clusters (7.19), e.g., should we identify $`(\mathrm{\Psi }_{}\mathrm{\Psi }_{})`$ with $`(J_{\alpha _1}(J_{\alpha _1}(J_{\alpha _2}J_{\alpha _2})))`$ or $`((J_{\alpha _1}J_{\alpha _1})(J_{\alpha _2}J_{\alpha _2}))`$? Ultimately, the ‘correct’ matrix $`𝐊_e`$ is selected by the requirement that the complete spectrum can be build out of the quasihole and electron operators or, more concretely, that the characters of $`(\widehat{𝔰𝔩_3})_{k,M}`$ can be written as a linear combination of UCPFs based on $`𝐊_\varphi 𝐊_e`$. A nontrivial (and highly selective) check is whether the central charge, given by (4.7), works out correctly, i.e., whether $`c_\varphi +c_e=8k/(k+3)`$, for the $`K`$-matrices (7.16) and the ‘appropriate’ generalization of (7.22) to higher levels and arbitrary $`M`$. We have checked this numerically for low values of $`k`$ and $`M`$ as well as exactly, for all $`k`$, in the $`M\mathrm{}`$ limit, in which case the central charge is entirely concentrated in the $`\varphi `$-sector. We refrain from giving the explicit matrices $`𝐊_e`$ until we have performed an additional simplifying reduction.
First observe that, for $`k=2`$, the matrix $`𝐊_e^{}`$ of Eq. (7.22) is invertible, in contrast to the matrix $`𝐊_\varphi ^M`$ of (7.18). One could therefore simply have started with $`𝐊_e^{}`$ and have obtained the dual sector by the duality transformations (4.15). This would result in a $`\varphi `$-sector, different from the one discussed above, with two physical quasiholes and three pseudoparticles. Unfortunately, this procedure breaks down, in general, for higher $`k`$ as the matrices $`𝐊_e`$, constructed according to the procedure outlined above, are no longer invertible. However, note that the matrix (7.22) can be reduced to an equivalent $`4\times 4`$ matrix by inverting the composite procedure – in this case by removing $`(\mathrm{\Psi }_{}\mathrm{\Psi }_{})`$ in the fourth column, since this column can be created by applying $`𝒞_{12}`$. This procedure works for general $`k>1`$ and leads to a $`2k\times 2k`$ electron $`K`$-matrix, for the composites (7.19) with either $`n_{}=0`$ or $`n_{}=0`$ (i.e. we lose the SU(2) multiplet structure), given by
$$𝐊_e=\left(\begin{array}{cccccccccc}2& 0& 2& 0& \mathrm{}& & 2& 0& 2& 1\\ 0& 2& 0& 2& \mathrm{}& & 0& 2& 1& 2\\ 2& 0& 4& 0& & & 4& 1& 4& 2\\ 0& 2& 0& 4& & & 1& 4& 2& 4\\ \mathrm{}& \mathrm{}& & & & & & & \mathrm{}& \mathrm{}\\ 2& 0& 4& 1& & & 2(k1)& k2& 2(k1)& k1\\ 0& 2& 1& 4& & & k2& 2(k1)& k1& 2(k1)\\ 2& 1& 4& 2& \mathrm{}& & 2(k1)& k1& 2k& k\\ 1& 2& 2& 4& \mathrm{}& & k1& 2(k1)& k& 2k\end{array}\right),$$
(7.23)
and
$`𝐭_e`$ $`=`$ $`(1,1;2,2;\mathrm{};k,k),`$
$`𝐬_e`$ $`=`$ $`(1,1;2,2;\mathrm{};k,k).`$ (7.24)
The generalization $`𝐊_e^M`$ to arbitrary $`M`$ follows by applying the shift map, in this case by adding the matrix $`M(\mathbb{𝟙}_2𝐃)`$ where $`\mathbb{𝟙}_2`$ is the identity matrix in two dimensions and $`(𝐃)_{ij}=ij`$ ($`1i,jk`$) (see for an explicit expression in the case $`k=2`$). This matrix is invertible, so we simply define $`𝐊_\varphi ^M=(𝐊_e^M)^1`$. A convenient permutation of rows and columns of $`𝐊_\varphi ^M`$ leads to the following matrix
$$(𝐊_\varphi ^M)^{\text{perm}}=\left(\begin{array}{cccccccc}& & & & & \mathrm{}& 0& \frac{1}{3}\\ & & & & & \mathrm{}& 0& \frac{2}{3}\\ & & 𝐀_2^1𝐀_{k1}& & & \mathrm{}& & \\ & & & & & \mathrm{}& \frac{2}{3}& 0\\ & & & & & \mathrm{}& \frac{1}{3}& 0\\ [1.5]8\\ 0& 0& & \frac{2}{3}& \frac{1}{3}& \mathrm{}& \frac{(4k1)M+6}{3(2kM+3)}& \frac{M}{3(2kM+3)}\\ \frac{1}{3}& \frac{2}{3}& & 0& 0& \mathrm{}& \frac{M}{3(2kM+3)}& \frac{(4k1)M+6}{3(2kM+3)}\end{array}\right),$$
(7.25)
containing two physical particles and $`2(k1)`$ pseudoparticles. Also,
$`𝐭_\varphi `$ $`=`$ $`(0,0;0,0;\mathrm{};\frac{1}{2kM+3},\frac{1}{2kM+3}),`$
$`𝐬_\varphi `$ $`=`$ $`(0,0;0,0;\mathrm{};1,1),`$ (7.26)
as one would expect. We have checked that the total central charge $`c_e+c_\varphi `$ for Eqs. (7.23) and (7.25) works out correctly, namely $`c_e+c_\varphi =8k/(k+3)`$. Moreover, we have checked for low values of $`k`$ that the equation for $`\lambda _{\text{tot}}`$, resulting from the IOW equations based on (7.25), are identical to those based on (7.16). Furthermore, in all formulations, the equations (4.10) and (4.11) are consistent with (7.1).
For $`k=2,3`$, we checked the small $`x`$ behaviour for $`\lambda _{\text{tot}}`$, Eq. (4.5). We again expect the constants $`\alpha `$ to be the quantum dimensions of the associated conformal field theory. Using some results in , these quantum dimensions are given by
$$\alpha _k=1+2\mathrm{cos}\left(\frac{2\pi }{k+3}\right).$$
(7.27)
For $`k=2`$, the equation for $`\lambda _{\text{tot}}`$ reads (upon taking $`x_{}=x_{}=x`$)
$$(\lambda _{\text{tot}}^{\frac{1}{2}}1)^2=x^2\lambda _{\text{tot}}^{\frac{8M+5}{8M+6}}+x\lambda _{\text{tot}}^{\frac{6M+4}{8M+6}}x\lambda _{\text{tot}}^{\frac{2M+1}{8M+6}},$$
(7.28)
which leads to the following small $`x`$ behaviour
$$\lambda _{\text{tot}}=1+2\left(\frac{1+\sqrt{5}}{2}\right)x+o(x^2),$$
(7.29)
in agreement with $`\alpha _2=(1+\sqrt{5})/2`$ from (7.27); the extra factor $`2`$ comes from the sum over the two physical particles, see Eq. (4.5). For $`k=3`$ we find
$$(\lambda _{\text{tot}}^{\frac{1}{2}}1)=x\lambda _{\text{tot}}^{\frac{8M+3}{6(6M+3)}}(\lambda _{\text{tot}}^{\frac{1}{6}}+1)^{\frac{1}{3}}(\lambda _{\text{tot}}^{\frac{1}{3}}+1)^{\frac{2}{3}},$$
(7.30)
which gives $`\alpha _3=2`$, consistent with (7.27). Note that for the abelian case $`k=1`$, we find for the small $`x_,`$-behaviour, using (7.5),
$$\lambda _{\text{tot}}=1+(x_{}+x_{})+o(x^2),$$
(7.31)
in agreement with (7.27) and the fact that for $`k=1`$ we have an abelian state.
As was the case for the spin polarized states of Sect. 6, also for the non-abelian spin singlet states a particle behaving as a supercurrent can be identified in the non-magnetic limit. The situation here is slightly more complicated than in the case of the spin polarized states discussed in Sect. 6. This is because in the formulation above, there is no candidate particle with the property that all the statistics parameters go to zero in the limit $`q0`$ (with $`q=1/\nu =M+3/2k`$). However, if one acts with $`𝒞_{2k1,2k}`$ on $`𝒮_M𝐊_e`$, with $`𝐊_e`$ given by Eq. (7.23), one introduces a composite with charge $`2k`$ and spin $`0`$, which has the desired properties. In the $`\varphi `$-sector, the particle content is changed to one quasihole and $`2k`$ pseudoparticles, of which a few carry spin.
The possibility to introduce a composite with the right properties enables one to repeat the discussion of Sect. 6, with the only difference that the flux quantum in this case equals $`h/2ke`$. So, also in this case, we can identify a supercurrent in the non-magnetic limit.
## 8. Discussion
In this paper we derived the $`K`$-matrix structure for two classes of so called non-abelian quantum Hall states, putting the results of on a firmer basis. In doing so, we extensively made use of a duality between the edge electron and quasihole excitations. The abelian formalism was extended to include electron spin, in order to be able to treat spin singlet states. Moreover, we showed that many results of the abelian $`K`$-matrix formulation for hierarchy states also hold for our generalized $`K`$-matrices, thereby justifying their name. We would like to stress that the non-abelian states of are not hierarchical states; the $`K`$-matrix structure is necessary as a bookkeeping device for the non-abelian statistics.
An important concept we did not discuss is the torus degeneracy ; it is not clear at the moment how to generalize this to the non-abelian case (some remarks are made in Appendix D). Another important issue to be settled has to do with the cases where the pseudoparticles do carry spin (or charge). These may arise by creating extra composites in the electron sector; by the duality, the $`\varphi `$ sector changes accordingly, and pseudoparticles carrying spin may arise. The formulas Eq. (2) then need a proper adjustment, because they do not give the same result any more, and the physical quantities like the filling factors need to be invariant under the introduction of extra composites. We would like to remark that a description in which the pseudoparticles do not carry spin or charge is possible in the cases we examined, and the various physical quantities were obtained correctly.
As for the Laughlin wave functions, one would like to have a Landau-Ginzburg field theory describing the excitations for the non-abelian states. The backbone of such a theory will be a Chern-Simons term, in which the gauge fields are coupled in a special way. We expect that the $`K`$-matrices derived in this paper will play a crucial role. From a Landau Ginzburg theory (using the $`K`$-matrices etc. from the electronic sector), one should be able to identify the possible excitations in the $`\varphi `$-sector, as vortex solutions of the classical equations of motion. Identifying this Landau-Ginzburg theory is left for future investigations (see for related studies).
Another interesting issue for the non-abelian states is the determination of the degeneracies of the states when extra flux is applied through the sample. These degeneracies can be calculated using conformal field theory techniques, and can, interestingly, be simulated on a computer using a special, ultra local, interaction for the electron interaction. For the Pfaffian, exact counting results were obtained in ; the more general Read-Rezayi states were treated in . Counting results for the NASS states will be given elsewhere .
Finally, while our discussion of fqH-bases of conformal field theories based on quasiparticles with a statistics matrix $`𝐊𝐊^1`$ was restricted to $`(\widehat{𝔰𝔩_n})_k`$ (for $`n=2,3`$), it is obvious that such a description generalizes to more general conformal field theories (see for more examples), even though these may not have an interpretation in the context of the fractional quantum Hall effect.
## Acknowledgments
We would like to thank Sathya Guruswamy, Nick Read and Ole Warnaar for useful discussions. This research was financially supported in part by the foundation FOM of the Netherlands and the Australian Research Council.
## Appendix A Basic hypergeometric series
Consider the basic hypergeometric series
$$\begin{array}{c}{}_{r}{}^{}\varphi _{s}^{}(a_1,\mathrm{},a_r;b_1,\mathrm{},b_s;q,z)=\hfill \\ \hfill \underset{m0}{}\frac{(a_1;q)_m(a_2;q)_m\mathrm{}(a_r;q)_m}{(q;q)_m(b_1;q)_m\mathrm{}(b_s;q)_m}\left((1)^mq^{\frac{1}{2}m(m1)}\right)^{1+sr}z^m,\end{array}$$
(A.1)
where
$$(a;q)_n=\underset{k=0}{\overset{n1}{}}(1aq^k).$$
(A.2)
We have the $`q`$-Pfaff-Saalschütz sum
$${}_{3}{}^{}\varphi _{2}^{}(a,b,q^n;c,abq^{1n}/c;q,q)=\frac{(c/a;q)_n(c/b;q)_n}{(c;q)_n(c/ab;q)_n},$$
(A.3)
Taking $`b=0`$ in (A.3) gives the $`q`$-Chu-Vandermonde sum
$${}_{2}{}^{}\varphi _{1}^{}(a,q^n;c;q,q)=\frac{(c/a;q)_n}{(c;q)_n}a^n.$$
(A.4)
Now, taking $`a=q^{m_1}`$, $`b=q^{M_1+M_2(m_1+m_1)+1}`$, $`c=q^{M_2(m_1+m_2)+1}`$ and $`n=m_2`$ in (A.3) gives
$$\begin{array}{c}\left[\genfrac{}{}{0.0pt}{}{M_1}{m_1}\right]\left[\genfrac{}{}{0.0pt}{}{M_2}{m_2}\right]=\hfill \\ \hfill \underset{m0}{}q^{(m_1m)(m_2m)}\left[\genfrac{}{}{0.0pt}{}{M_1m_2}{m_1m}\right]\left[\genfrac{}{}{0.0pt}{}{M_2m_1}{m_2m}\right]\left[\genfrac{}{}{0.0pt}{}{M_1+M_2(m_1+m_2)+m}{m}\right].\end{array}$$
(A.5)
Taking $`a=q^{m_1}`$, $`c=q^{M_2(m_1+m_2)+1}`$ and $`n=m_2`$ in (A.4) gives
$$\frac{1}{(q)_{m_1}}\left[\genfrac{}{}{0.0pt}{}{M_2}{m_2}\right]=\underset{m0}{}q^{(m_1m)(m_2m)}\frac{1}{(q)_m(q)_{m_1m}}\left[\genfrac{}{}{0.0pt}{}{M_2m_1}{m_2m}\right],$$
(A.6)
while taking $`a=q^{m_1}`$, $`n=m_2`$ and $`c=0`$ in (A.4) gives
$$\frac{1}{(q)_{m_1}(q)_{m_2}}=\underset{m0}{}q^{(m_1m)(m_2m)}\frac{1}{(q)_m(q)_{m_1m}(q)_{m_2m}}.$$
(A.7)
## Appendix B The principal subspace
In this appendix we review an important result of which is used throughout the paper. Consider an affine Lie algebra $`\widehat{𝔤}_k`$ (see, e.g., for notation and definitions). If $`L(\mathrm{\Lambda })`$ is the integrable highest weight module of $`\widehat{𝔤}_k`$ with highest weight $`\mathrm{\Lambda }`$ and highest weight vector $`v_\mathrm{\Lambda }`$, then the principal subspace $`W(\mathrm{\Lambda })L(\mathrm{\Lambda })`$ is defined to be the subspace generated from $`v_\mathrm{\Lambda }`$ by the negative modes of the positive simple root currents $`J_{\alpha _i}(z)`$.
The character of the principal subspace $`W(\mathrm{\Lambda })`$ of the integrable highest weight module $`L(\mathrm{\Lambda })`$ for $`\mathrm{\Lambda }=k_0\mathrm{\Lambda }_0+k_j\mathrm{\Lambda }_j`$ ($`1jn`$, $`k_0+k_j=k`$) of $`(\widehat{𝔰𝔩_{n+1}})_k`$ was determined in .<sup>2</sup><sup>2</sup>2The character of the principal subspace $`W(\mathrm{\Lambda })`$ for more general level $`k`$ modules $`L(\mathrm{\Lambda })`$ is apparently not yet known. It is given by the UCPF
$$\mathrm{ch}_W=\underset{𝐩}{}\left(z_i^{sp_i^{(s)}}\right)\frac{q^{\frac{1}{2}𝐩𝐊𝐩+𝐐_j𝐩}}{_i_s(q)_{p_i^{(s)}}},$$
(B.1)
where
$$𝐊=𝐀_n𝐁_k,𝐐_j=𝐞_j(\underset{k_0}{\underset{}{0,\mathrm{},0}},1,2,\mathrm{},k_j),$$
(B.2)
and $`z_i`$ denotes the (generalized) fugacity of the current $`J_{\alpha _i}`$. Also, $`(𝐀_n)_{ij}=2\delta _{ij}\delta _{i1,j}\delta _{i+1,j}`$ is the Cartan matrix of $`𝔰𝔩_{n+1}`$, i.e.
$$𝐀_n=\left(\begin{array}{ccc}2& 1& \\ 1& 2& 1& \\ & 1& 2& 1& \\ & & \mathrm{}& \mathrm{}& \mathrm{}& \\ & & & 1& 2& 1\\ & & & & 1& 2\end{array}\right),$$
(B.3)
and $`(𝐁_k)_{rs}=\mathrm{min}(r,s)_{|r,s=1,\mathrm{},k}`$, i.e.
$$𝐁_k=\left(\begin{array}{ccccc}1& 1& 1& \mathrm{}& 1\\ 1& 2& 2& \mathrm{}& 2\\ 1& 2& 3& \mathrm{}& 3\\ \mathrm{}& & & \mathrm{}& \mathrm{}\\ 1& 2& 3& \mathrm{}& k\end{array}\right).$$
(B.4)
Furthermore, in (B.1), we have written $`𝐩=(p_j^{(s)})_{j=1,\mathrm{},n}^{s=1,\mathrm{},k}`$ with respect to $`(𝐀_n)_{ij}(𝐁_k)_{rs}`$.
## Appendix C $`(\widehat{𝔰𝔩_2})_{k,M}`$ character
The UCPF character for $`(\widehat{𝔰𝔩_2})_{k=1,M}`$ was discussed in (see also ). Here we discuss the $`q`$-Pfaffian case, i.e. $`k=2`$. For convenience we put $`q=M+1`$.
### C.1. Quasihole sector
In , finitized partition sums $`X_L=X_{l=\frac{8Lq2}{16q}}`$ and $`Y_L=Y_{l=\frac{8L+q6}{16q}}`$ for the quasihole sector of the $`q`$-pfaffian CFT were introduced. $`X_L`$ ($`Y_L`$) are restricted by requiring that the total charge be an even (odd) multiple of $`\frac{1}{2q}`$. In , it was established that the following recursion relations hold
$`X_L`$ $`=`$ $`X_{L2q}+xq^{\frac{8Lq2}{16q}}(Y_L+Y_{Lq}),`$
$`Y_L`$ $`=`$ $`Y_{L2q}+xq^{\frac{8L+q6}{16q}}X_{L1},`$ (C.1)
or, equivalently,
$$X_L=X_{L2q}+q^{\frac{1}{2}}\left(X_{Lq}X_{L3q}\right)+x^2q^{\frac{2L1}{2q}}X_{L1}.$$
(C.2)
By putting $`X_L/X_{Lq}\lambda _{\text{tot}}^{1/2q}`$, for large $`L`$, we reproduce equation (6.4). To build the entire spectrum of the $`(\widehat{𝔰𝔩_2})_{k=2,M}`$ conformal field theory we need $`3q`$ sectors whose initial conditions are given in Table C.1. The vacua of the sectors are labeled by, respectively, charge and the $`𝔰𝔩_2`$ irrep in which they appear for $`M=0`$ (the labels $`\mathbb{𝟙}`$, $`\sigma `$ and $`\psi `$ stand for the $`𝔰𝔩_2`$ singlet, doublet and triplet, respectively, in analogy with the Ising model). The parameter $`r`$ takes the values $`r=0,1,\mathrm{},q1`$.
The solutions to (C.1) can be written in terms of finitized UCPFs with (cf. (6.3))
$$𝐊_\varphi =\left(\begin{array}{cc}1& \frac{1}{2}\\ \frac{1}{2}& \frac{q+1}{4q}\end{array}\right).$$
(C.3)
Indeed, the recursion relations (5.12), with $`𝐊=𝐊_\varphi `$ and $`𝐐+𝐮=0`$, are explicitly given by
$`P_{L_1,L_2}`$ $`=`$ $`P_{L_11,L_2}+q^{L_1\frac{1}{2}}P_{L_11,L_2+\frac{1}{2}},`$
$`P_{L_1,L_2}`$ $`=`$ $`P_{L_1,L_21}+xq^{L_1\frac{q+1}{8q}}P_{L_1+\frac{1}{2},L_2\frac{q+1}{4q}},`$ (C.4)
and lead to (C.1) upon identifying
$$X_L=q^{\frac{1}{4}Q_1^2}P_{0,\frac{L}{2q}},Y_L=q^{\frac{1}{4}Q_1^2\frac{1}{16}}P_{\frac{1}{2},\frac{L}{2q}+\frac{q1}{4q}}.$$
(C.5)
The values for $`𝐐_\varphi `$ in each sector are listed in Table C.1, while the parameters $`s=0,\mathrm{},2q1`$, in Table C.1, are given in Table C.2.
### C.2. Electron sector
For the electron sector of the $`q`$-pfaffian, the paper introduced the truncated partition sums $`\mathrm{\Omega }_L`$, which contain all states constructed from the edge electron operator $`\mathrm{\Psi }_s`$ with $`sL\frac{q1}{2}`$. It satisfies the recursion relation
$$\mathrm{\Omega }_L=\mathrm{\Omega }_{L1}+yq^{L\frac{1}{2}(q1)}\mathrm{\Omega }_{Lq}+y^2q^{2L(2q1)}\mathrm{\Omega }_{L2q}y^3q^{3L\frac{1}{2}(9q5)}\mathrm{\Omega }_{L3q},$$
(C.6)
and results in Eq. (6.7) by putting $`\mathrm{\Omega }_L/\mathrm{\Omega }_{L1}\mu _{\text{tot}}`$ for large $`L`$. In this case the recursion relation does not appear to be solved by finitized UCPFs. However, there exists a recursion relation, leading to the same equation for $`\mu _{\text{tot}}`$, that is solved by a finitized UCPF and differs from the solution to (C.6) by terms of order $`q^L`$ (i.e. belongs to the same universality class, see the discussion in Sect. 5.3) and thus gives the correct solution in the limit $`L\mathrm{}`$. The UCPF is based on the $`K`$-matrix (cf. (6.6))
$$𝐊_e=\left(\begin{array}{cc}q+1& 2q\\ 2q& 4q\end{array}\right).$$
(C.7)
The initial conditions and values for $`𝐐_e`$ in each sector are listed in Table C.3.
There is a slight subtlety in the case of the sectors $`|re/q,\psi `$. These vectors do not correspond to an extremal vector in the $`(\widehat{𝔰𝔩_2})_{k=2,M}`$ modules. Thus the results of App. B do not apply. While the exclusion statistics of the currents is unchanged, and hence the $`K`$-matrix is still given by (C.7), it can easily be shown that the extremal vectors in the modules cannot be reproduced by any two dimensional vector $`𝐐`$. In fact, to correctly reproduce the extremal vectors one needs an infinite dimensional vector $`𝐐`$ (given in Table C.3) with a corresponding infinite dimensional $`K`$-matrix $`𝐊_e^{(\mathrm{})}`$ that is equivalent to (C.7) by the composite construction. Specifically, one introduces derived matrices $`𝐊_e^{(n)}`$ and associated generalized fugacities $`𝐳^{(n)}`$ by
$$𝐊_e^{(1)}=𝒞_{12}𝐊_e=\left(\begin{array}{ccc}q+1& 2q+1& 3q+1\\ 2q+1& 4q& 6q\\ 3q+1& 6q& 9q+1\end{array}\right),𝐳^{(1)}=\left(\begin{array}{c}z\\ z^2\\ z^3\end{array}\right),$$
(C.8)
$$𝐊_e^{(2)}=𝒞_{23}𝐊_e^{(1)}=\left(\begin{array}{cccc}q+1& 2q+1& 3q+1& 5q+2\\ 2q+1& 4q& 6q+1& 10q\\ 3q+1& 6q+1& 9q+1& 15q+1\\ 5q+2& 10q& 15q+1& 25q+1\end{array}\right),𝐳^{(2)}=\left(\begin{array}{c}z\\ z^2\\ z^3\\ z^5\end{array}\right),$$
(C.9)
and, ultimately,
$$\begin{array}{c}𝐊_e^{(\mathrm{})}=\underset{n\mathrm{}}{lim}𝐊_e^{(n)}=\underset{n\mathrm{}}{lim}𝒞_{2,n}𝒞_{2,n1}\mathrm{}𝒞_{23}𝒞_{12}𝐊_e\hfill \\ \hfill =\left(\begin{array}{cccccccc}q+1& 2q+1& 3q+1& 5q+2& 7q+3& \mathrm{}& (2k+1)q+k& \mathrm{}\\ 2q+1& 4q& 6q+1& 10q+1& 14q+1& \mathrm{}& 2(2k+1)q+1& \mathrm{}\\ 3q+1& 6q+1& 9q+1& 15q+1& 21q+2& \mathrm{}& & \\ 5q+2& 10q+1& 15q+1& 25q+1& 35q+1& \mathrm{}& & \\ 7q+3& 14q+1& 21q+2& 35q+1& 49q+1& \mathrm{}& & \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ & & & & & & (2k+1)^2q+1& \\ & & & & & & & \mathrm{}\end{array}\right),\end{array}$$
(C.10)
while
$$𝐳^{(\mathrm{})}=(z,z^2;z^3,z^5,z^7,\mathrm{}).$$
(C.11)
For every finite $`n`$, the UCPF based on $`(𝐊_e^{(n)};𝐐_e^{(n)})`$, where $`𝐐_e^{(n)}`$ is the $`(n+2)`$-dimensional truncation of the vector $`𝐐_e`$ in Table C.3, gives an accurate description of the module up to some level (which appears to be at at least polynomially increasing with $`n`$). To describe the entire module accurately, one needs to take the limit $`n\mathrm{}`$.
### C.3. The character
Combining the $`3q`$ sectors in Tables C.1 and C.2 should reproduce the spectrum of the chiral $`(\widehat{𝔰𝔩_2})_{2,M}`$ conformal field theory. Consider the combination of UCPFs
$$Z_{\text{tot}}=\underset{k=1}{\overset{3q}{}}a_{(k)}Z_{\mathrm{}}(𝐊_e;𝐐_e^{(k)})Z_{\mathrm{}/2}(𝐊_\varphi ;𝐐_\varphi ^{(k)},𝐮_\varphi ^{(k)}),$$
(C.12)
where the coefficients $`a_{(k)}`$ and vectors $`𝐐_e^{(k)}`$, $`𝐐_\varphi ^{(k)}=𝐮_\varphi ^{(k)}`$ are given in Table C.4, and where
$$\begin{array}{c}Z_{\mathrm{}/2}(𝐊_\varphi ;𝐐_\varphi ^{(k)},𝐮_\varphi ^{(k)})\hfill \\ \hfill q^{\frac{1}{4}(Q_1^{(k)})^2}\left(Z(𝐊_\varphi ;𝐐_\varphi ^{(k)},𝐮_\varphi ^{(k)})+q^{\frac{1}{16}}Z(𝐊_\varphi ;𝐐_\varphi ^{(k)},𝐮_\varphi ^{(k)}\left(\begin{array}{c}1/2\\ 0\end{array}\right))\right),\end{array}$$
(C.13)
corresponds to the limit
$$\underset{L\mathrm{}}{lim}\underset{t=0}{\overset{2q1}{}}\left(X_{Lt}+Y_{Lt}\right).$$
(C.14)
We have numerically checked that (C.12) indeed equals the $`(\widehat{𝔰𝔩_2})_{k=2,M}`$ character
$$Z_{\text{tot}}=\frac{1}{(q)_{\mathrm{}}}\underset{n\text{}}{}\left(x^{2n}q^{\frac{1}{2q}n^2}\underset{k1}{}(1+q^{k\frac{1}{2}})+x^{2n+1}q^{\frac{1}{2q}(n+\frac{1}{2})^2+\frac{1}{16}}\underset{k1}{}(1+q^k)\right),$$
(C.15)
corresponding to a free fermion and a boson compactified on a circle of radius $`R^2=q`$. It should be possible to prove the equality of (C.12) and (C.15) along the lines of (see also App. D). Finally, we note that the number of summands in (C.12) equals the torus degeneracy for the $`q`$-Pfaffian computed in .
## Appendix D $`(\widehat{𝔰𝔩_3})_{k,M}`$ character
We will restrict the discussion in this section to $`(\widehat{𝔰𝔩_3})_{k,M}`$ for level $`k=1`$.
### D.1. Quasihole sector
The recursion relation for the quasiholes $`(\varphi _{},\varphi _{})`$ in $`(\widehat{𝔰𝔩_3})_{k,M}`$ for $`k=1`$ and $`M=0`$ was worked out in . The generalization to arbitrary $`M`$ reads
$$X_L=X_{L(2M+3)}+(x_{}+x_{})q^{\frac{2L(M+2)}{2(2M+3)}}X_{L(M+2)}+x_{}x_{}q^{\frac{2L1}{2M+3}}X_{L1}.$$
(D.1)
By putting $`X_L/X_{L1}\lambda _{\text{tot}}^{\frac{1}{2M+3}}`$ we recover the IOW-equation (7.5). To build the entire spectrum out of quasiholes and electrons we need $`3M+4`$ sectors whose initial conditions are given in Table D.1. The vacua of the sectors are labeled by, respectively, charge, spin, and the $`𝔰𝔩_3`$ irrep in which they occur for $`M=0`$. The parameter $`r`$ takes the values $`r=1,2,\mathrm{},M+1`$.
The solution to (D.1) can be written in terms of finitized UCPFs (see (5.10)). Indeed, the recursion relations (5.12) with (see (7.3))
$$𝐊_\varphi =\frac{1}{2M+3}\left(\begin{array}{cc}M+2& (M+1)\\ (M+1)& M+2\end{array}\right),$$
(D.2)
and $`𝐐+𝐮=(0,0)`$ are explicitly given by
$`P_{L_1,L_2}`$ $`=`$ $`P_{L_11,L_2}+x_{}q^{L_1\frac{M+2}{2(2M+3)}}P_{L_1\frac{M+2}{2M+3},L_2+\frac{M+1}{2M+3}}`$
$`P_{L_1,L_2}`$ $`=`$ $`P_{L_1,L_21}+x_{}q^{L_2\frac{M+2}{2(2M+3)}}P_{L_1+\frac{M+1}{2M+3},L_2\frac{M+2}{2M+3}}`$ (D.3)
Setting $`X_LP_{L/(2M+3),L/(2M+3)}`$ leads to (D.1). The values for $`𝐐=𝐮`$ in the various sectors, as determined by the initial conditions, are given in Table D.1.
### D.2. Electron sector
The recursion relations for the electrons $`(\mathrm{\Psi }_{},\mathrm{\Psi }_{})`$ are given by
$$\mathrm{\Omega }_L=\mathrm{\Omega }_{L1}+(y_{}+y_{})q^{L\frac{M}{2}}\mathrm{\Omega }_{L(M+2)}+y_{}y_{}q^{2L(2M+1)}\mathrm{\Omega }_{L(2M+3)}.$$
(D.4)
with initial conditions listed in Table D.2. They can be solved by finitized UCPFs with (see (7.9))
$$𝐊_e=\left(\begin{array}{cc}M+2& M+1\\ M+1& M+2\end{array}\right),$$
(D.5)
and $`𝐐+𝐮=(1,1)`$, by putting $`\mathrm{\Omega }_L=P_{L,L}`$. The values for $`𝐐_e`$ in the various sectors are listed in Table D.2.
### D.3. Character
Combining the $`3M+4`$ sectors in Tables D.1 and D.2 should reproduce the spectrum of the chiral $`(\widehat{𝔰𝔩_3})_{k=1,M}`$ conformal field theory. Indeed, consider the following combination of UCPFs
$$Z_{\text{tot}}=\underset{k=0}{\overset{3M+3}{}}a_{(k)}Z_{\mathrm{}}(𝐊_e;𝐐_e^{(k)})Z_{\mathrm{}}(𝐊_\varphi ;𝐐_\varphi ^{(k)}),$$
(D.6)
where the coefficients $`a_{(k)}`$ are defined in Table D.3.
We claim that (D.6) equals the $`(\widehat{𝔰𝔩_3})_{k=1,M}`$ character
$$Z_{\text{tot}}=\frac{1}{(q)_{\mathrm{}}^2}\underset{p_i\text{}}{}(x_{}^{p_1}x_{}^{p_2})q^{\frac{1}{2}𝐩𝐊_\varphi 𝐩}$$
(D.7)
corresponding to the partition function of two chiral bosons on the deformed weight lattice of $`𝔰𝔩_3`$. E.g., (D.7) for $`M=0`$ is precisely the Frenkel-Kac character (see, e.g., ) of the sum of the integrable highest weight modules of $`(\widehat{𝔰𝔩_3})`$ at level $`k=1`$.
To prove this claim, first observe that we can rewrite (D.6) as a sum over $`2M+3`$ sectors by using (5.8). Specifically,
$$Z_{\mathrm{}}(𝐊_\varphi ,\left(\begin{array}{c}\frac{2M+3r}{2M+3}\\ \frac{2M+3r}{2M+3}\end{array}\right))+x_{}q^{\frac{M+22r}{2(2M+3)}}Z_{\mathrm{}}(𝐊_\varphi ,\left(\begin{array}{c}\frac{M+2r}{2M+3}\\ \frac{M+2r}{2M+3}\end{array}\right))=Z_{\mathrm{}}(𝐊_\varphi ,\left(\begin{array}{c}\frac{2M+3r}{2M+3}\\ \frac{r}{2M+3}\end{array}\right)),$$
(D.8)
after which the claim follows by applying the statements of Theorem 4.1 and Corollary 5.2 in .
Note that even though we use $`3M+4`$ sectors in generating the entire spectrum from the recursion relations (D.1) and (D.4), the $`(\widehat{𝔰𝔩_3})_{k=1,M}`$ partition function (D.7) can be written in terms of UCPFs based on $`𝐊=𝐊_e𝐊_\varphi `$ using only $`2M+3`$ sectors. So, even though the UCPF form of a partition function is not unique, and we do not understand the precise relation between the number of sectors and the torus degeneracy in the sense of Wen et al. , it is satisfying to see that the number $`2M+3`$ equals $`det𝐊_e`$ which is the torus degeneracy for abelian fqH systems.
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# 1 Introduction
## 1 Introduction
Supersymmetric quantum mechanics which underlies the dynamics of non–relativistic and relativistic spinning particles and superparticles is one of the simplest examples of supersymmetric sigma–models and it has attracted a great deal of attention as a laboratory for studying problems appearing in more complicated supersymmetric field and string theories. For instance, one–dimensional and multidimensional $`N=4`$ supersymmetric quantum mechanics (SUSY QM) can be associated with $`N=1`$, $`D=4`$ supersymmetric field theories (including supergravity) subject to an appropriate dimensional reduction down to $`D=1`$.
A recent revival of interest in superconformal mechanics has been caused, in particular, by an observation made in the context of the AdS/CFT correspondence conjecture that the dynamics of a superparticle near the AdS horizon of an extreme Reissner–Nordström black hole of a large mass is described by a superconformal mechanics . Applications of supersymmetric mechanics to the theory of black holes and to other problems have been reviewed in . In conditions on geometry of curved backgrounds, in which $`N=1,2`$ and 4 superconformal invariant models of non-relativistic spinning particles can exist, have been studied in the $`N=1`$ superfield formalism.
Note that, as in the case of superstrings, the superconformal group of the particle superworldline is an infinite dimensional subgroup of the group of its superdiffeomorphisms. It becomes manifest in the worldline superfield formulations of relativistic spinning particles and superparticles , which can thus be regarded as examples of quantum mechanics with manifest superconformal symmetry.
The superconformal invariance and, in more general case, worldvolume superdiffeomorphisms impose restrictions on the geometry of the background also in models of relativistic particles and branes. For instance, in the case of superbranes it requires that a target–superspace background obeys superfield supergravity constraints (see for a recent review).
In the case of spinning particles this problem is connected with the problem of selfconsistent field theoretical description of interacting particles with spin higher than $`2`$. It is well known that the theory of interacting higher spin fields should be formulated in an anti–de–Sitter background (see for a review).
In it was shown that difficulties with constructing a model of a spinning particle moving in a gravitational background arise already for spins $`3/2`$ and $`2`$. These difficulties have been overcome in , where an action for spinning particles with spin higher than one were constructed in backgrounds of constant curvature (such as the AdS spaces)<sup>1</sup><sup>1</sup>1We thank Sergey Kuzenko for bringing these papers to our attention..
The spin 2 particle model of is based on a so called “large” $`N=4`$ superconformal algebra containing $`SO(4)`$ as the subalgebra of local internal symmetries. It is well known that there exists another (so called “small”) $`N=4`$ superconformal algebra with $`SU(2)`$ as the subalgebra of local internal symmetries. It is therefore tempting to study whether in a superfield formulation of a spin 2 particle dynamics, which is manifestly invariant under “small” $`N=4`$ superconformal symmetry, conditions imposed on a curved background can be less restrictive than in the case of .
In this paper we present results of this study. We consider relativistic spinning particle mechanics invariant under local $`N=4`$ supersymmetry with $`SU(2)_{local}\times SU(2)_{global}`$ internal symmetries, which is associated with the “small” $`N=4`$ superconformal algebra. In a flat background this model has been constructed and studied in . It was shown that (in four dimensions) its first quantized spectrum consists of three scalar and one spin $`2`$ states corresponding to the linearized limit of a conformal gravity model. The superfield action for this $`N=4`$ spinning particle is a localized (or superconformal) version of the action for $`N=4`$ supersymmetric quantum mechanics with a quadratic superpotential. This correspondence prompts us how to generalize the free $`N=4`$ superconformal spinning particle action to the description of a particle propagating in a gravitational background. For this one should consider supersymmetric quantum mechanics with an arbitrary superpotential and make it invariant under the $`N=4`$ superconformal transformations .
In it has been shown that the $`N=4`$ superfield formulation of multidimensional $`N=4`$ SUSY QM leads to a supersymmetric nonlinear sigma–model with a target–space metric being a second derivative of a single real–valued function (superpotential) $`A(x)`$
$$g_{MN}(x)=\frac{^2A(x)}{x^Mx^N},$$
(1.1)
i.e. for an arbitrary dimension $`D`$ and signature of the sigma–model manifold, parametrized by real scalar fields $`x^M`$ $`(M=0,1,\mathrm{},D1)`$, its metric should have a “Kähler–like” structure. The metric of a similar type appeared also as a metric of black hole moduli spaces considered recently in .
As has been announced in , the $`N=4`$ superconformal generalization of the model of in a manifold of Minkowski signature describes a relativistic spinning particle propagating in a gravitational background with the metric (1.1).
It has been known for a long time that supersymmetry requires sigma–model manifolds of chiral superfields to be Kähler, hyper–Kähler , special Kähler , or special Lagrangian manifolds . The geometrical structure of these manifolds has been under intensive investigation because of its relation to the compactification of string theory on Calabi–Yau manifolds and to duality symmetries of corresponding supergravity models (see for a review).
The essential difference of the metric (1.1) from a Kähler metric
$$g_{MN}(z,\overline{z})=\frac{^2K(z,\overline{z})}{z^M\overline{z}^N}$$
(1.2)
is that the latter is a Hermitian metric on a complex manifold, while the former is a real manifold metric. The reason why a real sigma–model manifold appears in the case of $`N=4`$ SUSY QM under consideration is that we construct the supersymmetric sigma model with the use of constrained real superfields and not with chiral ones as one usually do.
Some Kähler manifolds mentioned above also admit real-valued representation for the metric (1.1). For example, this is so for a metric of the special Kähler manifolds in a flat Darboux coordinate system . However, the class of the manifolds with the metric (1.1) is more general and includes manifolds which do not have complex structure.
In particular, we have found that in a certain coordinate system the metric on an anti–de–Sitter space of an arbitrary dimension $`D`$ ($`AdS_D`$) can be represented in the form (1.1). Other examples are hyperbolic manifolds of negative curvature on which M–theory and string theories can compactify . To the best of our knowledge this observation is a novel one. This result can presumably be useful for better understanding the structure of string and supergravity theories in AdS superbackgrounds and AdS/CFT correspondence.
The paper is organized as follows. In Section 2 we review the $`N=4`$ superconformal particle model of ref. . In Section 3 we generalize it to describe a spinning particle in a curved background with the metric (1.1). Properties of the $`AdS_D`$ space which follow from the potential structure of its metric are considered in Section 4. In Conclusion we discuss open problems and outlook.
## 2 The free $`N=4`$ superconformal particle model
We begin with a brief description of free spinning particle mechanics with $`SU(2)_{local}\times SU(2)_{global}`$ $`N=4`$ superconformal symmetry on the superworldline . To avoid confusion we should note that in one– and two–dimensional spaces the (super)conformal symmetry is infinite dimensional (i.e. the parameter of (super)conformal transformations is a holomorphic function of (super)worldsheet coordinates). The $`N=4`$ superconformal superalgebra with the local internal $`SU(2)`$ automorphisms contains four supercharges and is a subalgebra of a more general $`N=4`$ superconformal algebra with an internal local $`SO(4)`$ which contains eight supercharges . Thus, in our case $`N=4`$ counts all supercharges, while usually (in particular in higher dimensions) it corresponds only to super–Poincare charges and does not include the number of special superconformal generators.
To construct the superfield action in the worldline superspace $`(\tau ,\theta ^a,\overline{\theta }_a)`$ (with $`\tau `$ being a time parameter, and $`\theta ^a`$ and $`\overline{\theta }_a=(\theta ^a)^{},(a=1,2)`$ being two complex (or four real) Grassmann–odd coordinates) one introduces $`D`$ real “matter” superfields $`\mathrm{\Phi }^M(\tau ,\theta ^a,\overline{\theta }_a)`$ ($`M=0,1,\mathrm{},D1)`$ and a worldline supereinbein $`E(\tau ,\theta ^a,\overline{\theta }_a)`$ which have the following properties with respect to the $`SU(2)`$ $`N=4`$ superconformal transformations of the worldline superspace <sup>2</sup><sup>2</sup>2Our conventions for spinors are as follows: $`\theta _a=\theta ^b\epsilon _{ba},\theta ^a=\epsilon ^{ab}\theta _b,\overline{\theta }_a=\overline{\theta }^b\epsilon _{ba},\overline{\theta }^a=\epsilon ^{ab}\overline{\theta }_b,\overline{\theta }_a=(\theta ^a)^{},\overline{\theta }^a=(\theta _a)^{},(\theta \theta )\theta ^a\theta _a=2\theta ^1\theta ^2,(\overline{\theta }\overline{\theta })\overline{\theta }_a\overline{\theta }^a=(\theta \theta )^{},(\overline{\theta }\theta )\overline{\theta }_a\theta ^a,\epsilon ^{12}=\epsilon ^{21}=1,\epsilon _{12}=1`$.
$`\delta \tau `$ $`=`$ $`\mathrm{\Lambda }{\displaystyle \frac{1}{2}}\theta ^aD_a\mathrm{\Lambda }{\displaystyle \frac{1}{2}}\overline{\theta }_a\overline{D}^a\mathrm{\Lambda },`$
$`\delta \theta ^a`$ $`=`$ $`i\overline{D}^a\mathrm{\Lambda },\delta \overline{\theta }_a=iD_a\mathrm{\Lambda },`$ (2.1)
$`\delta \mathrm{\Phi }^M`$ $`=`$ $`\mathrm{\Lambda }\dot{\mathrm{\Phi }}^M+\dot{\mathrm{\Lambda }}\mathrm{\Phi }^Mi(D_a\mathrm{\Lambda })(\overline{D}^a\mathrm{\Phi }^M)i(\overline{D}^a\mathrm{\Lambda })(D_a\mathrm{\Phi }^M),`$ (2.2)
$`\delta E`$ $`=`$ $`\mathrm{\Lambda }\dot{E}\dot{\mathrm{\Lambda }}Ei(D_a\mathrm{\Lambda })(\overline{D}^aE)i(\overline{D}^a\mathrm{\Lambda })(D_aE),`$ (2.3)
where dot denotes the time derivative $`\frac{d}{d\tau }`$. The transformation law (2.3) for the superfields $`\mathrm{\Phi }^M`$ shows that these superfields are vector superfields in the one dimensional $`N=4`$ superspace, while the superfields $`E\mathrm{\Phi }^M`$ are scalars.
The superfields $`\mathrm{\Phi }^M`$ and $`E`$ obey the quadratic constraints
$`[D_a,\overline{D}^a]\mathrm{\Phi }^M`$ $`=`$ $`0,`$
$`D^aD_a\mathrm{\Phi }^M`$ $`=`$ $`0,`$
$`\overline{D}_a\overline{D}^a\mathrm{\Phi }^M`$ $`=`$ $`0,`$ (2.4)
and
$`[D_a,\overline{D}^a]{\displaystyle \frac{1}{E}}`$ $`=`$ $`0,`$
$`D^aD_a{\displaystyle \frac{1}{E}}`$ $`=`$ $`0,`$
$`\overline{D}_a\overline{D}^a{\displaystyle \frac{1}{E}}`$ $`=`$ $`0,`$ (2.5)
where
$$D_a=\frac{}{\theta ^a}\frac{i}{2}\overline{\theta }_a\frac{}{\tau },\overline{D}^a=\frac{}{\overline{\theta }_a}\frac{i}{2}\theta ^a\frac{}{\tau },$$
(2.6)
are the supercovariant derivatives, and the infinitesimal superfield
$`\mathrm{\Lambda }(\tau ,\theta ,\overline{\theta })`$ $`=`$ $`a(\tau )+\theta ^a\overline{\alpha }_a(\tau )\overline{\theta }_a\alpha ^a(\tau )+\theta ^a(\sigma _i)_a^b\overline{\theta }_bb_i(\tau )`$ (2.7)
$`{\displaystyle \frac{i}{2}}(\theta ^a\dot{\overline{\alpha }}_a(\tau )+\overline{\theta }_a\dot{\alpha }^a(\tau ))\overline{\theta }\theta +{\displaystyle \frac{1}{8}}(\overline{\theta }\theta )^2\ddot{a}(\tau )`$
contains the parameters of local reparametrizations $`a(\tau )`$, local supertranslations $`\alpha ^a(\tau )`$, $`\overline{\alpha }_a(\tau )`$ and local $`SU(2)`$ rotations $`b_i(\tau )`$ of the worldline superspace. It is constrained by the same relations (2.4) as $`1/E`$ and $`\mathrm{\Phi }^M`$ (($`\sigma _i)^b_a`$ are the Pauli matrices, $`i=1,2,3`$).
The constraints (2.4)–(2.5) can be explicitly solved , the solution being described by the superfields
$`\mathrm{\Phi }^M(\tau ,\eta ,\overline{\eta })`$ $`=`$ $`{\displaystyle \frac{1}{e(\tau )}}x^M(\tau )+\theta ^a\overline{\psi }_a^M(\tau )\overline{\theta }_a\psi ^{Ma}(\tau )+\theta ^a(\sigma _i)_a^b\overline{\theta }_bT_i^M(\tau )`$ (2.8)
$`{\displaystyle \frac{i}{2}}(\theta ^a\dot{\overline{\psi }}_a^M(\tau )+\overline{\theta }_a\dot{\psi }^{Ma}(\tau ))\overline{\theta }\theta +{\displaystyle \frac{1}{8}}(\overline{\theta }\theta )^2{\displaystyle \frac{d^2}{d\tau ^2}}\left({\displaystyle \frac{1}{e(\tau )}}x^M\right),`$
and
$`{\displaystyle \frac{1}{E}}(\tau ,\theta ,\overline{\theta })`$ $`=`$ $`{\displaystyle \frac{1}{e(\tau )}}+\theta ^a\overline{\lambda }_a^{}(\tau )\overline{\theta }_a\lambda ^a(\tau )+\theta ^a(\sigma _i)_a^b\overline{\theta }_bt_i^{}(\tau )`$ (2.9)
$`{\displaystyle \frac{i}{2}}(\theta ^a\dot{\overline{\lambda ^{}}}_a(\tau )+\overline{\theta }_a\dot{\lambda }^a(\tau ))\overline{\theta }\theta +{\displaystyle \frac{1}{8}}(\overline{\theta }\theta )^2{\displaystyle \frac{d^2}{d\tau ^2}}{\displaystyle \frac{1}{e(\tau )}}.`$
The leading components $`x^M(\tau )`$ of the superfields $`\mathrm{\Phi }^M`$ are associated with coordinates of the particle trajectory in a $`D`$–dimensional flat target space–time, the Grassmann–odd vectors $`\psi ^{Ma}(\tau )`$ and $`\overline{\psi ^{}}_a^M(\tau )`$ correspond to particle spin degrees of freedom and $`T_i^M(\tau )`$ are auxiliary fields. The superfield $`1/E`$ describes an $`N=4`$ worldline supergravity multiplet consisting of the einbein (“graviton”) $`e(\tau )`$, two complex “gravitini” $`\lambda ^a(\tau )`$ and $`\overline{\lambda }_a^{}(\tau )`$, and the $`SU(2)`$ gauge field $`t_i^{}(\tau )`$. Upon an appropriate field redefinition (see eqs. (3.2) of the next section) we shall pass from “primed” to “unprimed” component fields.
The $`N=4`$ superfield action for a relativistic spinning particle in a flat target space has the following form
$$S=8𝑑\tau d^2\theta d^2\overline{\theta }E\mathrm{\Phi }^M\mathrm{\Phi }^N\eta _{MN}$$
(2.10)
where $`\eta _{MN}=diag(,+,\mathrm{},+)`$ is the Minkowski metric. The components of $`E`$ play the role of Lagrange multipliers. Their presence implies that the dynamics of the particle is subject to relativistic constraints, in particular, the particle is massless $`(p_Mp^M=0)`$ . The Dirac quantization of the model (2.10) shows that its quantum spectrum consists of one spin $`2`$ and three spin $`0`$ particle states and it can be regarded as a linearized spectrum of a conformal gravity .
## 3 The spinning particle in a curved background
Let us now generalize the model of the previous section to describe a spinning particle propagating in the gravitational background. To this end we replace (2.10) with the most general action functional which respects the $`N=4`$ superconformal symmetry
$$S=8𝑑\tau d^2\theta d^2\overline{\theta }E^1A(E\mathrm{\Phi }^M),$$
(3.1)
where $`A(E\mathrm{\Phi }^M)`$ is an arbitrary function (called the superpotential) of $`E\mathrm{\Phi }^M`$. Recall that $`E\mathrm{\Phi }^M`$ transform as scalar superfields with respect to (2.3), while $`\mathrm{\Phi }^M`$ and $`\frac{1}{E}`$ are vectors. Note also that $`E^1A(E\mathrm{\Phi }^M)`$ can be regarded as a rank one homogeneous function in a $`D+1`$ dimensional space with $`x^D=E^1`$. A consequence of such a structure of the superfield action (3.1) is the fact that only $`D`$ of the bosonic coordinates in the $`D+1`$ dimensional space describe dynamical degrees of freedom. The einbein $`e(\tau )`$ and its superpartners are auxiliary fields as in the free particle case (2.10).
Integrating (3.1) over the Grassmann coordinates $`\theta ^a`$, $`\overline{\theta }_a`$ and making the following redefinition of the component fields
$$\lambda ^a=e^{\frac{3}{2}}\lambda ^a,\overline{\lambda }_a=(\lambda ^a)^{},t_i=2e(t_{}^{}{}_{i}{}^{}+e\lambda ^b(\sigma _i)_b^a\overline{\lambda }_a^{}),\psi ^{Ma}=\sqrt{e}(\psi ^{Ma}x^M\lambda ^a),$$
$$\overline{\psi }_a^M=(\psi ^{Ma})^{},T_i^M=2\sqrt{e}(T_i^Mx^Mt_i^{}+\frac{\sqrt{e}}{2}\lambda ^b(\sigma _i)_b^a\overline{\psi }_a^M+\frac{\sqrt{e}}{2}\psi ^{Mb}(\sigma _i)_b^a\overline{\lambda }_a^{}),$$
(3.2)
one obtains the component action
$$S=𝑑\tau (KV),$$
(3.3)
where
$`K`$ $`=`$ $`{\displaystyle \frac{1}{2e}}g_{MN}(\dot{x}^Mi\overline{\lambda }_a\psi ^{Ma}+i\overline{\psi }_a^M\lambda ^a)(\dot{x}^Ni\overline{\lambda }_b\psi ^{Nb}+i\overline{\psi }_b^N\lambda ^b)`$ (3.4)
$`+ig_{MN}(\overline{\psi }_a^M\dot{\psi }^{aN}+\psi ^{aM}\dot{\overline{\psi }_a^N})`$
is the kinetic term and
$`V`$ $`=`$ $`{\displaystyle \frac{1}{2}}g_{MN}T_i^MT_i^N+2\sqrt{e}\mathrm{\Gamma }_{LMN}\psi ^{Mb}(\sigma _i)_b^a\overline{\psi }_a^LT_i^Nt_ig_{MN}\psi ^{Nb}(\sigma _i)_b^a\overline{\psi }_a^M`$
$`+2\mathrm{\Gamma }_{LMN}(\lambda ^a\overline{\psi }_a^L\overline{\psi }_b^M\psi ^{bN}+\overline{\lambda }_b\psi ^{Lb}\psi ^{Ma}\overline{\psi }_a^N)+e(_L\mathrm{\Gamma }_{MNP})(\overline{\psi }_a^L\overline{\psi }^{Ma})(\psi ^{Nb}\psi _b^P)`$
describes fermionic interactions. In eqs. (3.4) and (3)
$$g_{MN}(x)=\frac{^2}{x^Mx^N}A(x)_{MN}^2A(x),A(x^M)=A(E\mathrm{\Phi }^M)|_{\theta ,\overline{\theta }=0}$$
(3.6)
is the metric of a sigma–model $`D`$–dimensional manifold parametrized by the worldline scalar fields $`x^M(\tau )`$ and
$$\mathrm{\Gamma }_{LMN}(x)=\frac{1}{2}_{LMN}^3A(x)$$
(3.7)
is the (totally symmetric) Christoffel connection associated with $`g_{MN}`$ (i.e. $`𝒟_Lg_{MN}=_Lg_{MN}\mathrm{\Gamma }_{LM}^Pg_{PN}\mathrm{\Gamma }_{LN}^Pg_{PM}=0`$).
The Riemann curvature of this manifold has the form
$$R_{LM,NP}=\mathrm{\Gamma }_{LP}^Q\mathrm{\Gamma }_{QMN}\mathrm{\Gamma }_{LN}^Q\mathrm{\Gamma }_{QMP}.$$
(3.8)
Upon solving for the equations of motion of the auxiliary fields $`T_i^M`$, substituting the solution back into eqs. (3.3)–(3) and performing Legendre transformations one arrives at the first order form of the spinning particle action
$$S=𝑑\tau \left[p_M\dot{x}^M+i(\psi _M^a\dot{\overline{\psi }_a^M}+\overline{\psi }_{Ma}\dot{\psi }^{Ma})H\right],$$
(3.9)
where $`p_M`$ is the momentum canonically conjugate to $`x^M`$, and the Hamiltonian $`H`$ of the system has the following structure
$$H=e(\tau )H_0+i\lambda ^a(\tau )\overline{Q}_a+i\overline{\lambda }_a(\tau )Q^at_i(\tau )L_i,$$
(3.10)
with
$`H_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}g^{MN}p_Mp_N+R_{LN,PM}(\overline{\psi }_a^L\overline{\psi }^{Ma})(\psi ^{Nb}\psi _b^P)+R_{MP,NL}(\overline{\psi }_a^L\psi ^{Ma})(\overline{\psi }_b^N\psi ^{Pb})`$ (3.11)
$`+𝒟_LG_{MNP}(\overline{\psi }_a^L\overline{\psi }^{Ma})(\psi ^{Nb}\psi _b^P),`$
$$\overline{Q}_a=\overline{\psi }_a^Mp_M+i\mathrm{\Gamma }_{LMN}\overline{\psi }_c^L\overline{\psi }^{Mc}\psi _a^N,$$
(3.12)
$$Q^b=\psi ^{Ib}p_I+i\mathrm{\Gamma }_{LMN}\overline{\psi }^{Lb}\psi ^{Mc}\psi _c^N$$
(3.13)
and
$$L_i=g_{MN}\psi ^{Nb}(\sigma _i)_b^a\overline{\psi }_a^M$$
(3.14)
being associated with constraints on the dynamics of the relativistic particle caused by the worldline superreparametrization invariance of the model. The constraints are of the first class since they form a closed $`N=4`$ supersymmetry algebra
$$\{\overline{Q}_a,Q^b\}=i\delta _a^bH_0,[L_i,L_j]=ϵ_{ijk}L_k,$$
$$[L_i,\overline{Q}_a]=\frac{i}{2}(\sigma _i)_a^c\overline{Q}_c,[L_i,Q^a]=\frac{i}{2}(\sigma _i)_c^aQ^c$$
(3.15)
with respect to the following graded Dirac brackets (which are obtained upon solving for the second class constraints on the canonical fermionic momenta $`\pi _{Ma}=i\overline{\psi }_{Ma}`$ and $`\overline{\pi }_M^a=i\psi _M^a`$ derived from eq. (3.9))
$$[x^M,p_N]=\delta _N^M,\{\psi ^{aM},\overline{\psi }_b^N\}=\frac{i}{2}\delta _b^ag^{MN},[p_M,p_N]=2iR_{MN,PL}\overline{\psi }_a^P\psi ^{aL},$$
$$[p_M,\psi _N^a]=\mathrm{\Gamma }_{MNP}\psi ^{aP},[p_M,\overline{\psi }_N^a]=\mathrm{\Gamma }_{MNP}\overline{\psi }^{aP},$$
(3.16)
We observe that $`p_M`$ have properties of covariant momenta when acting on fermionic variables $`\psi ^{Ma}`$ and $`\overline{\psi }_a^M`$.
The superalgebra (3.15) of the constraints (3.11)–(3.14) generates the $`SU(2)_{local}\times SU(2)_{global}`$ $`N=4`$ superconformal transformations (2.3) of the components of the superfields $`\mathrm{\Phi }^M`$.
We have thus shown that the $`N=4`$ worldline superfield action (3.1), which reduces to (3.9)–(3.14) upon integrating over Grassmann–odd coordinates and eliminating auxiliary fields, describes the dynamics of an $`N=4`$ superconformal spinning particle in a curved background whose geometry is characterized by eqs. (3.7)–(3.8).
We should note that the last terms in (3.11)–(3.13), containing the Christoffel connection, are non-covariant with respect to general coordinate transformations of the background. The reason is that background diffeomorphisms acting on the superfields $`\mathrm{\Phi }^M`$, in general, are incompatible with the constraints (2.4)–(2.5). This, in particular, means that if a background metric (3.6) admits isometries, not all of them will be symmetries of the actions (3.1) and (3.9). It is an interesting open problem to study whether the model under consideration can be modified in such a way that only target–space covariant terms remain in the action.
## 4 The potential structure of the anti–de–Sitter metric
It is curiously enough that the anti–de–Sitter spaces belong to the class of the manifolds whose metric in a certain coordinate system acquires the form (1.1). To show this consider first a coordinate system
$$X^M=(X^\mu ,\rho ),\mu =0,\mathrm{},D2$$
(4.1)
in which the metric of a $`D`$–dimensional AdS space has a conformally flat form (for simplicity we put the AdS radius to one)
$$ds^2=\frac{1}{\rho ^2}\left(\eta _{\mu \nu }dX{}_{}{}^{\mu }dX^\nu +d\rho ^2\right),$$
(4.2)
where $`\eta _{\mu \nu }=(1,1,\mathrm{},1)`$.
Now perform a coordinate transformation to the new set of variables
$$x^M=(x^\mu ,r)$$
(4.3)
such that
$$X^\mu =\frac{x^\mu }{r},\rho =\frac{1}{\sqrt{r}}.$$
(4.4)
The passage from $`\rho `$ to $`r`$ has proved to be convenient for the analysis of the properties of the potential $`A(x)`$ considered below.
In the coordinate system (4.3) the AdS metric $`g_{MN}`$ takes the form
$$g_{\mu \nu }=\frac{\eta _{\mu \nu }}{r},g_{\mu r}=\frac{\eta _{\mu \nu }x^\nu }{r^2},g_{rr}=\frac{x^\mu x^\nu \eta _{\mu \nu }}{r^3}+\frac{1}{4r^2},$$
(4.5)
where the index $`r`$ of the metric tensor components corresponds to the coordinate $`r`$.
One can easily check that the metric (4.5) is a second derivative of the following function
$$A(x)=\frac{x^\mu x^\nu \eta _{\mu \nu }}{2r}\frac{1}{4}\mathrm{ln}r.$$
(4.6)
Thus we have shown that $`AdS_D`$ is one of the manifolds of the type (1.1), where the $`N=4`$ superconformal spinning particle can live.
By passing note that if in the action (3.1) we take $`A(E\mathrm{\Phi }^M)`$ in the form (4.6) and put $`\mathrm{\Phi }^\mu =0`$ and $`\mathrm{\Phi }^r=1`$ we shall arrive at the action
$$S=2𝑑\tau d^2\theta d^2\overline{\theta }\frac{\mathrm{ln}E}{E}$$
which describes a one–dimensional $`N=4`$ superconformal mechanics considered in .
The potential (4.6) generating the metric on $`AdS_D`$ is not unique. Another form of the potential arises when one performs the following change of variables (4.1)
$$X^\mu =\left(\frac{x^\mu }{r}\right)^{m_\mu }\rho =\frac{1}{\sqrt{r}},$$
(4.7)
where $`m_\mu 0,\frac{1}{2}`$ is a set of real numbers, namely
$$A=\frac{m_0^2}{2m_0(2m_01)}\frac{(x^0)^{2m_0}}{r^{2m_01}}+\underset{i=1}{\overset{D2}{}}\frac{m_i^2}{2m_i(2m_i1)}\frac{(x^i)^{2m_i}}{r^{2m_i1}}\frac{1}{4}\mathrm{ln}r.$$
(4.8)
More generally, we could make, for instance, a “logarithmic transformation”
$$X^\mu =\mathrm{ln}(\frac{x^\mu }{r})\rho =\frac{1}{\sqrt{r}},$$
(4.9)
for which the corresponding potential has the form
$$A=r\mathrm{ln}\frac{x^0}{r}\underset{i=1}{\overset{D2}{}}r\mathrm{ln}\frac{x^i}{r}\frac{1}{4}\mathrm{ln}r.$$
(4.10)
The coordinate transformation (4.4) is singled out by the requirement that it is a single–valued and that a Lorentz subgroup $`SO(1,D2)`$ of the $`AdS_D`$ isometry group $`SO(2,D1)`$ acts linearly on both the “old” coordinates $`X^\mu `$ of (4.1) and the “new” coordinates $`x^\mu `$ of (4.3), (4.4).
So we shall further discuss some amusing properties of $`AdS_D`$ associated with its potential structure (4.6) in the coordinate system (4.3).
The group $`SO(2,D1)`$ of the isometry transformations of $`AdS`$ coordinates, which leave the form of the $`AdS`$ metric invariant, is known to act as a conformal group on a $`(D1)`$-dimensional boundary of $`AdS_D`$. In the coordinate system (4.1) the boundary (which is a $`(D1)`$–dimensional Minkowski space) is associated with the coordinates $`X^\mu `$. Under infinitesimal $`SO(2,D1)`$ transformations $`X^\mu `$ and $`\rho `$ vary as follows
$`\delta X^\mu `$ $`=`$ $`a^\mu +a^{\mu \nu }X^\lambda \eta _{\nu \lambda }+a_DX^\mu +a_K^\mu X^\nu X^\lambda \eta _{\nu \lambda }2(a_K^\nu X^\lambda \eta _{\nu \lambda })X^\mu +a_K^\mu \rho ^2,`$
$`\delta \rho `$ $`=`$ $`(2a_K^\mu X^\nu \eta _{\mu \nu }a_D)\rho ,`$ (4.11)
where the $`SO(2,D1)`$ parameters $`a^\mu ,a^{\mu \nu },a_D,a_K^\mu `$ are, respectively, the parameters of $`D1`$ translations, $`SO(1,D2)`$ rotations, dilatation and conformal boosts, acting as conformal transformations in a $`(D1)`$–dimensional slice of $`AdS_D`$ parametrized by $`X^\mu `$.
From (4.4) and (4) one gets the infinitesimal $`SO(2,D1)`$ transformations of $`x^\mu `$ and $`r`$ (4.3)
$`\delta x^\mu `$ $`=`$ $`a^\mu r+a^{\mu \nu }x^\lambda \eta _{\nu \lambda }a_Dx^\mu +a_K^\mu {\displaystyle \frac{x^\nu x^\lambda \eta _{\nu \lambda }}{r}}+2(a_K^\nu x^\lambda \eta _{\nu \lambda }){\displaystyle \frac{x^\mu }{r}}+a_K^\mu ,`$
$`\delta r`$ $`=`$ $`4a_K^\mu x^\nu \eta _{\mu \nu }2a_Dr.`$ (4.12)
Under (4) the potential (4.6) varies as follows
$$A(x^{})=A(x)+\delta A(x),$$
$$\delta A(x)=\delta x^M_MA(x)=a^\mu x^\nu \eta _{\mu \nu }+\frac{1}{2}a_D+a_K^\mu x^\nu \eta _{\mu \nu }\frac{x^\lambda x^\sigma \eta _{\lambda \sigma }}{r^2}.$$
(4.13)
One can check that the form of the metric (4.4) remains invariant under the action of the $`SO(2,D1)`$ transformations (4), so that they are indeed the isometries of this $`AdS`$ metric. However, the superfield action (3.1) is invariant only under the subgroup of $`SO(2,D1)`$ generated by $`D1`$ translations $`a^\mu `$, $`SO(1,D2)`$ Lorents rotations $`a^{\mu \nu }`$ and dilatations $`a_D`$ which transform the superfields $`E\mathrm{\Phi }^M`$ in the same way as $`x^M`$ in (4). As can be seen from the form of the variation of $`x^M`$ (and respectively of $`E\mathrm{\Phi }^M`$) with respect to conformal boosts $`a_K^\mu `$, the corresponding term does not satisfy the superfield constraints (2.4) and (2.5), and, hence, the transformed $`\mathrm{\Phi }^M`$ will not do so as well. This is the reason of the appearance of noncovariant terms depending on the Christoffel connection in the component actions (3.3) and (3.9)–(3.14).
An interesting property of the potential (4.6) is that the contraction of its partial derivatives with the coordinates (4.3) are constants starting from the second derivative
$$x^Mx^N_{MN}^2A(x)=x^Mx^Ng_{MN}=\frac{1}{4}$$
$$x^Lx^Mx^N_{LMN}^3A(x)=2x^Lx^Mx^N\mathrm{\Gamma }_{LMN}=\frac{1}{2}$$
$$x^{M_1}\mathrm{}x^{M_{n+1}}_{M_1\mathrm{}M_{n+1}}^{n+1}A(x)=(1)^{n+1}\frac{n+1}{4},n=1,\mathrm{},\mathrm{}.$$
(4.14)
To get the relation (4.14) one should note that under the following rescaling of the coordinates (4.3) $`x^M(1+ϵ)x^M`$ (where $`ϵ`$ is a numerical parameter)<sup>3</sup><sup>3</sup>3One should not confuse this rescaling with dilatation isometry (4) which acts as follows $`x^\mu (1+ϵ)x^\mu `$ and $`r(1+ϵ)^2r`$. the potential (4.6) takes the form
$`A_ϵ`$ $`=`$ $`(1+ϵ){\displaystyle \frac{x^\mu x^\nu \eta _{\mu \nu }}{2r}}{\displaystyle \frac{1}{4}}\mathrm{ln}r{\displaystyle \frac{1}{4}}\mathrm{ln}(1+ϵ)`$ (4.15)
$`=`$ $`(1+ϵ){\displaystyle \frac{x^\mu x^\nu \eta _{\mu \nu }}{2r}}{\displaystyle \frac{1}{4}}\mathrm{ln}r{\displaystyle \frac{1}{4}}ϵ+{\displaystyle \frac{1}{4}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^{n+1}{\displaystyle \frac{n+1}{(n+1)!}}ϵ^{n+1},`$
where on the right hand side of (4.15) we have expanded $`\mathrm{ln}(1+ϵ)`$ in series of $`ϵ`$.
On the other hand
$$A_ϵ=A(x+ϵx)=A(x)+ϵx^M_MA(x)+\underset{n=1}{\overset{\mathrm{}}{}}\frac{ϵ^{n+1}}{(n+1)!}x^{M_1}\mathrm{}x^{M_{n+1}}_{M_1\mathrm{}M_{n+1}}^{n+1}A(x).$$
(4.16)
Comparing (4.15) with (4.16) we get (4.14).
A local basis in a tangent space of the $`AdS_D`$ manifold can be described by the following vielbein one–form $`e^A=dx^Me_M^A(x)`$ ($`A=0,1\mathrm{},D1`$)
$$e^\alpha =dx^\mu \delta _\mu ^\alpha r^{\frac{1}{2}},e^r=dx^\mu x_\mu r^{\frac{3}{2}}+dr\frac{1}{2r}$$
(4.17)
determined such that $`g_{MN}=e_M^Ae_N^B\eta _{AB}`$ and $`\eta _{AB}=(,+\mathrm{},+)`$.
Using (4.17) it is easy to calculate the determinant of the metric (4.4)
$$detg_{MN}=(dete_M^A)^2=\frac{1}{4r^{D+1}}.$$
(4.18)
One more observation concerns the form of the covariant derivative of the $`AdS_D`$ Christoffel connection (3.7) appeared in (3.11). A direct computation results in the following relation
$$𝒟_L\mathrm{\Gamma }_{MNP}=\frac{1}{8}_{LMNP}^4A(x)g_{LM}g_{NP}g_{MN}g_{LP}g_{NL}g_{MP}.$$
(4.19)
We see that noncovariance of (4.19) is in a certain sense ”concentrated” in the fourth partial derivative of $`A(x)`$.
Note that the results of this section do not depend on the signature of the metric $`\eta _{\mu \nu }`$ in (4.2). For instance, we could equally well choose $`\eta _{\mu \nu }`$ to be Euclidean. Then we would deal with Euclidean AdS, or hyperbolic spaces considered recently in the context of string and M–theory compactifications .
## 5 Discussion
To conclude, in this paper we have considered the classical dynamics of a spinning particle governed by the action invariant under the $`SU(2)_{local}\times SU(2)_{global}`$ $`N=4`$ superconformal transformations of the particle superworldline. We have shown that the $`N=4`$ superconformal invariance allows the particle to propagate in a curved background with a “Kähler–like” metric generated by a real superpotential $`A`$, and we have found that the anti–de–Sitter and hyperbolic spaces belong to this class of manifolds.
There are several directions of the extension of the results of this paper. One of them is the quantum description of the $`N=4`$ superconformal particle model, which can be carried out following either the lines of or using path integral quantization methods. The latter procedure seems to be more attractive, since it may lead to deeper understanding of the model, for instance, in the context of the AdS/CFT correspondence conjecture.
In particular, it is interesting to study both the classical and quantum dynamics of the $`N=4`$ superconformal spinning particle moving in backgrounds which are direct products of $`AdS_D`$ and Kähler manifolds. Particle motion on the Kähler manifolds can be described by making a multidimensional generalization of the $`N=4`$ supersymmetric quantum mechanics considered in . For this, in addition to $`\mathrm{\Phi }^M`$, one should introduce a number of chiral superfields $`\mathrm{\Psi }^n(\tau ,\theta ,\overline{\theta })`$ ($`\overline{D}^a\mathrm{\Psi }^n=0`$)
$`\mathrm{\Psi }^n(\tau ,\theta ,\overline{\theta })`$ $`=`$ $`z^n(\tau )+\theta ^a\chi _a^n(\tau )+{\displaystyle \frac{i}{2}}\overline{\theta }\theta \dot{z}^n(\tau )+\theta \theta F^n(\tau )`$ (5.1)
$`{\displaystyle \frac{i}{4}}\theta \theta \overline{\theta }_a\dot{\chi }^{na}(\tau ){\displaystyle \frac{1}{16}}\overline{\theta }\overline{\theta }\theta \theta \ddot{z}^n(\tau ),`$
and their complex conjugate antichiral superfields $`\overline{\mathrm{\Psi }}^n(\tau ,\theta ,\overline{\theta })`$. The superfields $`\mathrm{\Psi }`$ and $`\overline{\mathrm{\Psi }}`$ transform as scalars under the $`N=4`$ superconformal transformations (2.1).
We can add to the action (3.1) the following $`N=4`$ superconformal invariant action constructed from $`\mathrm{\Psi }^n`$ and $`\overline{\mathrm{\Psi }}^n`$
$$S_K=2𝑑\tau d^2\theta d^2\overline{\theta }\frac{1}{E}K(\mathrm{\Psi },\overline{\mathrm{\Psi }}),$$
(5.2)
where $`K`$ is a Kähler superpotential.
When the superpotential $`A(E\mathrm{\Phi }^M)`$ is chosen in the form (4.6), the sum of the actions (3.1) and (5.2) describes a spinning particle propagating in an $`AdS_D\times K_{2n}`$ background, where $`K_{2n}`$ is a Kähler manifold with a metric (1.2). For instance, the case $`n=1`$ and $`K_2=\mathrm{ln}(1+\mathrm{\Psi }\overline{\mathrm{\Psi }})`$ corresponds to a two-dimensional sphere $`S^2`$, which is known to be a Kähler manifold. A detailed analysis of these models will be given elsewhere.
Acknowledgments. We are grateful to Dmitri Fursaev, Paul Howe, Armen Nersessian, George Papadopoulos, Paolo Pasti, Volodya Rubtsov, Mario Tonin and Peter West for interest to this work and helpful discussions. Work of A.P. and M.T. was supported in part by the Russian Foundation of Fundamental Research, under the grant 99-02-18417 and the joint grant RFFR-DFG 99-02-04022, and by a grant of the Committee for Collaboration between Czech Republic and JINR. Work of V.O.R. was partially supported by CNPq and a grant by FAPESP. Work of D.S. was partially supported by the European Commission TMR Programme ERBFMPX-CT96-0045 to which the author is associated. M.T. is grateful to the Abdus Salam International Centre for Theoretical Physics, Trieste, where a part of this work was done.
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# 1 Introduction
## 1 Introduction
Every well determined feature of the cosmic ray energy spectrum will have considerable impact on theories of the origin, acceleration and propagation of cosmic rays. After 40 years effort by many groups, the details of the energy spectrum above 10<sup>17</sup> eV are still limited by statistics, systematics and resolution. Experimentalists have been searching for the existence of a cut off on the energy spectrum above 10<sup>20</sup> eV for more than 30 years. This cutoff could result from the interaction of cosmic ray protons or nuclei and the 2.7<sup>o</sup> K black body radiation if the sources are distant enough.
The detection of these extremely high energy cosmic rays is necessarily indirect because of the extremely low flux. The earth’s atmosphere makes the low flux detectable by converting the cosmic ray primaries into extensive air showers(EAS) of various secondary particles at ground level by detecting the secondary particles or the Cerenkov light they produce. Alternatively, one can detect the atmospheric nitrogen fluorescence light induced by the passage of the shower. This technique employed by the Fly’s Eye detector and its successor the High Resolution Fly’s Eye (HiRes) are the only way capable of measuring longitudinal shower developments individually, thus allowing a direct estimation of each shower’s primary energy.
The basic mechanism contribute to the generation of the light signal seen by Fly’s Eye or HiRes detector is nitrogen fluorescence light which relates directly to the number of charged particles in an EAS. The fluorescence light is emitted isotropically from the shower, allowing for detection of showers at large distances. In HiRes detector, UHE shower can bee seen as far as 40 km’s away. This makes a very big aperture, and has high statistics for the cosmic ray energy spectrum. The disadvantage is that, the multiply scattering effect will become important when the shower is far away. In this paper, we will study this effect on the energy estimation of this technique.
The multiply scattering of light by molecules and aerosols in the atmosphere gives rise to a radiance field about a pointlike source. This contribution, or called aureole, depends on many parameters, such as the single scattering albedo, the optical thickness, the scattering phase function, the off angle of the detector from the source direction, the field view solid angle of the detector, and also the integration time window of the detector. This effect has been studied for the last 25 years by the people who are interested in solar blind ultraviolet communication and warning systems or atmospheric study. However most studies were devoted mainly to the total aureole effect . There were few studies on the time-dependent aureole effect. But they were either for the case of very large optical thickness (fog, cloud, etc) or only for very low order of scattering which is for the case of small optical thickness .
Since the shower detected by HiRes can be in the range from a few km to more than 40 km, on good weather ( less aerosol) or bad weather ( dense aerosol). And also the detector is integrating the signal within about 5 micro second time window. We will present calculations of the time-dependent aureole radiance field about an impulsive isotropic source in a scattering and absorbing medium. The approach used is generally applicable to any case, and does not involve any approximation. The calculation of the temporal characteristics of the scattering radiation are based on the work of Trakhovsky et al. .
In Section 2 we will present the recursive approach to calculate the scattering radiation effect initially developed by Zachor , and generalized by Trakhovsky et al to time-dependent case. In Section 3 we will present the Monte Carlo way to calculate the temporal scattering radiation. The results are given in Section 4 for both cases in different atmosphere conditions and detector setup.
And then we apply the result for any cosmic ray shower detected by HiRes detector, and calculate the effect on the energy estimate. This is done in Section 5.
## 2 Recursive Approach
As shown in Fig. 1. The first-order scattering term can be decomposed into three steps: direct transmission from the source to the volume element dV, scattering inside dV in the detector direction, and direct transmission from dV to the detector. Assume an isotropic point source which emits at an instant t = 0 an inpuls of total photo Q. The irradiance (power/unit area) incident on element $`dV=r_2^2dr_2d\mathrm{\Omega }`$ is:
$$=\frac{Q}{4\pi r_1^2}exp[\alpha +\beta )r_1].$$
(1)
in which $`\beta `$ is the volume scattering coefficient, $`\alpha `$ is the absorption coefficient. Accordind to the definition of scattering coefficient $`\beta `$ and single-scattering phase function $`P\left(cos\theta \right)`$, the scattered intensity of radiation by dV toward detector is:
$$d_s=\beta P(cos\theta )dV.$$
(2)
The irradiance incident on the detector is:
$$d_s=\frac{d_s}{r_2^2}exp[(\alpha +\beta )r_2].$$
(3)
Then the radiance (photo/unit area/unit solid angle) received at O is:
$$d𝒩_s=\frac{d_s}{d\mathrm{\Omega }}=\frac{Q\beta P\left(cos\theta \right)}{4\pi r_{1}^{}{}_{}{}^{2}}exp\left[\left(\alpha +\beta \right)\left(r_1+r_2\right)\right]dr_2$$
(4)
Since we are interested in a time-resolved measurement. Following Reilly and Warde . we chose a prolate spheroidal coordinate system with a source positioned at one local point and a detector at the other. It may be shown that the parameters of Fig. 1. ($`r_1,r_2,\theta ,\gamma `$) are transformed into prolate spheroidal coordinates using the following relationships:
$`r_1`$ $`=`$ $`{\displaystyle \frac{R}{2}}\left(\xi +\eta \right),`$ (5)
$`r_2`$ $`=`$ $`{\displaystyle \frac{R}{2}}\left(\xi \eta \right),`$ (6)
$`cos\theta `$ $`=`$ $`{\displaystyle \frac{2\xi ^2\eta ^2}{\xi ^2\eta ^2}},`$ (7)
$`\gamma `$ $`=`$ $`cos^1\left({\displaystyle \frac{1\xi \eta }{\xi \eta }}\right).`$ (8)
let $`ct=r_1+r_2`$. t is the scattered photon time of flight. Using Eqs. (5) and (6) we obtain:
$$\xi =\frac{r_1+r_2}{R}=\frac{ct}{R},$$
(9)
Transforming Eq. (8) we have
$`\eta `$ $`=`$ $`{\displaystyle \frac{1\xi cos\gamma }{\xi cos\gamma }},`$ (10)
Using above Eqs. then Eq. (4) can be mathematically transferred to
$`d𝒩_s(R,\gamma ,t)`$ $`=`$ $`{\displaystyle \frac{cQ\beta P\{cos[\theta (R,\gamma ,t)]\}}{2\pi R^2}}{\displaystyle \frac{exp[(\alpha +\beta )ct]}{\xi ^22\xi cos\gamma +1}}dt.`$ (11)
Where $`\theta `$ is calculated by Eqs. (7), (9) and (10). Let
$$N_1(R,\gamma ,t)=\frac{d𝒩_s(R,\gamma ,t)}{dt}.$$
(12)
is defined as the temporal radiance. In order to remove the singularity in Eq. (11), and also for conveniency of calculating high order scattering, we define the apparent temporal radiance as:
$$B_n(R,\gamma ,t)=\frac{4\pi R^2sin\gamma }{Q}\sigma ^nN_n(R,\gamma ,t)$$
(13)
Then Eq. (11) can be transferred to:
$$B_1(R,\gamma ,t)=ck_{ext.}\frac{2P\{cos[\theta (R,\gamma ,t)]\}sin\gamma }{(\xi ^22\xi cos\gamma +1)}exp(k_{ext.}ct).$$
(14)
where $`k_{ext.}=\alpha +\beta `$ is the total extinction coefficiency, $`\sigma =\frac{\beta }{\alpha +\beta }=\frac{\beta }{k_{ext.}}`$ is the single-scattering albedo.
Following the same way as , assume the apparent temporal radiance of (n-1) th order incidented at instant $`t^{}`$ at an angle $`\gamma ^{}`$ on a volume element $`dV^{}`$ located at distance $`R^{}`$ from the source through a field view solid angle $`d\mathrm{\Omega }^{}(\gamma ^{},\varphi ^{})`$ (see Fig. 2) is $`B_{n1}(R^{},\gamma ^{},t^{})`$.
Geometrically it is easy to get the following relationship:
$$R^{}=(R^2+r_{}^{}{}_{}{}^{2}2Rr^{}cos\gamma )^{1/2}.$$
(15)
$$ϵ^{}=cot^1\left(\frac{cos\gamma \frac{r^{}}{R}}{sin\gamma }\right),$$
(16)
The scattering angle $`\theta ^{}`$ is defined by
$$cos\theta ^{}=cos\gamma ^{}cosϵ^{}+sin\gamma ^{}sinϵ^{}cos\varphi ^{}.$$
(17)
then the nth order apparent temporal radiance is
$$B_n(R,\gamma ,t)=sin\gamma _0^D^{}𝑑r^{}_0^\pi 𝑑\gamma ^{}\left(\frac{R}{R^{}}\right)^2B_{n1}(R^{},\gamma ^{},t^{})\overline{P(cos\theta ^{})}.$$
(18)
where $`\overline{P(cos\theta ^{})}`$ is an azimuthally integrated scattering phase function :
$$\overline{P(cos\theta ^{}})=_0^{2\pi }d\varphi ^{}P(cos\theta ^{}).$$
(19)
where
$$D^{}=\frac{R}{2}(\xi \eta ).$$
(20)
$$SB+BC=ct^{}=ctr^{}.$$
(21)
$`\xi ,\eta `$ is defined in Eqs. (9) and (10).
With the above equations, theoretically we can calculate any order of scattering radiance for any given geometry setup (R,$`\gamma `$,t) and any atmosphere ( $`\sigma `$, $`P(cos\theta )`$, $`k_{ext.}`$ ). But in the real calculation, it will take tremendous of CPU time to calculate the radiance high than 3rd order of scattering, because of the 3-dimension integration in Eq. (18).
The total apparent temporal scattering radiance is:
$$B(R,\gamma ,t)=\underset{n=1}{\overset{n=\mathrm{}}{}}\sigma ^nB_n(R,\gamma ,t)$$
(22)
For simplification purpose, we define
$$c=k_{ext.}=1.$$
(23)
Then in the following, both the distance and time are in the unit of extinction length.
We then define regular grids for both the geometry setup and the atmosphere: R is from 0.0 to 6.0 with interval of 0.1; $`\gamma `$ is from $`0.01^o`$ to $`179.99^o`$ with interval $`0.2^o`$ when $`\gamma `$ $``$ $`2^o`$, with interval $`2^o`$ when $`\gamma `$ $`>`$ $`2^o`$; $`tR`$ is from 0 to 2 with interval 0.1; $`\sigma `$ is from 0.0 to 1.0 with interval 0.2; The atmospheric phase function is a weighted average of single-scattering phase functions $`P_R(cos\theta )`$ and $`P_A(cos\theta )`$, which represent, respectively, the Rayleigh and aerosol components. The weights are the corresponding scattering coefficients:
$$P(cos(\theta ))=\frac{\beta _RP_R(cos(\theta ))+\beta _AP_A(cos(\theta ))}{\beta _R+\beta _A}=\rho P_R(cos(\theta ))+(1\rho )P_A(cos(\theta )).$$
(24)
$$\rho =\frac{\beta _R}{\beta _R+\beta _A}.$$
(25)
$`\rho `$ is from 0.0 to 1.0 with interval of 0.2.
The grids chosed above to calculate the scattering radiance are simply trying to cover all the different situations the HiRes experiment will meet, and also considering the CPU time and calculation errors. To calculate the nth order, we only need to interpolate the result of the (n-1)th order, and then calculate the integration in Eq. (18) by using Gauss-Legendre procedure. The error of this calculation is about a few percent.
## 3 Monte Carlo Calculation
We made Monte Carlo calculations for the scattering radiance also. There are many references on this method . Here we follow the way as . The basic idea of MC method is to decompose light into a set of pencils of light that are called photons for the sake of brevity. The program follows the path of each photon inside the medium. Since the source is an isotropic point source, so it is very simple to setup the MC process (it can be generated to more completed case):
The geometry for the MC is shown in Fig. 3. Source S is assumed to be placed at the center of a sphere of radius $`R_{max}`$. The photon trajectory is then simulated by successive straight lines between collisions with scattering and absorbing centers inside the sphere of radius $`R_{max}`$. Each interaction between a photon and a scattering center obeys the law of single scattering. The scattering coefficient is used in computing the probability of the photon’s being scattered between distance l to (l+dl) as $`\beta exp(\beta l)dl`$. To take into account of absorption along the path, we weight a photon by factor $`W_i`$ that is initially set at 1 and then multiplied by $`exp(\alpha l)`$ at each collision. After the scattering, the new direction of the photon is determined by the scattering phase function $`P(cos\theta )`$. To take the advantage of the spherical symmetry, we look the whole spherical surface with radius R as a detector with normal direction.
When a photon crosses the detector’s surface, we record 2 quantities: the real path length ct, the angle $`\gamma `$ between the photon’s direction and the local normal direction of the detector. Then the weight of the photon $`W_i`$ is:
$$W_i=exp(\alpha ct).$$
(26)
If a photon does not hit the detector at time $`tt+dt`$, within direction $`\gamma \gamma +d\gamma `$, then
$$W_i=0.$$
(27)
The total number of photons hit the detector at time $`tt+dt`$, within direction $`\gamma \gamma +d\gamma `$, with the detector aerie $`dS`$ is:
$$_{mc}(R,\gamma ,t)dtd\mathrm{\Omega }dS=\underset{i=1}{\overset{i=Q_T}{}}W_i(R,\gamma ,t)$$
(28)
where $`Q_T`$ is the total number of input photons. (in the real calculation we simulate up to $`10^8`$ photons). From the spherical symmetry of the detector, it is easy to get the detector aerie:
$$dS=4\pi R^2cos(\gamma ).$$
(29)
and the field view solid angle $`d\mathrm{\Omega }`$:
$$d\mathrm{\Omega }=2\pi _\gamma ^{\gamma +d\gamma }sin(\gamma ^{})𝑑\gamma ^{}.$$
(30)
Same as the definition in Eqs. (13) and (22) From Eqs. (26), (27), (28), (29), (30) we get the apparent temporal radiance:
$$B_{mc}(R,\gamma ,t)=\frac{4\pi R^2sin\gamma }{Q_T}\frac{_{mc}(R,\gamma ,t)dtd\mathrm{\Omega }dS}{dtd\mathrm{\Omega }dS}=\frac{\underset{i=1}{\overset{i=Q_T}{}}W_i(R,\gamma ,t)}{Q_T}\frac{sin\gamma }{cos\gamma }\frac{1}{2\pi _\gamma ^{\gamma +d\gamma }sin(\gamma ^{})𝑑\gamma ^{}}.$$
(31)
In order to get good statistics, we simulate as many photons as possible within reasonable CPU time.
As described in Section 1. the atmosphere is completely defined by extinction coefficient $`k_{ext}`$, single-scattering albedo $`\sigma `$, and total phase function $`P(cos\theta )`$ which is described in Eqs. (24), (25). In the following we will describe how we model the phase functions.
The Rayleigh scattering phase function is modeled as the simplified form:
$$P_R(cos\theta )=\frac{3}{16\pi }(1+cos^2\theta )$$
(32)
Usually the aerosol phase function is described by a modified Henyey-Greenstein function, with an additional parameter f that gives rise to a backward peak:
$$P_A(cos\theta )=\frac{1g^2}{4\pi }[\frac{1}{(1+g^22gcos\theta )^{3/2}}+\frac{f(3cos^2\theta 1)}{2(1+g^2)^{3/2}}].$$
(33)
where $`g`$ is asymmetry parameter.
Since we are interested in the HiRes experiment, which is set in the desert of western USA. Here we will use a desert aerosol phase function calculated from Mie scattering theory with a aerosol particle size distribution function $`a^4`$, where $`a`$ is aerosol particle size. The phase functions are shown in Fig. 4. This phase function is very close to the real aerosol phase function at HiRes site .
## 4 Results
The recursive calculation was performed for up to 15 order for 6 optical depth.
In Figs. 5(a), 5(b), we show the first a few order of scattering for $`\sigma =0.8`$, $`\rho =0.8`$ $`R=1`$, and $`time`$ is defined as $`tR`$, which is the traveling time of scattering light t different from the direct light traveling time R. In 5(a) $`\gamma =2^o`$, in 5(b) $`\gamma =20^o`$. The signal is defined from Eqs. (18), (22), and the unit has been scaled by $`1/4\pi `$. In Figs. 6(a), 6(b), we show the first a few order of scattering for the same condition as Figs. 5. except $`R=4`$. From these Figs. we can see when the view angle $`\gamma `$ become bigger, or source detector distance R becomes longer, the high order scattering becomes important.
In Figs. 7(a), 7(b), we show the dependency of the total scattering radiance on the single scattering albedo $`\sigma `$. $`\sigma `$ is from 0.2 to 1.0 from bottom to up (when $`\sigma `$ is 0.0 there is not scattering, so we did not plot it out). The other conditions are $`\rho =0.8`$, $`R=1`$, 7(a) $`\gamma =2^o`$, in 7(b) $`\gamma =20^o`$. It is very clear when the single scattering albedo is bigger the scattering is stronger. In Figs. 8(a), 8(b), we show the dependency of the total scattering radiance on the rate of aerosol scattering coefficiency $`\rho `$. $`\rho `$ is from 0.0 to 1.0 from top to bottom. $`R=1`$, $`\sigma =0.8`$, 8(a) $`\gamma =2^o`$, 8(b) $`\gamma =20^o`$. Remember the aerosol phase function is more forward distributed( see Fig. 4). This can explain why the scattering is stranger when $`\rho `$ is high at small time range.
In Figs. 9(a), 9(b), we show the comparisons of the result from Monte Carlo to the iterative method. for $`\sigma =0.8`$, $`\rho =0.8`$ $`R=1`$, In 9(a) $`\gamma =3^o`$, in 9(b) $`\gamma =19^o`$. In Figs. 10(a), 10(b), we show the same thing as Fig 9. except $`R=4`$. The difference between the two methods are a few percent. except at large optical depth, where it reachs 10%. The difference mainly come from the iterative calculation, as pointed out in Sec. 2, by the interplating and Gauss-Legendre integration. The statistical fluctuations of MC results is poor even we used 48 hours of CPU on an 500-MHZ machine. From application point view, the iterative calculation is a more useful way. The error from iterative calculation can be lowed by more notes of interplation and integration. As will be pointed out later in Set. 5. the current result is good enough for the application on HiRes experiment.
## 5 Application
In this section, we will use the above iterative calculation results to study the multiply scattering effect in the energy estimate of cosmic ray by fluorescence technique, such as HiRes experiment. In the past 10 years, HiRes group has developed a detailed MC code to simulate the detector response for any given EAS.
The EAS was generated from CORSIKA package . The basical idea for the HiRes detector MC is that, start from the input shower, calculate the light produced at every stage, which includes fluorescence light and Cerenkov light. The fluorescence light then transmit to the detector. The Cerenkov light is mainly in the direction of the input shower, which can be scattered into the detector by Rayleigh scattering or aerosol scattering. The fluorescence light is directly related to number of charged particle at given stage, but the Cerenkov light is more complicated, and related to the history of the shower development. Fortunately, the Cerenkov light is much weak than the fluorescence light when the shower does not directly shoot to the detector. When the shower are within the direction of the detector, we simply drop it for better understanding of the data. All the good events we chose is only fluorescence light dominated. Then the above calculation can be directly used here.
Let $`𝒬(x)dx`$ is the total fluorescence light generated at stage $`xx+dx`$, F(x) is the detector coefficiency relate to the source at point x. r(x) is the distance from the source to the detector.
Then the direct light from stage $`xx+dx`$ to the detector is:
$$S_{dir.}(x)=F(x)\frac{𝒬(x)dx}{4\pi r^2(x)}exp(k_{ext.}r(x)).$$
(34)
If we define the shower start at time $`t=0`$, the time the direct light flight to the detector is: (Fig. 11.)
$$t(x)=\frac{r(x)+l(x)}{c}$$
(35)
let
$$c=k_{ext.}=1$$
(36)
Since the electronics of the detector will integrate all the signal come to the detector from point x within field view solid angle $`\mathrm{\Omega }(x)`$, within a fixed time window W, the scattering light from point $`x^{}`$ will also be integrated if it is in the time window and in the right direction. The geometric parameters for calculation of this effect is in Fig. 11. This effect can be represented as:
$$S_{sct.}(x)=F(x)_{t(x^{})t(x)+W/2}𝑑x^{}_{t(x)W/2}^{t(x)+W/2}𝑑t^{}_0^{\mathrm{\Omega }(x)}𝑑\mathrm{\Omega }\frac{𝒬(x^{})}{4\pi r^2(x^{})sin(\gamma (x^{},x))}B[r(x^{}),\gamma (x^{},x),t^{}].$$
(37)
in which $`B(r(x^{}),\gamma (x^{},x),t)`$ is calculated from Eqs. (14), (18), (22).
In the real calculation, we first build up a table about $`B(r(x^{}),\gamma (x^{},x),t)`$ as described in Sec. 2. For Eq. (37), we simple interplate the result from the table for any given geometry $`r(x^{}),\gamma (x^{},x),t`$ at any atmosphere ( $`k_{ext.},\sigma `$ ) condition. We chose W as 5.6 micro second, which is the real time integration window for HiRes detector. The detector openning field view solid angle $`\mathrm{\Omega }(x)`$ is a lit complicated, which is related to the stage of the shower and also the detector pixel.
The Calculation also depends on the atmosphere model, the following results based on the standard US atmosphere model. The aerosol part based on an exponentially decay model with an scale height of 1.2 Km, the extinction length on the ground is 10 Km.
The results are show in Figs. 12, 13, 14 15. for vertical EAS showers with primary energy $`10^{20}`$ eV at 10 km, 20 km, 30 km and 40 km away from the detector respectively. In 12(a), 13(a), 14(a), 15(a) we show the longitudinal shower profile at every stage from top of atmosphere to the ground. The solid lines are the total signal, the dashed lines are the direct light, the dot lines show the multiply scattering light. The X axis is the zenith angle for given stage, the Y axis is the signal measured at the detector site. Because of the detector trigger threshold condition, only the parts of shower around $`\pm 10^o`$ from the maximum signal point will be seen by the detector and gotten trigged.
In 12(b), 13(b), 14(b), 15(b) we show the rate of the multiply scattering light contribution to the directly light, define as $`S_{sct.}(x)`$ / $`S_{dir.}(x)`$ from Eqs. (34), (37). We can see the effect is around $`10\%`$ within the range the detector can see, and becomes stronger when the shower is farther. At the stage close to the ground, because both the atmosphere and aerosol density is higher, the extinction length is shorter, so the optical depth is longer the multiply scattering effect becomes stronger. For 12(b), the multiply scattering effect at early stage is stronger, this is simply because the geometrical distance is longer for early stage than later stage with a factor $`1/sin\theta `$ (Fig. 11.), which makes the optical depth longer even when the density of atmosphere become lower when it go up. The thickness of atmosphere is fixed, when the shower is farther like 13(b), 14(b) and 15(b). this geometrical effect becomes less important.
In order to see directly how the multiply scattering effect contribute to the energy estimate of this experiment. We draw the longitudinal profile versus the shower depth X in unit of $`gm/cm^2`$, as Fig (16). The shower to detector distance is 30 Km away. The aeries under the lines are directly proportional to the primary energy of the Cosmic ray as:
$$E_{primary}𝑑Xsig(X).$$
(38)
We can easily get the contribution to the primary energy estimation. The result shown in Fig. (17) is the rate of the aerie from scattering light to the aerie from the total light in percentage. Here we also studied how much this effect is, based on different aerosol models. We changed the horizontal aerosol extinction length from 10 Km to 20 Km, with the scale height change from 1.2 Km to 1.5 Km. We can see the scatteing light contributes to the energy estimation in about 10%.
It is interesting to see that when the shower to detector distance is moved from 10 Km to 40 Km, the scattering effect does not increase dramatically. This is because we only intergrate the signal within a fixed time window.
More generally, the EAS can not be verticle, we also simulated showers with arbitrary zenith angle, and the multiply scattering effect on the shower longitudinal profile will change shower by shower depend on the shower gemetry. But the effect on the primary energy estimate is still about 10%. In the real HiRes data analysis, we will calculate this effect for every shower.
As point in Sec. 4, there is about $`10\%`$ error in the scattering radiance calculation. But since the multiply scattering light itself is only about $`10\%`$ for the total light, so the error transferred from the scaterring radiance calculation is only about $`1\%`$, which is much small than the systematics of this experiment.
## 6 conclusion
The properties of an aureole about a point source that are due to atmospheric scattering were calculated by two approaches: the Monte Carlo method and iterative technique. For the iterative technique we have an applicable approach to calculate as high order of scattering as you want within reasonable calculation time for any kind of atmosphere and any detector setup. The results match very well with the MC results within error of about 10%. The iterative method is more useful than MC method if we are interested in the temporal aureole problem. We studied the properties of the temporal aureole depend mainly on single-scattering albedo, scattering phase function, optical depth between source and detector, detector view angle.
Then we apply the multiply scattering property of a point source to the EHE cosmic ray shower detected by the atmosphere fluorescence technique. And found the contribution of the scattering light is about $`10\%`$ , which depends on how far the shower is from the detector in the unit of optical depth. This means the energy estimate will be corrected by $`10\%`$ because this effect.
This approach can also be generalized to anisotropic cases. Such as the Lidar system used to calibrate the atmosphere . In the world there are some Cerenkov light detector like . Since for anisotropic case the multiply scattering effect will become stronger , in order to understand better of the system, the multiply scattering effect should be studied in those cases.
## 7 Acknowledgments
The author thanks Prof. W. Lee, Prof. P. Sokolsky, Dr. c. Zen, Dr. L. Wiencke for the creative discussions. The support from High Energy Astrophysics Institute at University of Utah is gratefully acknowledged during his visiting time.
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# NONCOMMUTATIVE RATIONAL FUNCTIONS AND FARBER’S INVARIANTS OF BOUNDARY LINKS
## 1 Introduction
Alexander polynomials of knots and links are constructed from the Alexander module over $`𝐙[t_1,\mathrm{},t_\mu ]`$ (where $`\mu `$ is the number of components of the link) which corresponds to the maximal abelian covering of the link complement. They contain essentially all commutative information about the link (or rather its fundamental group).
Farber in constructed invariants of $`\mu `$-component boundary links with values in the algebra of noncommutative rational functions. When $`\mu =1`$ (i.e. when link is a knot) this invariant is equivalent to the Alexander polynomial.
Farber associated to any $`n`$-dimensional boundary link a sequence of rational noncommutative power series $`\varphi _i`$, $`i=1,\mathrm{},n`$ providing rather strong link invariants. To compute $`\varphi _i`$ one has to calculate a finite number of integers (the traces of certain linear maps acting on the homology of a Seifert manifold).
The core of Farber’s construction is a characteristic rational power series $`\varphi _M`$ associated to any finite-dimensional module $`M`$ over a certain $`k`$-algebra $`P_\mu `$ with $`\mu +1`$ generators having a subalgebra isomorphic to the algebra $`kx_1,\mathrm{},x_\mu `$ of noncommutative polynomials. When $`k`$ is an algebraically closed field of zero characteristic and the $`P_\mu `$-module $`M`$ is semisimple, the series $`\varphi _M`$ determines $`M`$ up to an isomorphism.
In this paper we construct a characteristic rational power series $`\chi _N`$ for any finite-dimensional $`kx_1,\mathrm{},x_\mu `$-module $`N`$. When $`k`$ is an algebraically closed field of arbitrary characteristic and $`N`$ is semi-simple, the series $`\chi _N`$ determines $`N`$ up to an isomorphism.
If $`N`$ is also a $`P_\mu `$-module and the actions of $`P_\mu `$ and $`kx_1,\mathrm{},x_\mu `$ are compatible, the series $`\varphi _N`$ and $`\chi _N`$ determine each other.
Note that our construction of $`\chi _N`$ is simpler than the construction of $`\varphi _N`$ and this gives us a hope that this invariant can be applied to a wider class of links.
Alexander polynomial can be computed in terms of minors of the Alexander matrix and it would be very interesting to have a similar interpretation of its noncommutative generalization. In this paper we make the first step in the direction of applying noncommutative determinants to obtain invariants of links. We show that characteristic series $`\chi `$ and $`\varphi `$ can be expressed by means of quasideterminants, a noncommutative generalization of determinants introduced by Gelfand and Retakh (see, for example, ).
The first author was suppported in part by Arkansas Science and Technology Authority and the second one by NSERC(Canada).
## 2 Characteristic functions of finite-dimensional $`kX`$-modules
Let $`k`$ be a field and $`X=\{x_1,\mathrm{},x_\mu \}`$ be a set of noncommuting variables. Denote by $`kX`$ (resp. $`kX`$) the $`k`$-algebra of noncommutative polynomials (resp. formal series) in $`x_1,\mathrm{},x_\mu `$. Denote by $`X^{}`$ the free monoid generated by $`X`$. Then $`X^{}`$ is a $`k`$-basis of $`k<X>`$. The elements of $`X^{}`$ will be called words.
###### Definition 2.1
The ring $``$ of noncommutative rational series is defined as the smallest $`k`$-subalgebra of $`kX`$ satisfying the following properties.
(i) $`kX`$;
(ii) if $`g`$ and $`g`$ is invertible in $`kX`$ then $`g^1`$.
General theory of rational series may be found in . Here we remind just basic facts.
Rational series are called sometimes noncommutative rational functions. They can be characterized in terms of Fox derivatives.
###### Definition 2.2
The Fox derivative $`_i`$ on $`kX`$ is a $`k`$-endomorphism defined by
$$_i(x_jw)=\delta _{ij}w,_i1=0$$
for any word $`wX^{}`$.
It satisfies the Leibniz rule: for formal series $`f`$ and $`g`$
$$_i(fg)=(_if)g+ϵ(f)_ig,$$
where $`ϵ(f)`$ means the constant term of $`f`$.
Fox derivatives turn $`kX`$ into a right $`kX`$-module. This structure may also be described by the following construction: the action of polynomial $`P`$ on series $`S`$ is given by $`(SP,w)=(S,Pw)`$, for $`wX^{}`$, where $`(,)`$ denotes the canonical pairing between $`kX`$ and $`kX`$, defined by
$$(S,P)=\underset{wX^{}}{}s_wp_w,\text{if}S=\underset{wX^{}}{}s_ww,P=\underset{wX^{}}{}p_ww.$$
In this notations $`wx_i=_iw`$.
###### Remark 2.3
Let $`\mathrm{\Lambda }=k[F_\mu ]`$ be the group algebra of the free group $`F_\mu `$ with $`\mu `$ generators $`g_1,\mathrm{},g_\mu `$. The Magnus embedding
$$\mathrm{\Lambda }kX,g_i1+x_i,g_i^1\underset{n0}{}(x_i)^n$$
maps $`\mathrm{\Lambda }`$ into $``$.
By Leibniz rule the image of $`\mathrm{\Lambda }`$ under the Magnus embedding is a $`kX`$-submodule of $``$.
The structure of $`kX`$-module on $`kX`$ that we consider here is different from the one obtained by the canonical embedding $`kXkX`$.
Rational series can be characterized in terms of the Fox derivatives. The following proposition belongs to Schützenberger and Fliess .
###### Proposition 2.4
A series $`f`$ is rational if and only if the $`k`$-vector space spanned by all Fox derivatives $`_{i_k}\mathrm{}_{i_2}_{i_1}f`$ of $`f`$ is finite-dimensional.
Therefore, to each rational function $`\chi `$ there corresponds a finite-dimentional vector space $`M_\chi `$ spanned by the Fox derivatives of $`\chi `$ of all orders. The space $`M_\chi `$ has a natural $`kX`$-module structure given by $`x_i_i,i=1,\mathrm{},\mu `$.
Vice versa, to each finite-dimensional $`kX`$-module $`M`$ there corresponds a rational function $`\chi _M`$.
###### Definition 2.5
Let $`M`$ be a $`kX`$-module, finite dimensional over $`k`$, and $`u_i\text{End}_k(M)`$ the endomorphism of $`M`$ corresponding to the action of $`x_ikX`$.
Define a $`k`$-algebra homomorphism
$$\alpha :kX\text{End}_kM,x_iu_i.$$
The characteristic function $`\chi _MkX`$ of $`M`$ is defined by
$$\chi _M=\underset{wX^{}}{}Tr(\alpha (w))w.$$
###### Example 2.6
Let $`k`$ be an algebraically closed field and $`\mu =1`$, i.e. $`M`$ is a $`k[x]`$-module. Let $`d=\text{dim}(M)`$ and $`\lambda _i,i=1,\mathrm{},d`$ be eigenvalues of $`\alpha (x)`$. Then
$$\chi _M=\underset{1id}{}\frac{1}{1\lambda _ix}.$$
It is easy to prove the following fact.
###### Proposition 2.7
The characteristic function $`\chi _M`$ is a rational series which is additive for short exact sequences of $`kX`$-modules.
The following theorem shows that a simple $`kX`$-module can be recovered from its characteristic function.
###### Theorem 2.8
Let $`k`$ be an algebraically closed field and $`M`$ be a simple finite-dimensional $`kX`$-module. Then $`M_{\chi _M}`$ is isomorphic (as a $`kX`$-module) with the direct sum of $`d=\mathrm{dim}(M)`$ copies of $`M`$.
We will deduce this result from a theorem of Fliess which we state now.
Let $`M`$ be a finite dimensional right $`kX`$-module. Then the dual space $`M^{}`$ is a left $`kX`$-module. Let $`mM`$ (resp. $`\varphi M^{}`$) generate $`M`$ (resp. $`M^{}`$) under the $`kX`$-action.
Define
$$S=\underset{wX^{}}{}(m(\alpha (w)))\varphi w,$$
where $`\alpha `$ is defined as before ($`x_i`$ acts on the right on $`M^{}`$, $`i=1,\mathrm{},\mu `$ and similarly $`\varphi `$ acts on the right on $`M`$).
###### Theorem 2.9 (Fliess )
Module $`M`$ is isomorphic to $`M_S`$ as right $`kX`$-modules and, under this isomorphism $`m`$ (resp. $`\varphi `$) corresponds to $`S`$ (resp. to the constant term map $`ϵ:kXk`$).
The proof is mechanical: we define the isomorphism by $`m^{}_{wX^{}}(m^{}(\alpha w))\varphi w`$. Then a straightforward verification shows that it is well-defined, injective and surjective. Note that $`k`$ needs not to be algebraically closed here.
Proof of Theorem 2.8. Let $`\{m_1,\mathrm{},m_d\}`$ be a basis of $`M`$, and $`\{\varphi _1,\mathrm{},\varphi _d\}`$ be the dual basis of $`M^{}`$. Define $`S_i=_{wX^{}}(m_i(\alpha w))\varphi _i`$. Then $`\chi _M=_{i=1}^dS_i`$.
Let $`N`$ be a direct sum of $`d`$ copies of $`M`$, $`m=(m_1,\mathrm{},m_d)N`$, $`\varphi =(\varphi _1,\mathrm{},\varphi _d)N^{}`$. Then $`\chi _M`$ corresponds to the triple $`N,m,\varphi `$ as is described in the Theorem 2.8, so that $`M_{\chi _M}`$ is isomorphic to $`N`$, if we can verify the hypothesis of this theorem.
Since $`k`$ is algebraically closed, and $`M`$ is simple $`kX`$-module, the subalgebra of $`\text{End}_k(M)`$ is, by Burnside’s theorem, all of $`\text{End}_k(M)`$. Note that this algebra coincides with $`\{\alpha P|PkX\}`$. Hence, we may find $`PkX`$ such that $`m_i(\alpha P)=m_i`$ and $`m_j(\alpha P)=0`$ for $`ji`$. Thus $`m(\alpha P)=(0,\mathrm{},0,m_i,0,\mathrm{},0)`$ which implies that $`m`$ generates all of $`N`$ under this action, $`M`$ being simple. For the other hypothesis of Theorem 2.8 it is similar: one replaces $`M`$ by its dual $`M^{}`$ on which the $`u_j`$’s act simply on the left.
###### Remark 2.10
Note that the result is not true if $`k`$ is not algebraically closed. The smallest example is $`k=\mathrm{R}`$, $`M=ke_1ke_2`$, and action of $`u`$ given by imitating multiplication by complex number $`i`$, i.e.
$$e_1u=e_2,e_2u=e_1.$$
Then $`\chi =\chi _M=22x^2+2x^42x^6+\mathrm{}`$, and $`M_\chi `$ is isomorphic to $`M`$ and not to $`MM`$.
There is certainly a version of Theorem 2.8 when $`k`$ is not algebraically closed, where the “arithmetic” of $`k`$ plays some role. Also, for effective computational purposes, $`k=𝐐`$ seems to be the best field.
If $`M`$, instead of being simple, is only semi-simple, there is a variant of Theorem 2.8.
###### Theorem 2.11
Let $`M`$ be a semi-simple module over $`kX`$. If $`M=\underset{r=1}{\overset{q}{}}M_r,`$ where $`M_r,r=1,\mathrm{},q`$ is simple, then $`M_{\chi _M}`$ is isomorphic to $`\underset{r=1}{\overset{q}{}}\underset{r=1}{\overset{\mathrm{dim}(M_r)}{}}M_r.`$
As an application of Theorem 2.8 we obtain the following important result.
###### Theorem 2.12
Let $`k`$ be algebraically closed. Two finite-dimensional semi-simple $`kX`$-modules are isomorphic if and only if their characteristic functions coincide.
## 3 $`P_\mu `$-modules
In his study of boundary links and links modules M. Farber introduced an algebra $`P_\mu `$ with $`\mu +1`$ generators and defined characteristic functions of finite-dimensional $`P_\mu `$-modules.
Let $`P_\mu `$ be a $`k`$-algebra defined by $`\mu +1`$ generators $`z,\pi _i,i=1,\mathrm{},\mu ,`$ and relations
$$\pi _i\pi _j=\delta _{ij}\pi _i,\pi _1+\mathrm{}+\pi _\mu =1.$$
Modules over $`P_\mu `$ are automatically $`kX`$-modules via the homomorphism $`\delta :kXP_\mu `$ defined by $`x_iz\pi _i`$. Denote $`\delta (x_i)`$ by $`_i`$, $`i=1,\mathrm{},\mu `$.
###### Proposition 3.1
The homomorphism $`\delta `$ is an embedding.
Proof. Indeed, this follows from the two simple facts that, first, the image of $`\delta `$ in $`P_\mu `$ coincides with the subalgebra $`zP_\mu P_\mu `$, and that, second, this subalgebra is freely generated by $`\mu `$ elements $`z,z\pi _1,\mathrm{},z\pi _{\mu 1}`$. Using the identity $`_i\pi _i=1`$, every word of type $`z^{a_1+1}\pi _{i_1}z^{a_2+1}\pi _{i_2}\mathrm{}z^{a_k},a_j0,1i_j\mu 1`$ can be written as a polynomial in $`_i=\delta (x_i)=z\pi _i`$ in a unique way.
Recall that $`X^{}`$ is the free monoid generated by the alphabet $`X=\{x_1,\mathrm{},x_\mu \}`$. Every word $`wX^{}`$ can be uniquely written as $`w=x_jw^{}`$ for some $`j`$ and $`w^{}X^{}`$. Define action of the generators of $`P_\mu `$ on $`X^{}`$ as
$$\pi _i(x_jw)=\delta _{ij}x_jw,i,j=1,\mathrm{},\mu ,z(x_kw)=w.$$
This gives a $`P_\mu `$-module structure on the augmentation ideal $`kX_+`$ of $`kX`$ (i.e. the formal series without constant terms).
Denote by $`\mathrm{\Lambda }_0`$ the image of the Magnus embedding of $`\mathrm{\Lambda }=k[F_\mu ]kX`$. The action of $`P_\mu `$ on $`X^{}`$ defines a $`P_\mu `$-module structure on $`kX/\mathrm{\Lambda }_0`$.
###### Proposition 3.2
Let $``$ be the ring of noncommutative rational series. Then $`/\mathrm{\Lambda }_0`$ is invariant under $`P_\mu `$-action on $`kX/\mathrm{\Lambda }_0`$.
Farber introduced the following notion of characteristic function for finite-dimensional $`P_\mu `$-modules.
###### Definition 3.3
Let $`A`$ be a finite-dimensional $`P_\mu `$-module. Its characteristic function is the series
$$\varphi _A=\underset{k=1}{\overset{\mu }{}}\underset{\alpha }{}\text{Tr}(\pi _k_\alpha )x^\alpha x_k,$$
where $`\alpha =(i_1,\mathrm{},i_p)`$, $`_\alpha =_{i_p}\mathrm{}_{i_1}`$, $`x^\alpha =x_{i_1}\mathrm{}x_{i_p}`$.
###### Proposition 3.4
The characteristic function is a rational series which is additive for short exact sequences of $`P_\mu `$-modules.
###### Corollary 3.5
Let $`M`$ be finite-dimensional $`P_\mu `$-module and let $`C_i`$, be its distinct composition factors appearing with multiplicities $`m_i1`$, $`i=1,\mathrm{},n`$. Then
$$\varphi _M=\underset{i=1}{\overset{n}{}}m_i\varphi _{C_i}.$$
###### Example 3.6
If $`\mu =1`$ then
$$\varphi _A=\underset{j=1}{\overset{n}{}}\frac{x}{1+\lambda _jx},$$
where $`n=\text{dim}_kA`$ and $`\lambda _1,\mathrm{},\lambda _n`$ are the eigenvalues of the operator $`z`$.
In general, function $`\varphi _A`$ is an invariant capable of capturing only semi-simple information about $`P_\mu `$-modules. To study non-semisimple modules we will need more subtle invariants.
Farber proved the following result similar to our Theorem 2.12.
###### Theorem 3.7
Suppose $`k`$ is an algebraically closed field of characteristic zero. Let $`A`$ and $`B`$ be finite-dimensional semi-simple $`P_\mu `$-modules. Then $`\varphi _A=\varphi _B`$ if and only if $`A`$ is isomorphic to $`B`$.
This result follows from the version of our Theorem 2.8 for $`P_\mu `$-modules. It shows that the $`P_\mu `$-submodule of $`/\mathrm{\Lambda }_0`$ generated by $`\chi _A`$ is closely related to $`A`$.
###### Definition 3.8
A $`P_\mu `$-module $`M`$ is called primitive if one of the generators $`\pi _1,\mathrm{},\pi _\mu ,z`$ of $`P_\mu `$ acts as the identity on $`M`$ and the other generators act trivially.
A $`P_\mu `$-module $`A`$ is called primitive-free if it has no primitive composition factors.
###### Theorem 3.9 ()
If $`k`$ is algebraically closed and $`A`$ is simple non-primitive $`P`$-module of dim$`A=d`$, then $`P_\mu \varphi _A/\mathrm{\Lambda }_0`$ is $`P_\mu `$-isomorphic to $`dA`$, the direct sum of $`d`$ copies of $`A`$.
###### Remarks 3.10
(1) If one takes an alphabet with $`\mu +1`$ letters, then the characteristic series $`\chi _M`$ on this alphabet for a $`P_\mu `$-module is equivalent to Farber’s series; i.e. the knowledge of one is equivalent to the knowledge of the other.
(2) It it easy to construct example of simple $`P_\mu `$-modules which are not $`kx_1,\mathrm{},x_\mu `$-simple under the homomorphism $`x_iz\pi _i`$. Let $`M`$ be a three-dimensional vector space with a basis $`e_i`$, $`i=1,2,3`$. Define an action of algebra $`P_3`$ on $`M`$ by setting $`\pi _ie_j=\delta _{ij}e_j,i,j=1,2,3`$. The action of $`z`$ in this basis is given by a matrix
$$\left[\begin{array}{ccc}\hfill 1& \hfill 0& \hfill 1\\ \hfill 1& \hfill 1& \hfill 2\\ \hfill 0& \hfill 1& \hfill 1\end{array}\right].$$
One can see that $`M`$ is a simple $`P_3`$-module under this action, and the two-dimensional image of $`z`$ is invariant under action of $`_i`$ for $`i=1,2,3`$.
Therefore, it is not clear whether the main results about semi-simple $`P_\mu `$-modules of this section follow directly from the results obtained in Section 2.
## 4 Quasideterminants and characteristic functions
In this section we express characteristic functions constructed in Sections 2 and 3 via quasideterminants.
Quasideterminants were introduced by Gelfand and Retakh (see, for example, ) and are defined as follows. Let $`A`$ be an $`m\times m`$-matrix over an algebra $`R`$. For any $`1i,jm`$, let $`r_i(A)`$, $`c_j(A)`$ be the i-th row and the j-th column of $`A`$. Let $`A^{ij}`$ be the submatrix of $`A`$ obtained by removing the i-th row and the j-th column from $`A`$. For a row vector $`r`$ let $`r^{(j)}`$ be $`r`$ without the j-th entry. For a column vector $`c`$ let $`c^{(i)}`$ be $`c`$ without the i-th entry. Assume that $`A^{ij}`$ is invertible. Then the quasideterminant $`|A|_{ij}R`$ is defined by the formula
$$|A|_{ij}=a_{ij}r_i(A)^{(j)}(A^{ij})^1c_j(A)^{(i)},$$
where $`a_{ij}`$ is the $`ij`$-th entry of $`A`$. Let $`B=(b_{ij})`$ be a matrix of order $`n`$ with formal entries and $`E_n`$ a unit $`n`$-matrix. The following proposition was proved in .
###### Proposition 4.1
In the ring $`𝐙b_{ij}`$ of formal series with integer coefficients generated by noncommuting variables $`b_{ij}`$ we have
$$|E_nB|_{ii}^1=1+\underset{k_1,\mathrm{},k_p}{}b_{ik_1}b_{k_1k_2}\mathrm{}b_{k_pi},$$
where the sum is over all $`1k_1,\mathrm{},k_pn`$, $`p=1,2,\mathrm{}`$.
Let $`k`$ be a field, $`X=\{x_1,\mathrm{},x_\mu \}`$ set of noncommuting variables, and $`M`$ be an $`n`$-dimensional $`kX`$-module. Choose a basis of $`M`$ and let $`A_i`$, $`i=1,\mathrm{},\mu ,`$ be the matrix of the operator $`x_i`$ in this basis.
The following proposition gives an expression of the characteristic function $`\chi _M`$ of $`M`$ via quasideterminants.
###### Proposition 4.2
The sum
$$\underset{i=1}{\overset{n}{}}|E_nx_1A_1\mathrm{}x_\mu A_\mu |_{ii}^1$$
does not depend on the choice of the basis and is equal to the characteristic function $`\chi _M`$ of the $`kX`$-module $`M`$.
A similar result holds for Farber’s characteristic functions.
Let $`M`$ be a finite-dimensional $`P_\mu `$-module. Fix a basis of $`M`$ and denote by $`B_i`$, $`i=1,\mathrm{},\mu `$ the matrix of the the operator $`(1z)\pi _i`$ in this basis.
###### Proposition 4.3
The sum
$$\underset{i=1}{\overset{n}{}}|E_nx_1B_1\mathrm{}x_\mu A_\mu |_{ii}^1$$
does not depend on the choice of the basis and is equal to the characteristic function $`\varphi _M`$ of the $`P_\mu `$-module $`M`$.
## 5 Link modules
Let $`F_\mu `$ be a free group with $`\mu `$ generators $`t_i`$ and $`\mathrm{\Lambda }=k[F_\mu ]`$ its group ring. A finitely-generated left $`\mathrm{\Lambda }`$-module $`M`$ is called link module if $`\text{Tor}_q^\mathrm{\Lambda }(k,M)=0`$ for all $`q0`$, where $`k`$ is regarded as a right $`\mathrm{\Lambda }`$-module via the augmentation map.
The following characterization of link modules allows to consider them as $`P_\mu `$-modules.
###### Theorem 5.1 (Sato )
$`M`$ is a link module if and only if every element $`m`$ of $`M`$ has a unique representation
$$m=\underset{i=1}{\overset{\mu }{}}(t_i1)m_i,m_iM.$$
###### Corollary 5.2
In the notations of the theorem, formulas $`\pi _im=(t_i1)m_i,zm=m_i`$, make every link module a $`P_\mu `$-module. In addition, $`_im=m_i`$.
Note that every $`P_\mu `$-module is also a $`kx_1,\mathrm{},x_\mu `$-module under the canonical embedding $`\delta :kx_1,\mathrm{},x_\mu P_\mu `$.
###### Definition 5.3 ()
A finite-dimensional $`P_\mu `$-submodule $`A`$ of a link module $`M`$ is called lattice if $`A`$ generates $`M`$ over $`\mathrm{\Lambda }`$.
It is easy to prove the following fact.
###### Theorem 5.4 ()
Every link module contains a lattice. Intersection of all lattices in $`M`$ is a lattice which is called the minimal lattice of $`M`$.
###### Definition 5.5
If $`M`$ is a link module, we set $`\varphi _M=\varphi _A`$ and $`\chi _M=\chi _A`$, where $`A`$ is the minimal lattice in $`M`$.
###### Theorem 5.6 (Farber )
Let $`M`$ and $`N`$ be semi-simple link modules with $`\varphi _M=\varphi _N`$. If the field $`k`$ is algebraically closed and $`\mathrm{char}(k)=0`$, then $`M`$ and $`N`$ are isomorphic as $`P_\mu `$-modules.
The proof of the theorem is based on Theorem 3.7.
## 6 Boundary links
An $`n`$-dimensional $`\mu `$-component link is an oriented smooth submanifold $`\mathrm{\Sigma }`$ of $`S^{n+2}`$, where $`\mathrm{\Sigma }=\mathrm{\Sigma }_1\mathrm{}\mathrm{\Sigma }_\mu `$ is an ordered disjoint union of $`\mu `$ submanifolds of $`S^{n+2}`$, each diffeomorphic to $`S^n`$. It is called a boundary link if there exists an oriented submanifold $`V`$ of dimension $`n+1`$ of $`S^{n+2}`$, such that $`V=V_1\mathrm{}V_\mu `$ and $`V_i=\mathrm{\Sigma }_i,i=1,\mathrm{},\mu `$. If, in addition, each $`V_i`$ is connected, we say that $`V`$ is a Seifert manifold for $`\mathrm{\Sigma }`$.
The homology groups of the Seifert manifold $`H_{}(V;k)`$ have several natural operations. The first $`\mu `$ operations are the projections $`\pi _i:H_{}(V;k)H_{}(V_i;k)`$. The last operation $`z`$ is defined as follows. For $`Y=S^{n+2}V`$ let $`I_\pm :VY`$ be small shifts in the direction of positive and negative normals to $`V`$, respectively. The map
$$i_+i_{}:H_k(V)H_k(Y)$$
is an isomorphism for any $`k=0,1,2,\mathrm{}`$ (see ) and we define $`z`$ by $`z=(i_+i_{})^1i_+`$.
Let $`\mathrm{\Sigma }`$ be a link in $`S^{n+2}`$ and $`X=S^{n+2}T(\mathrm{\Sigma })`$ the complement of its tubular neighborhood.
Fix a base point $`X`$. Then connecting the meridians of $`\mathrm{\Sigma }`$ (small loops around each component $`\mathrm{\Sigma }_i`$) with $``$ we obtain elements $`m_1,\mathrm{},m_\mu \pi _1(X,)`$ defined up to conjugation. If $`\mathrm{\Sigma }=V`$ is a boundary link then there is an epimorphism (defined up to conjugation) $`\sigma :\pi _1(X,)F_\mu `$ defined as follows. If $`\alpha `$ is a loop in $`X`$ intersecting $`V`$ transversally, first in component $`V_{i_1}`$, then $`V_{i_2}`$, etc, we set $`\sigma [\alpha ]=t_{i_k}^{\epsilon _k}\mathrm{}t_{i_2}^{\epsilon _2}t_{i_1}^{\epsilon _1}`$, where $`\epsilon _i=\pm 1`$ is the local intersection index between $`\alpha `$ and $`V`$ at $`i`$-th intersection point.
Consider the covering $`\stackrel{~}{X}X`$ corresponding to the kernel of $`\sigma `$. The group of deck transformations of this covering space is $`F_\mu `$ and therefore the homology $`H(\stackrel{~}{X};k)`$ is a $`\mathrm{\Lambda }=k[F_\mu ]`$-module. For $`0in`$ this module is a link module.
There is a canonical lifting map $`S^{n+2}\mathrm{\Sigma }\stackrel{~}{X}`$ giving a homomorphism of $`P_\mu `$-modules $`f:H_i(V;k)H_i(\stackrel{~}{X};k)`$. The image of $`f`$ is a lattice. If $`H_i(V;k)`$ has no primitive $`P_\mu `$-submodules, then $`f`$ is a monomorphism and the image of $`H_i(V;k)`$ is a minimal sublattice of $`H_i(\stackrel{~}{X};k)`$ (cf. ).
Applying the constructions of Sections 2 and 3 to the minimal lattices in homology groups $`H_i(\stackrel{~}{X},k)`$ we obtain sequences of rational series $`\chi _i`$ and $`\varphi _i`$, $`i=1,\mathrm{},\mu `$. Thus we have a sequence of noncommutative invariants associated to every boundary link.
These invariants are stronger than the well-known commutative invariants. In particular, Farber constructed an example of a link one of whose Alexander modules vanishes, but the corresponding characteristic function is non-trivial.
###### Example 6.1
For $`\mu =1`$ the invariants $`\varphi `$, $`\chi `$, and the Alexander polynomial determine each other. Let $`\mathrm{\Delta }_i(t)`$ be the Alexander polynomial of $`H_i(\stackrel{~}{X};𝐐)`$ where $`\stackrel{~}{X}`$ is an infinite cyclic cover of the complement of the knot. Then if
$$\mathrm{\Delta }_i(t)=\underset{j=1}{\overset{n}{}}(t\nu _j),\nu _i𝐂,$$
the characteristic functions are
$$\chi (x)=\underset{j=1}{\overset{n}{}}1/(1\lambda _jx)$$
and
$$\varphi (x)=\underset{j=1}{\overset{n}{}}x/(1+\lambda _jx),$$
where
$$\lambda _j=1/(1\nu _j),j=1,\mathrm{},n.$$
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# Optical Absorption of an Interacting Many-Polaron Gas
## I Introduction
As is the case for polar semiconductors and ionic crystals , insight into the nature of polarons in high-temperature superconductors can be gained by studying the optical properties of these materials. The goal of the present paper is to present a theory of the optical conductivity of a system of continuum polarons at any density, including many-body effects between the constituent charge carriers, for small and intermediate values of the electron-phonon coupling constant and for zero temperature. The method we develop here is based on the variational method introduced by Lemmens, Devreese, Brosens (LDB) for the ground state energy of the many-polaron gas. The advantage of this approach over other theories of many-polaron optical absorption is that it allows to include the many-body effects in the system of the constituent charge carriers of the polaron gas on the level of the dynamical structure factor of the underlying electron (or hole) system. Thus it is possible to select the level of approximation used in the treatment of the many-polaron gas by choosing the appropriate expression of the dynamical structure factor for the electron (or hole) system.
Recently the infrared spectrum of cuprates has been the subject of intensive investigations , especially in the case of the neodymium-cerium cuprate family Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4-y</sub> (NCCO) . Several optical absorption features in the infrared cuprate spectrum have been tentatively associated with large polarons or with a mixture of large and small (bi)polarons . These comparisons with polaron theory were derived using a single-polaron picture, so that the density (doping) dependence of the optical absorption spectra could not be studied in detail. The many-body theory of the $`N`$-polaron spectrum, presented here, allows to study the density (doping) dependence of optical absorption spectra. As a first application of the many-polaron optical absorption theory introduced here, a preliminary comparison is presented between the theoretical many-polaron optical absorption derived in the current work and the mid-infrared spectrum of the neodymium-cerium cuprates recently determined experimentally by Calvani and co-workers .
## II Optical absorption in the many-polaron system
### A The LDB variational wave function for a many-polaron system
The Hamiltonian of a system of $`N`$ interacting continuum polarons is given by:
$$H_0=\underset{j=1}{\overset{N}{}}\frac{p_j^2}{2m_b}+\underset{𝐤}{}\mathrm{}\omega _{\text{LO}}a_𝐤^+a_𝐤+\underset{𝐤}{}\underset{j=1}{\overset{N}{}}\left[e^{i𝐤.𝐫_j}a_𝐤V_𝐤+e^{i𝐤.𝐫_j}a_𝐤^+V_𝐤^{}\right]+\frac{e^2}{2\epsilon _{\mathrm{}}}\underset{j=1}{\overset{N}{}}\underset{\mathrm{}(j)=1}{\overset{N}{}}\frac{1}{|𝐫_i𝐫_j|},$$
(1)
where $`𝐫_j,𝐩_j`$ represent the position and momentum of the $`N`$ constituent electrons (or holes) with band mass $`m_b`$; $`a_𝐤^+,a_𝐤`$ denote the creation and annihilation operators for longitudinal optical (LO) phonons with wave vector $`𝐤`$ and frequency $`\omega _{\text{LO}}`$; $`V_𝐤`$ describes the amplitude of the interaction between the electrons and the phonons; and $`e`$ is the elementary electron charge. The ground state energy of this many-polaron Hamiltonian has been studied before by LDB , for weak and intermediate strengths of the electron-phonon coupling, by introducing a variational wave function:
$$|\psi _{\text{LDB}}=U|\varphi |\phi _{\text{el}},$$
(2)
where $`|\phi _{\text{el}}`$ represents the ground-state many-body wave function for the electron (or hole) system and $`|\varphi `$ is the phonon vacuum, and $`U`$ is a many-body unitary operator which determines a canonical transformation for a fermion gas interacting with a boson field:
$$U=\mathrm{exp}\left\{\underset{j=1}{\overset{N}{}}\underset{𝐤}{}\left(f_𝐤a_𝐤e^{i𝐤.𝐫_j}f_𝐤^{}a_𝐤^+e^{i𝐤.𝐫_j}\right)\right\}.$$
(3)
In the limit of one fermion, $`U`$ reduces to a canonical transformation inspired by Tomonaga and applied later by several workers, after Lee, Low and Pines , but always for one particle-theories. In LDB , this canonical transformation was extended and used to establish a many-fermion theory. The $`f_𝐤`$ were determined variationally resulting in
$$f_𝐤=\frac{V_𝐤}{\mathrm{}\omega _{LO}+{\displaystyle \frac{\mathrm{}^2k^2}{2m_bS(𝐤)}}},$$
(4)
for a system with total momentum $`𝐏=_j𝐩_j=0`$. In this expression, $`S(𝐤)`$ represents the static structure factor of the constituent interacting many electron or hole system :
$$NS(𝐤)=\underset{j=1}{\overset{N}{}}\underset{j^{}=1}{\overset{N}{}}e^{i𝐤.(𝐫_j𝐫_j^{})}.$$
(5)
The angular brackets $`\text{}`$ represent the expectation value with respect to the ground state. It may be emphasized that (3), although it appears like a straightforward generalization of the one-particle transformation in , represents -especially in its implementation- a nontrivial extension of a one-particle approximation to a many-body system. As noted in the introduction, the main advantage of the LDB many-polaron variational approach lies in the fact that the many-body effects in the system of charge carriers (electrons or holes) are completely contained in the structure factor of the electron (or hole) gas. This advantage will be carried through into the calculation of the optical properties of the interacting gas of continuum polarons which is the subject of the current paper.
### B Kubo formula for the optical conductivity of the many-polaron gas
The many-polaron optical conductivity is the response of the current-density, in the system described by the Hamiltonian (1), to an applied electric field (along the $`x`$-axis) with frequency $`\omega `$. This applied electric field introduces a perturbation term in the Hamiltonian (1), which couples the vector potential of the incident electromagnetic field to the current-density. As is well known, within linear response theory, the optical conductivity can be expressed through the Kubo formula as a current-current correlation function :
$$\sigma (\omega )=i\frac{Ne^2}{\text{V}m_b\omega }+\frac{1}{\text{V}\mathrm{}\omega }_0^{\mathrm{}}e^{i\omega t}[J_x(t),J_x(0)]𝑑t.$$
(6)
In this expression, V is the volume of the system, and $`J_x`$ is the $`x`$-component of the current operator $`𝐉,`$ which is related to the momentum operators of the charge carriers:
$$𝐉=\frac{q}{m_b}\underset{j=1}{\overset{N}{}}𝐩_j=\frac{q}{m_b}𝐏,$$
(7)
with $`q`$ the charge of the charge carriers ($`+e`$ for holes, $`e`$ for electrons) and $`𝐏`$ the total momentum operator of the charge carriers. The real part of the optical conductivity at temperature zero, which is proportional to the optical absorption coefficient, can be written as a function of the total momentum operator of the charge carriers as follows :
$$Re[\sigma (\omega )]=\frac{1}{\text{V}\mathrm{}\omega }\frac{e^2}{m_b^2}Re\left\{_0^{\mathrm{}}e^{i\omega t}[P_x(t),P_x(0)]𝑑t\right\}.$$
(8)
Integrating by parts twice, the real part of the optical conductivity of the many-polaron system can be written with a force-force correlation function:
$$Re[\sigma (\omega )]=\frac{1}{\text{V}\mathrm{}\omega ^3}\frac{e^2}{m_b^2}Re\left\{_0^{\mathrm{}}e^{i\omega t}[F_x(t),F_x(0)]𝑑t\right\},$$
(9)
with $`𝐅=(i/\mathrm{})[H_0,𝐏]`$. The commutator of the Hamiltonian (1) with the total momentum operator of the charge carriers simplifies to
$$𝐅=i\underset{𝐤}{}\underset{j=1}{\overset{N}{}}𝐤\left(e^{i𝐤.𝐫_j}a_𝐤V_𝐤+e^{i𝐤.𝐫_j}a_𝐤^+V_𝐤^{}\right).$$
(10)
This result for the force operator clarifies the significance of using the force-force correlation function rather than the momentum-momentum correlation function. The operator product $`F_x(t)F_x(0)`$ is proportional to $`|V_𝐤|^2`$, the charge carrier - phonon interaction strength. This will be a distinct advantage for any expansion of the final result in the charge carrier - phonon interaction strength, since one power of $`|V_𝐤|^2`$ is factored out beforehand. Denoting $`\rho _𝐤=_{j=1}^Ne^{i𝐤.𝐫_j}`$, the real part of the optical conductivity takes the form:
$$Re[\sigma (\omega )]=\frac{1}{\text{V}\mathrm{}\omega ^3}\frac{e^2}{m_b^2}Re\{\underset{𝐤,𝐤^{}}{}k_x.k_x^{}_0^{\mathrm{}}e^{i\omega t}\left[\begin{array}{c}e^{iH_0t/\mathrm{}}\left(\rho _𝐤a_𝐤V_𝐤+\rho _𝐤a_𝐤^+V_𝐤^{}\right)e^{iH_0t/\mathrm{}},\\ \left(\rho _𝐤^{}a_𝐤^{}V_𝐤^{}+\rho _𝐤^{}a_𝐤^{}^+V_𝐤^{}^{}\right)\end{array}\right]_0dt\}.$$
(11)
Up to this point, no approximations other than linear response theory have been made.
### C LDB canonical transformation for the optical conductivity
The expectation value appearing in the right hand side of expression (11) for the real part of the optical conductivity is calculated now with respect to the LDB many-polaron wave function $`|\psi _{\text{LDB}}`$ (2):
$`𝒥(𝐤,𝐤^{})`$ $`=`$ $`\psi _{\text{LDB}}\left|[e^{iH_0t/\mathrm{}}\left(\rho _𝐤a_𝐤V_𝐤+\rho _𝐤a_𝐤^+V_𝐤^{}\right)e^{iH_0t/\mathrm{}},\left(\rho _𝐤^{}a_𝐤^{}V_𝐤^{}+\rho _𝐤^{}a_𝐤^{}^+V_𝐤^{}^{}\right)]\right|\psi _{\text{LDB}}`$ (12)
$`=`$ $`\phi _{\text{el}}\left|\varphi \left|\left[\begin{array}{c}e^{iH^{}t/\mathrm{}}U^1\left(\rho _𝐤a_𝐤V_𝐤+\rho _𝐤a_𝐤^+V_𝐤^{}\right)Ue^{iH^{}t/\mathrm{}},\\ U^1\left(\rho _𝐤^{}a_𝐤^{}V_𝐤^{}+\rho _𝐤^{}a_𝐤^{}^+V_𝐤^{}^{}\right)U\end{array}\right]\right|\varphi \right|\phi _{\text{el}},`$ (15)
where $`U`$ is the many-polaron canonical transformation defined in (3) and $`H^{}=U^1H_0U`$ is the transformed Hamiltonian obtained in . The canonical transformation of the force term is
$$U^1\left(\rho _𝐤a_𝐤V_𝐤+\rho _𝐤a_𝐤^+V_𝐤^{}\right)U=\rho _𝐤\left(a_𝐤f_𝐤^{}\rho _𝐤\right)V_𝐤+\rho _𝐤\left(a_𝐤^+f_𝐤\rho _𝐤\right)V_𝐤^{}.$$
(16)
The terms of lowest order in the electron-phonon interaction amplitude $`|V_𝐤|^2`$ are given by
$`𝒥(𝐤,𝐤^{})`$ $`=`$ $`|V_𝐤|^2\delta _{\mathrm{𝐤𝐤}^{}}\phi _{\text{el}}\left|\varphi \left|e^{iH^{}t/\mathrm{}}\rho _𝐤a_𝐤e^{iH^{}t/\mathrm{}}\rho _𝐤a_𝐤^+\rho _𝐤a_𝐤e^{iH^{}t/\mathrm{}}\rho _𝐤a_𝐤^+e^{iH^{}t/\mathrm{}}\right|\varphi \right|\phi _{\text{el}}`$ (17)
$`=`$ $`2i|V_𝐤|^2\delta _{\mathrm{𝐤𝐤}^{}}Im\left[\phi _{\text{el}}\left|\varphi \left|e^{iH^{}t/\mathrm{}}\rho _𝐤a_𝐤e^{iH^{}t/\mathrm{}}\rho _𝐤a_𝐤^+\right|\varphi \right|\phi _{\text{el}}\right].`$ (18)
Taking the expectation value with respect to the phonon vacuum, we find
$$𝒥(𝐤,𝐤^{})=2i|V_𝐤|^2\delta _{\mathrm{𝐤𝐤}^{}}Im\left\{e^{i\omega _{\text{LO}}t}\phi _{\text{el}}\left|e^{iH^{}t/\mathrm{}}\rho _𝐤e^{iH^{}t/\mathrm{}}\rho _𝐤\right|\phi _{\text{el}}\right\}.$$
(19)
This can be substituted in the expression (11) for the real part of the optical conductivity, which becomes
$$Re[\sigma (\omega )]=2\frac{1}{\text{V}\mathrm{}\omega ^3}\frac{e^2}{m_b^2}Im\left\{\underset{𝐤}{}k_x^2|V_𝐤|^2_0^{\mathrm{}}e^{i\omega t}Im\left[e^{i\omega _{\text{LO}}t}\phi _{\text{el}}\left|e^{iH^{}t/\mathrm{}}\rho _𝐤e^{iH^{}t/\mathrm{}}\rho _𝐤\right|\phi _{\text{el}}\right]dt\right\}.$$
(20)
The right hand side of (20) can be written in a more compact form by introducing the dynamical structure factor of the electron (or hole) system.
### D General expression
To find the formula for the real part of the optical conductivity in its final form, we introduce the standard expression for the dynamical structure factor of the system of charge carriers interacting through a Coulomb potential,
$$S(𝐪,w)=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\phi _{\text{el}}\left|\frac{1}{2}\underset{j,\mathrm{}}{}e^{i𝐪.(𝐫_j(t)𝐫_{\mathrm{}}(0))}\right|\phi _{\text{el}}e^{iwt}𝑑t.$$
(21)
Rewriting expression (20) with the dynamical structure factor of the electron (or hole) gas results in:
$$Re[\sigma (\omega )]=\frac{n}{\mathrm{}\omega ^3}\frac{e^2}{m_b^2}\underset{𝐤}{}k_x^2|V_𝐤|^2S(𝐤,\omega \omega _{\text{LO}}),$$
(23)
where $`n=N/`$V is the density of charge carriers. As noted before, $`V_𝐤`$ is the electron-phonon interaction amplitude and $`k_x`$ is the $`x`$-component of the wave vector. Formula (23) for the optical absorption of the many-polaron system has an intuitively appealing form.
In the theory of one-polaron optical absorption for weak coupling constants $`\alpha `$, the optical absorption coefficient as obtained from Fermi’s golden rule is :
$$\text{1 polaron}:Re[\sigma (\omega )]\omega ^3\underset{𝐤}{}k_x^2|V_𝐤|^2\delta [\mathrm{}k^2/(2m_\text{b})(\omega \omega _{\text{LO}})].$$
(23)
At low densities, the dynamical structure factor $`S(q,\nu )`$ is strongly peaked around $`q^2/2=\nu `$ and is close to zero everywhere else . Substituting a delta-function $`\delta (q^2/2\nu )`$ for the dynamical structure factor in formula (23) it is easily seen that the one-polaron limit (23) is retrieved. The one-polaron result (23) is derived by considering a process in which the initial state consists of a photon of energy $`\mathrm{}\omega `$ and a polaron in its ground state, and the final state consists of an emitted LO phonon with energy $`\mathrm{}\omega _{\text{LO}}`$ and the polaron, scattered into a state with momentum $`𝐤`$ and kinetic energy $`(\mathrm{}k)^2/(2m_\text{b})=\mathrm{}(\omega \omega _{\text{LO}}).`$ The many-polaron result, formula (23), is a generalization of this one-polaron picture. The contribution which corresponds to the scattering of a polaron into the momentum state $`𝐤`$ and energy $`\mathrm{}(\omega \omega _{\text{LO}})`$ is now weighed by the dynamical structure factor $`S(k,\omega \omega _{\text{LO}})`$ of the electron (or hole) gas.
Formula (23) is reminiscent of the Hopfield formula describing the effect of impurities on the optical absorption of metals, which in turn is related to the expression obtained by Ron and Tzoar for the optical absorption in a quantum plasma. Formula (23) also represents a generalization of the results obtained by Gurevich, Lang and Firsov . These authors focused their attention on the many-body effects related to the Fermi exclusion statistics, whereas the present analysis will extend the results of to study the influence of plasmons and further many-body effects in the system of the constituent electrons or holes, as discussed in the next section.
The advantage of the LDB canonical transformation method for the evaluation of the ground state energy of a polaron gas is that the many-body effects are contained in the static structure factor of the electron (or hole) system, appearing in the analytical expression for the energy. The corresponding advantage of the canonical transformation in the present case, for the optical conductivity, is that many-body effects are again incorporated through a structure factor, now the dynamical structure factor of the electron (or hole) system. The level of approximation made in treating the many-body nature of the polaron system is determined by the choice of the dynamical structure factor of the electron or hole system.
In the present treatment of the electron-phonon interactions, the terms to leading order in $`|V_𝐤|^2`$ are automatically taken into account through the variational formulation based on LDB . As noted before, the use of the force-force correlation function allows to express the real part of the optical conductivity as $`Re[\sigma (\omega )]=\alpha (\omega ,\alpha ),`$ so that to lowest order in $`\alpha `$, $`Re[\sigma (\omega )]=\alpha (\omega ,\alpha =0)`$ where the electron-phonon interaction is no longer present in the factor $``$ which includes the many-body effects of the electron (or hole) system. A possible way to take into account higher-order terms in the electron-phonon interactions would be to include electron-phonon coupling effects at the level of the dynamical structure factor appearing in $``$, in a manner similar to Mahan’s treatment of the polaron spectral function , or to include multiple-phonon final states in the calculation .
### E Scaling relation for the optical absorption in two and three dimensions
The modulus squared of the Fröhlich electron-phonon interaction amplitude is given by
$$|V_𝐤|^2=\{\begin{array}{c}\frac{(\mathrm{}\omega _{\text{LO}})^2}{k^2}\frac{4\pi \alpha }{\text{V}}\sqrt{\frac{\mathrm{}}{2m_b\omega _{\text{LO}}}}\text{ in 3D}\hfill \\ \frac{(\mathrm{}\omega _{\text{LO}})^2}{k}\frac{2\pi \alpha }{\text{A}}\sqrt{\frac{\mathrm{}}{2m_b\omega _{\text{LO}}}}\text{in 2D,}\hfill \end{array}$$
(24)
where $`\alpha `$ is the (dimensionless) Fröhlich coupling constant determining the coupling strength between the charge carriers and the longitudinal optical phonons, and A is the surface of the 2D system . In what follows, we will use polaron units ($`\mathrm{}=m_b=\omega _{\text{LO}}=1`$). The sum over wave vectors in (23) can be written as an integral, so that for the three dimensional case (with dynamical structure factor $`S_{\text{3D}}`$) we find:
$$Re[\sigma _{\text{3D}}(\omega )]=ne^2\text{ }\frac{2}{3}\alpha \frac{1}{2\pi \omega ^3}\underset{0}{\overset{\mathrm{}}{}}𝑑q\text{ }q^2S_{\text{3D}}(q,\omega \omega _{\text{LO}}),$$
(25)
and for the two-dimensional case:
$$Re[\sigma _{\text{2D}}(\omega )]=ne^2\text{ }\frac{\pi }{2}\alpha \frac{1}{2\pi \omega ^3}\underset{0}{\overset{\mathrm{}}{}}𝑑q\text{ }q^2S_{2\text{D}}(q,\omega \omega _{\text{LO}}).$$
(26)
From these expressions, it is clear that the scaling relation
$$Re[\sigma _{\text{2D}}(\omega ,\alpha )]=Re[\sigma _{\text{3D}}(\omega ,3\pi \alpha /4)]$$
(27)
which holds for the one-polaron case introduced in ref. , is also valid for the many-polaron case if the corresponding 2D or 3D dynamical structure factor is used.
## III Results and discussion
### A General results
The expressions (25,26) allow us to derive results both for a three-dimensional and for a two dimensional polaron gas at $`T=0`$. The choice of a dynamical structure factor for the electron (or hole) system allows furthermore to study the different levels of approximation (Hartree-Fock, RPA,…) in the treatment of the many-electron or many-hole system. The results presented in this section were obtained using the material parameters of GaAs (for the two-dimensional case) and ZnO (for the three-dimensional case). These material parameters are summarized in Table I.
Fig. 1 shows the Hartree-Fock and the RPA result for the 2D many-polaron optical absorption spectrum (for GaAs, at a density $`n=10^{12}`$ cm<sup>-2</sup>). For reference, the dashed curve represents the familiar one-polaron result. In a first step, we discuss the result obtained by using the Hartree-Fock expression for the dynamical structure factor of the electron (or hole) system in the expressions (25,26).
The Fermi statistics cause the polarons to fill up a Fermi sphere up to $`k_\text{F}=[n/(2\pi )]^{1/2}`$. The optical absorption of the polaron gas resulting from this system is represented by the full curve labeled ‘Hartree-Fock’ in Fig. 1. The spectral weight at frequencies between $`\omega _{\text{LO}}`$ and $`1.4`$ $`\omega _{\text{LO}}`$ in Fig. 1 is reduced as compared to the single polaron case, whereas at higher frequencies it is enhanced. A kink appears in the spectrum at $`\omega =\omega _{\text{LO}}+E_\text{F}/\mathrm{}`$, as indicated by the dotted vertical line in Fig. 1.
This can be understood as follows. The absorption process is characterized by an initial state consisting of a polaron gas filling up the Fermi sphere (up to energy $`E_\text{F}`$, at $`T=0`$) and a photon with given energy $`\mathrm{}\omega `$, and by a final state made up of an emitted LO phonon with energy $`\mathrm{}\omega _{\text{LO}}`$ and a polaron gas such that one polaron state inside the Fermi sphere is not occupied and one polaron state with energy $`E>E_\text{F}`$ is occupied. The incident photon can only excite polarons out of the Fermi sea for which $`\mathrm{}\omega >h\omega _{\text{LO}}+E_\text{F}`$. A straightforward calculation in 2D shows that the fraction of polaron states in the Fermi sphere which can interact with a photon of energy $`\mathrm{}\omega `$ is given by
$$\{\begin{array}{c}0\text{ for }\omega <\omega _{\text{LO}}\hfill \\ \frac{\mathrm{}(\omega \omega _{\text{LO}})}{E_\text{F}}\text{ for }\omega _{\text{LO}}<\omega <E_\text{F}/\mathrm{}+\omega _{\text{LO}}\hfill \\ 1\text{ for }\omega >E_\text{F}/\mathrm{}+\omega _{\text{LO}}\hfill \end{array}.$$
(28)
For photon frequencies between $`\omega _{\text{LO}}`$ and $`\omega _{\text{LO}}+E_\text{F}/\mathrm{}`$, the number of polarons which cannot participate in the optical absorption process due to the Pauli exclusion principle, decreases linearly. At $`\omega =\omega _{\text{LO}}+E_\text{F}/\mathrm{}`$, all polarons can participate. This leads to a kink in the function (28) describing the number of polarons which can interact with the photon of given energy $`\mathrm{}\omega `$, as a function of $`\omega `$. This is also the origin of the kink in the optical absorption. This kink in the 2D many-polaron optical absorption spectrum at $`\mathrm{}\omega =\mathrm{}\omega _{\text{LO}}+E_\text{F}`$ was already noted in .
The full curve labeled with ‘RPA’ in Fig. 1 is obtained by using the random phase approximation (RPA) for the dynamical structure factor of the electron (or hole) system. It illustrates the combined effects of the Fermi statistics, discussed in the previous paragraph, and screening in the electron (or hole) system. In comparison to the Hartree-Fock curve, the main effect is an overall reduction of the spectral weight at frequencies $`\omega >\omega _{\text{LO}}`$. There is however a second effect, which is the appearance of a contribution related to plasmons - this is the subject of the next subsection.
### B Plasmon-phonon contribution
The RPA dynamical structure factor for the electron (or hole) system can be separated in two parts, one related to continuum excitations of the electrons (or holes) $`S_{\text{cont}}`$ , and one related to the undamped plasmon branch :
$`S_{\text{RPA}}(q,\omega )=A_{\text{pl}}(q)\delta [\omega \omega _{\text{pl}}(q)]+S_{\text{cont}}(q,\omega ),`$
where $`\omega _{\text{pl}}(q)`$ is the wave number dependent plasmon frequency and $`A_{\text{pl}}`$ is the strength of the undamped plasmon branch . The insets of Fig. 2 depict the regions in the $`q`$-$`v`$ plane ($`q=k/k_F,v=m_\text{b}\omega /(\mathrm{}k_F^2)`$) where the RPA dynamical structure factor is different from zero. The contribution (after substitution of $`S_{\text{RPA}}`$ in (25)-(26)) in the many-polaron optical absorption) deriving from the undamped plasmon branch $`A_{\text{pl}}(q)\delta [\omega \omega _{\text{pl}}(q)]`$ will be denoted as ‘plasmon-phonon’ contribution. The physical process related to this contribution is the emission of both a phonon and a plasmon in the scattering process.
Figure 2 shows the result for the optical absorption of the many-polaron gas for the 2D case (GaAs, left panel) and the 3D case (ZnO, right panel). For reference, the dashed curves show the one-polaron result. The full curves show the many-polaron results in the random-phase approximation. The shaded gray areas indicate the plasmon-phonon contribution.
Now examine the 3D case (the right panel of Fig. 2). The frequency of the undamped plasmon mode lies between $`\omega _1`$ and $`\omega _\text{2}`$ where $`\omega _\text{1}=\omega _{\text{pl}}=\sqrt{4\pi ne^2/m_\text{b}}`$ is the frequency of the plasmon branch at $`q=0`$, and $`\omega _2`$ is the frequency at which the branch of the undamped plasmons enters the Landau damping region (whose edge is given by $`\omega =\mathrm{}q^2/(2m_b)+\mathrm{}k_\text{F}q/(2m_b)`$). The corresponding plasmon-phonon contribution to the optical absorption ‘starts’ at $`\omega _{\text{LO}}+\omega _1`$ and ‘ends’ at $`\omega _{\text{LO}}+\omega _2`$. These frequencies are indicated by vertical dotted lines in the right panel of Fig. 2.
In the 2D case, the undamped plasmon branch is acoustic-like; for $`q0`$, $`\omega _{\text{pl}}0`$. Consequently, the phonon-plasmon peak in this case extends from $`\omega _{\text{LO}}`$ up to $`\omega _{\text{LO}}+\omega _2`$ where $`\omega _2`$ is the frequency at which the undamped plasmon branch enters the region of the continuum excitations of the 2D (RPA) electron gas.
In Fig. 3, the evolution of the many-polaron optical absorption spectrum is shown as the density of electrons (or holes) is increased. Two effects can be observed for increasing density: the reduction of the optical absorption above $`\omega >\omega _{\text{LO}}`$ and the shift towards higher frequencies of the plasmon-phonon contribution, both in 2D (left panel) and in 3D (right panel). The f-sum rule is nevertheless satisfied due to the presence of a central $`\delta (\omega )`$ peak in the optical absorption of the polaron gas at $`T=0`$.
### C Comparison with other theories
In earlier work, Wu, Peeters and Devreese studied the influence of screening on the electron-phonon interaction in a two-dimensional electron gas based on a memory function approach using a perturbation expansion in the electron-LO phonon coupling constant. The results of this perturbative approach of for the optical conductivity in 2D, also in the RPA framework, are consistent with the results derived from the present method based on the variational LDB unitary transformation. These authors found an enhancement of the optical absorption at the frequency where the undamped plasmon branch reaches the region of continuum excitations of the electron gas. The present, variational, method extends these results by taking into account the entire undamped plasmon branch.
Recently, Cataudella, De Filippis and Iadonisi investigated the optical properties of the many-polaron gas by calculating the correction due to electron-phonon interactions to the RPA dielectric function of the electron gas, starting from the Feynman polaron model and ref. . An aspect of the present method is that it is not restricted to the random-phase approximation for the treatment of the many-body effects between the charge carriers. Cataudella et al. also find a suppression of the optical absorption with increasing density. To our knowledge, the plasmon-phonon contribution was not revealed by the work of Cataudella et al. .
### D Comparison to the infrared spectrum of Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>2-y</sub>
Calvani and collaborators have performed doping-dependent measurements of the infrared absorption spectra of the high-T<sub>c</sub> material Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>2-y</sub> (NCCO). The region of the spectrum examined by these authors (50-10000 cm<sup>-1</sup>) is very rich in absorption features: they observe is a “Drude-like” component at the lowest frequencies, and a set of sharp absorption peaks related to phonons and infrared active modes (IRAV, up to about 1000 cm<sup>-1</sup>) possibly associated to small (Holstein) polarons . Three distinct absorption bands can be distinguished: the ‘d-band’ (around 1000 cm<sup>-1</sup>), the Mid-Infrared band (MIR, around 5000 cm<sup>-1</sup>) and the Charge-Transfer band (CT, around 10<sup>4</sup> cm<sup>-1</sup>) . Of all these features, the d-band and, at a higher temperatures, the Drude-like component have (hypothetically) been associated with large polaron optical absorption .
For the lowest levels of Ce doping, the d-band can be most clearly distinguished from the other features. The experimental optical absorption spectrum (up to 3000 cm<sup>-1</sup>) of Nd<sub>2</sub>CuO<sub>2-δ</sub> ($`\delta <0.004`$), obtained by Calvani and co-workers , is shown in Fig. 4 (shaded area) together with the theoretical curve obtained by the present method (full, bold curve) and, for reference, the one-polaron optical absorption result (dotted curve). At lower frequencies (600-1000 cm<sup>-1</sup>) a marked difference between the single polaron optical absorption and the many-polaron result is manifest. The experimental d-band can be clearly identified, rising in intensity at about 600 cm<sup>-1</sup>, peaking around 1000 cm<sup>-1</sup>, and then decreasing in intensity above that frequency. At a density of $`n=1.5`$ $`10^{17}`$ cm<sup>-3</sup>, we found a remarkable agreement between our theoretical predictions and the experimental curve. The experimentally determined material parameters used in the present calculation are summarized in Table I. A background contribution, taken to be constant over the frequency range of the d-band, was substracted in Figure 4. The lack of experimental data on several material constants leaves us with three adjustable parameters: the electron-phonon coupling constant $`\alpha ,`$ the band mass $`m_\text{b}`$, and the density of charge carriers. These parameters were chosen as follows:
* A set of theoretical optical absorption spectra were generated for different values of the band mass and densities (for $`m_\text{b}=0.1,0.2,0.5,0.8,1.0,2.0`$ and $`n=\{0.1,0.2,0.5,1.0,1.2,1.5,2.0\}\times 10^{17}`$ cm<sup>-3</sup>).
* For each of those spectra, the coupling constant $`\alpha `$ was chosen so as to fit the tail region of the experimental optical absorption spectrum best, using a least squares fitting procedure (the tail region is relatively insensitive to many-polaron effects).
* The best fitting curve, using again a least squares evaluation of the goodness of fit, was selected; we found fair agreement with $`m_\text{b}=0.5`$ $`m_\text{e}`$ and $`n=1.5`$ $`10^{17}`$ cm<sup>-3</sup> at $`\alpha =2.1.`$
This comparison with experiment could not be performed at higher doping content: the frequency region of the d-band usually contains a strong non-uniform contribution of other optical absorption features (such as the onset of the MIR, the tail of the Drude contribution, and the IRAV and phonon modes). To take these other contributions into account, additional adjustable parameters would have to be introduced making a comparison less convincing. Fortunately, experimental results are available in the form of the normalized first frequency moment of the optical absorption spectrum (after substraction of MIR and CT band) $`Re[\sigma _{\text{exp}}(\omega )]`$ :
$$\omega =\frac{{\displaystyle \underset{0}{\overset{\omega _{\text{max}}}{}}}\omega Re[\sigma _{\text{exp}}(\omega )]d\omega }{{\displaystyle \underset{0}{\overset{\omega _{\text{max}}}{}}}Re[\sigma _{\text{exp}}(\omega )]d\omega }$$
(29)
where $`\omega _{\text{max}}=10000`$ cm<sup>-1</sup> . Calvani and co-workers determined $`\omega `$ for NCCO samples with a varying cerium doping content. Increasing the cerium doping will inject electrons in the copper-oxide planes of the material, and increase the 2D charge carrier density in these planes. A comparison of this experimental normalized first frequency moment to the theoretical one presents the advantage that fewer parameters need to be adapted: only the density and the electron band mass have to be taken from experiment or (if experimental values are lacking) fitted.
The carrier density can be estimated numerically from the effective carrier concentrations in the different samples and from a measurement of the two-dimensional Fermi velocity performed for one of the samples . As for the other cuprates, the band mass of the electrons in NCCO has not yet been determined experimentally , and remains as an adjustable parameter.
Fig. 5 represents the comparison between the present theory and experiment. The squares with error bars show the experimental results for differently doped samples of NCCO, reported in . The dashed curve shows the normalized first frequency moment of the theoretical optical absorption spectrum, integrated over the entire frequency range ($`\omega _{\text{max}}\mathrm{}`$). The tail region of the many-polaron optical absorption still carries a significant weight, just as it does in the one-polaron optical absorption. It is necessary to include the cutoff frequency. The full curve represents the theoretical first frequency moment with a cutoff frequency $`\omega _{\text{max}}=10000`$ cm<sup>-1</sup>, which corresponds to the experimental cutoff .
There exists a fair agreement between the theoretical and the experimental values of the normalized first frequency moment for the five samples with lowest density, which have a cerium doping content of $`x<0.12`$. These correspond to the squares to the left of the dotted vertical line in Fig. 5. For the four remaining samples $`(x>0.12)`$ a discrepancy between the theoretically predicted first frequency moment for unpaired polarons and the observed first frequency moment appears. It has been observed experimentally that the weight of the low-frequency component in these samples (with $`x>0.12`$) is significantly larger that the corresponding weight in samples with $`x<0.12`$ . This was interpreted in as a consequence of an insulator-to-metal transition taking place around the cerium-doping level of $`x=0.12`$. Therefore, it seems reasonable to assume that above this doping level $`x`$ a change in the nature of the charge carriers takes place. One could hypothesize that, as the formation of bipolarons is stabilized with increasing density of the polaron gas , bipolarons start playing a role in the optical absorption spectrum. In a variety of other cuprates and manganates, the presence of bipolarons has also been invoked to interpret a number of response-related properties .
## IV Conclusions
Starting from the many-polaron canonical transformations and the variational many-polaron wave function (LDB) introduced in we have derived a formula for the optical absorption coefficient $`Re[\sigma (\omega )]`$ of a many-polaron gas. We find that $`Re[\sigma (\omega )]`$ can be expressed in a closed analytical form in terms of the dynamical structure factor $`S(q,\omega )`$ of the electron (or hole) system, equation (26) in 2D and (25) in 3D. In the present approach, the electron-phonon coupling and the electron-electron many-body effects formally decouple in the expression for $`Re[\sigma (\omega )]`$. Therefore, the many-body effects in the electron (hole) system can be taken into account by employing any desired approximation to the dielectric response (Hartree-Fock, RPA, etc.) of this electron (hole) system.
In the present work, the dynamical structure factor $`S(q,\omega )`$ of the electron (or hole) gas was considered both in the Hartree-Fock and the RPA approximation.
The main effect of the Pauli exclusion principle on the optical absorption of the polaron gas turns out to be a shift of the oscillator strength towards higher frequencies. This effect can be understood in terms of the available initial and final states in the polaron-photon scattering process and naturally invokes the Fermi energy $`E_\text{F}`$ of the electron (hole) gas.
The main effects in the case of the RPA approximation are an overall reduction of the optical absorption at frequencies $`\omega >\omega _{\text{LO}}`$ and the introduction of a novel absorption feature which we identified as a plasmon-phonon peak. This plasmon-phonon peak shifts to higher frequencies with increasing density, such that a double peak structure can appear in the 3D many-polaron optical absorption spectrum, consistent with the observed bimodal polaronic band in cadmium oxide .
As a first application of the method presented here, we chose to investigate the optical absorption of the interacting polaron gas in the RPA framework. For Nd<sub>2</sub>CuO<sub>2-δ</sub> ($`\delta <0.004`$), similarities were observed (see Fig. 4) between the line shape of the experimental d-band and the many-polaron optical absorption as calculated here. To study the density dependence, measurements (performed by Calvani and co-workers for a family of Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4-y</sub> materials) of the first frequency moment of the optical absorption were compared to the results of the present theory. We find a fair agreement for the samples with the lowest densities (cerium doping $`x<0.12`$). A softening of the first frequency moment of the optical absorption for the samples with higher densities (cerium doping $`x>0.12`$) is consistent with a change in the nature of the charge carriers at a doping content $`x=0.12`$ inferred in ref. from infrared absorption experiments.
## Acknowledgments
The authors like to acknowledge S. N. Klimin and V. M. Fomin for helpful discussions and intensive interactions. We are indebted to P. Calvani for fruitful discussions and for communication of experimental data. We thank F. Brosens and L. F. Lemmens for discussions. One of us, J.T., (“Postdoctoraal Onderzoeker van het Fonds voor Wetenschappelijk Onderzoek – Vlaanderen”), is supported financially by the Fonds voor Wetenschappelijk Onderzoek – Vlaanderen (Fund for Scientific Research – Flanders). Part of this work is performed in the framework of the “Interuniversity Poles of Attraction Program – Belgian State, Prime Minister’s Office – Federal Office for Scientific, Technical and Cultural Affairs” (“Interuniversitaire Attractiepolen – Belgische Staat, Diensten van de Eerste Minister – Wetenschappelijke, Technische en Culturele Aangelegenheden”), and in the framework of the FWO projects 1.5.545.98, G.0287.95, 9.0193.97, WO.025.99N and WO.073.94N (Wetenschappelijke Onderzoeksgemeenschap, Scientific Research Community of the FWO on “Low Dimensional Systems”), and in the framework of the BOF NOI 1997 and GOA BOF UA 2000 projects of the Universiteit Antwerpen.
## Table
Table I: Material parameters used in the various figures. The physical parameters for GaAs correspond to those of the GaAs-AlGaAs heterostructure ; the material parameters for ZnO are taken from . The physical parameters for the neodymium-cerium cuprate are taken from and . “n.a.” (“not applicable”) means that not enough data are available to estimate this material parameter.
material parameters: GaAs ZnO Nd<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>2</sub> phonon frequency $`\mathrm{}\omega _{\text{LO}}=`$ $`36.77`$ meV $`73.27`$ meV $`74`$ meV dielectric constants $`\epsilon _0=`$ $`12.83`$ $`8.15`$ n.a. $`\epsilon _{\mathrm{}}=`$ $`10.9`$ $`4.00`$ ca. $`3.`$ band mass $`m_\text{b}=`$ $`0.0657`$ $`m_\text{e}`$ $`0.24`$ $`m_\text{e}`$ n.a. coupling constant $`\alpha =`$ $`0.068`$ $`0.849`$ n.a. polaron length unit $`a_{\text{HO}}=`$ $`5.616`$ nm $`2.082`$ nm n.a. bohr radius $`a_\text{B}=`$ $`8.7797`$ nm $`0.882`$ nm n.a.
## V Figure captions
Figure 1 : The real part of the optical conductivity (proportional to the optical absorption coefficient) of an interacting large-polaron gas is shown as a function of frequency, for a two dimensional gas (GaAs) from equation 26. The material parameters are given in Table I. The dashed curve represents the one-polaron result, the full curve labeled ‘Hartree-Fock structure factor’ shows the result using the Hartree-Fock approximation to the dynamical structure factor of the electron (hole) system, and the full curve labeled ‘RPA structure factor’ is the result in the Random Phase approximation. The dotted vertical line indicates the threshold frequency above which all polarons can be scattered into unoccupied final states and participate in the absorption process. The broad gray peak in the RPA curve is the plasmon-phonon contribution (see also figure 2).
Figure 2 : The real part of the optical conductivity is shown as a function of frequency for an interacting large-polaron gas in the 2D case (left panel) and the 3D case (right panel). The material parameters used here are given in Table I. The dashed curves represent the single polaron spectra. The full curve represents the many-polaron spectrum. In this figure, the plasmon-phonon contribution to the optical many-polaron spectrum is shown as a shaded area. This contribution arises from a process where a polaron, with the absorption of a photon, emits a phonon and a plasmon. The inset shows the regions in the $`q`$-$`v`$ plane where the dynamical structure factor $`S(q,v)`$ of the electron (or hole) system used in the optical absorption formulas (25,26) differs from zero; the Landau damping region and the undamped plasmon branch can be distinguished.
Figure 3 : The real part of the optical conductivity is shown as a function of frequency for different densities of an interacting large-polaron gas, in the 2D case (left panel) and the 3D case (right panel). The material parameters used for this figure are given in Table I. For increasing density, the optical conductivity is reduced. Another effect in the RPA approximation is the presence of a peak related to the undamped plasmon branch (see figure 2) which shifts according to the plasma frequency.
Figure 4 : The infrared absorption of Nd<sub>2</sub>CuO<sub>2-δ</sub> ($`\delta <0.004`$) is shown as a function of frequency, up to 3000 cm<sup>-1</sup>. The experimental results of Calvani and co-workers is represented by the thin black curve and by the shaded area. The so-called ‘d-band’ rises in intensity around 600 cm<sup>-1</sup> and increases in intensity up to a maximum around 1000 cm<sup>-1</sup>. The dotted curve shows the single polaron result. The full black curve represents the theoretical results obtained in the present work for the interacting many-polaron gas with $`n=1.5`$ $`10^{17}`$ cm<sup>-3</sup>, $`\alpha =2.1`$ and $`m_\text{b}=0.5`$ $`m_\text{e}`$.
Figure 5 : The normalized first frequency moment of the optical absorption spectra is shown as a function of the density (expressed through the Fermi wave vector). The squares represent the experimental results of Calvani and co-workers in a family of Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub> materials. The dashed curve shows the results from the theoretical two-dimensional many-polaron optical absorption, obtained by integrating all frequencies in the calculation of the first frequency moment. The full curve shows the theoretical results obtained by integrating up to a cut-off frequency, which is chosen at $`10000`$ cm<sup>-1</sup> and which corresponds to the maximum frequency in the experiment . The material parameters are listed in the figure, and the effect of choosing a different electron band mass is illustrated in the inset. The points with $`x<0.12`$, to the left of the vertical dotted line, show agreement with the theoretical result from the many-polaron theoretical optical absorption, but it is clear that the experimental cutoff frequency has to be taken into account. For the samples with $`x>0.12`$, a discrepancy between the theoretically predicted first frequency moment and the observed first frequency moment is consistent with a possible insulator-to-metal transition at $`x=0.12`$ .
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# Cosmological Models from Quintessence
## I Allowed potentials for quintessence
Quintessence has been recently invoked as an advantageous alternative to the cosmological constant in order to explain the apparent accelerating expansion of the universe which has stirred cosmologists after observations and measurements of distant supernovae \[2-4\]. The existence of a quintessential field has been related to supersymmetric models , the problem of fine-tuning of the cosmological constant , or supergravity models . The bare standard cosmological model (BSCM), without any constant cosmological term or vacuum fields, predicts the existence of an expanding universe which can be closed, open or flat, but always decelerating. However, if a positive cosmological constant is added to the field equations, then the expansion of the universe may become accelerating. Actually, as early as 1975, Gunn and Tinsley, while discussing observations on the Hubble diagram and constraints on the matter density of the universe and ages of galaxies, found a series of allowable, if not compelling, cosmological models with nonzero, positive cosmological constant, which were accelerating. Recent observations on distant supernovas have resurrected the spirit of these early conclusions and led to the strong suspect that, in spite of the fact that the BSCM gives satisfactory explanations to many other observations, it probably is incomplete or even incorrect .
When the cosmological constant, $`\mathrm{\Lambda }`$, is interpreted as the energy density of vacuum for an equation of state $`p=\rho `$, if it is positive, then most inflationary models are suitably pinpointed. However, it is largely known that $`\mathrm{\Lambda }`$ is not free of fundamental problems . Actually, the so-called cosmological constant problem is one of the most challenging questions in fundamental physics, as it is very hard to envisage any consistent mechanism that dynamically explains how the vacuum energy density can be lowered from its most natural value at around the Planck scale down to its observationally allowed value, $`ϵ10^{47}`$ GeV<sup>4</sup>. Although quintessence models do not solve this problem, they may improve the related fine-tuning problem in the sense that they can explain a tiny value for the vacuum energy density with a scale comparable with the scales of high energy physics. Besides, these models give rise to an accelerating universe by using a vacuum dynamically adjustable, time-dependent scalar field that is spatially (in-)homogeneous and evolves slowly enough so that the kinetic term of the energy density is always smaller than the potential energy term. It is worth noticing that this is not necessarily required in tracker models of quintessence (see e.g Ref. ). Indeed, in the case of an overshoot the kinetic energy dominates at high redshift. If we disregard such tracker models, the resulting negative pressure will then correspond to an equation of state $`p=\omega \rho `$ where the free-parameter $`\omega `$ can take on any values $`0\omega >1`$. Thus, the cosmological constant will correspond to the extreme case $`\omega =1`$.
Recently, however, di Prieto and Demaret have shown that, if we restrict ourselves to a constant equation of state none of the vacuum scalar-field potentials, $`V(\varphi )`$, currently used in quintessence models, such as the exponential , cosine and inverse power-law potentials can satisfy the constraint on $`V(\varphi )`$ implied by field equations and conservation laws, i.e.
$`{\displaystyle \frac{V^{}}{V_0^{}}}=`$
$$\pm \left[\mathrm{\Omega }_\varphi \left(\frac{V}{V_0}\right)^2+\mathrm{\Omega }_M\left(\frac{V}{V_0}\right)^{\frac{\omega +2}{\omega +1}}+\mathrm{\Omega }_k\left(\frac{V}{V_0}\right)^{\frac{3\omega +5}{3(\omega +1)}}\right]^{\frac{1}{2}},$$
(1.1)
where $`V_0`$ and $`V_0^{}`$ are the current values of the scalar-field potential, $`V(\varphi )`$, and its derivative with respect to the field, $`V^{}=dV(\varphi )/d\varphi `$, respectively, and the $`\mathrm{\Omega }_i`$’s (with $`i=\varphi ,M,k`$) are the dimensionless density parameters for the scalar field, ordinary matter and topological curvature. Apart from a solution for any $`\omega `$ in the flat case, Pietro and Demaret were nonetheless able to find some solutions to constraint (1.1) for particular values of the parameter $`\omega `$. Thus, for $`\omega =1/3`$, they obtained
$$V(\varphi )=V_0\left\{\frac{\sqrt{2}}{4ϵ_0}\mathrm{sinh}\left[\pm ϵ_0(\varphi \varphi _0)+\nu _0\right]\right\}^4,$$
(1.2)
where
$`ϵ_0={\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\mathrm{\Omega }_\varphi +\mathrm{\Omega }_k}{2\mathrm{\Omega }_M}}},\nu _0=\mathrm{arcsin}\left(2\sqrt{2}ϵ_0\right),`$
and, for $`\omega =2/3`$,
$`V(\varphi )=V_0\{{\displaystyle \frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_\varphi }}\mathrm{sinh}[\pm (\varphi \varphi _0)+\delta _0]`$
$$+\frac{\mathrm{\Omega }_k}{4\mathrm{\Omega }_\varphi }(\frac{\mathrm{\Omega }_k}{\mathrm{\Omega }_M}e^{(\varphi \varphi _0)\sigma _0}2)\}^{\frac{1}{2}},$$
(1.3)
where
$`e^{\sigma _0}={\displaystyle \frac{2\mathrm{\Omega }_\varphi +2\sqrt{\mathrm{\Omega }_\varphi }+\mathrm{\Omega }_k}{\mathrm{\Omega }_M}}`$
$`\delta _0=\sigma _0+{\displaystyle \frac{1}{2}}\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }_M}{4\mathrm{\Omega }_\varphi }}\right).`$
We furthermore note that, besides solutions (1.2) and (1.3), there are a whole family of scalar potentials $`V(\varphi )`$ defined in terms of the Jacobian elliptic functions , $`J_e`$, which satisfy the constraint (1.1) for $`\omega =1/6`$. Such solution can be generally written as
$$V(\varphi )=V_0\left\{J_e[\alpha _0(\varphi \varphi _0),m]\right\}^{10},$$
(1.4)
in which $`\alpha _0=\alpha _0(\mathrm{\Omega }_i)`$ is a given constant whose form depends on the particular elliptic function being considered, and $`m=m(\mathrm{\Omega }_i)1`$ is the characteristic parameter (modulus) of the corresponding elliptic function. For example, taking for $`J_e`$ the function cn, we have
$`\alpha _0=\sqrt{{\displaystyle \frac{7\mathrm{\Omega }_k(\mathrm{\Omega }_\varphi +2\mathrm{\Omega }_k)}{200\mathrm{\Omega }_\varphi (\mathrm{\Omega }_k+\mathrm{\Omega }_\varphi )}}},`$
$`m={\displaystyle \frac{\mathrm{\Omega }_k}{\mathrm{\Omega }_\varphi +2\mathrm{\Omega }_k}},`$
or for the function sd
$`\alpha _0=\sqrt{{\displaystyle \frac{7\mathrm{\Omega }_M}{200\mathrm{\Omega }_\varphi }}},`$
$`m={\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{2\mathrm{\Omega }_\varphi ^2}{7}}\right),`$
and similar, but distinct expressions of $`\alpha _0`$ and $`m`$ for the remaining 10 Jacobian elliptic functions.
It appears clearly of interest to investigate whether the new potentials (1.2)-(1.4) are actually able to predict accelerating cosmological models which can match recently obtained data from observations on distant supernova Ia, discussing their physical relevance as well. It is the aim of this paper to carry out such an investigation, incorporating other possible new solutions from a generalized quintessence model which simultaneously accommodates both a vacuum scalar field $`\varphi `$ and a varying cosmological term $`\mathrm{\Lambda }`$. In this paper, we shall restrict ourselves to equations of state with a constant $`\omega `$, both in conventional and generalized quintessence models, disregarding tracker models , where time varying equations of state are invoked and a general inverse-power law potential for the quintessence field is assumed.
## II Generalized quintessence model
The field equations corresponding to a Friedmann-Robertson-Walker spacetime with ordinary matter which is not coupled to a homogeneous (quintessence) scalar field $`\varphi `$, to which we add a varying cosmological term $`\mathrm{\Lambda }`$ can be written as
$$\frac{\dot{R}^2}{R^2}+\frac{k}{R^2}=\frac{1}{3}\kappa ^2\left(\rho _\varphi +\rho _M+\rho _\mathrm{\Lambda }\right)$$
(2.1)
$$2\frac{\ddot{R}}{R}+\frac{\dot{R}^2}{R^2}+\frac{k}{R^2}=\kappa ^2\left(p_\mathrm{\Lambda }p_\varphi \right)$$
(2.2)
$$\ddot{\varphi }+3\dot{\varphi }\frac{\dot{R}}{R}=V^{},$$
(2.3)
where the overhead dot means time derivative, $`{}_{}{}^{}d/d\varphi `$, $`\kappa ^2=8\pi G_N`$, $`k`$ is the topological curvature and we have defined the scalar field such that
$$\kappa ^2\rho _\varphi =\frac{1}{2}\dot{\varphi }^2+V(\varphi )$$
(2.4)
$$\kappa ^2p_\varphi =\frac{1}{2}\dot{\varphi }^2V(\varphi ).$$
(2.5)
As usual, the ordinary matter is assumed to obey the equation of state for an ordinary fluid, $`p_M=0`$. As pointed out before, the scalar quintessence field will be assumed to behave like a perfect fluid with equation of state
$$p_\varphi =\omega \rho _\varphi ,1<\omega 0.$$
(2.6)
The generalization implied by the quintessence field with respect to the case of a pure cosmological constant can be manisfested by noting that the particular value of the constant parameter $`\omega =1`$ corresponds to the cosmological constant case when the field $`\varphi `$ becomes a constant as well .
The conservation laws that the involved fields are here assumed to satisfy are as follows. First of all, we note that, since there is no interaction between the scalar field and the other fields involved, we can take all these laws as separable from each other. For ordinary matter, $`M`$, and scalar field, $`\varphi `$, we should then have for all values of $`\omega `$, except $`\omega =1`$
$$\rho _M=\rho _{M0}\left(\frac{R_0}{R}\right)^3,\rho _\varphi =\rho _{\varphi 0}\left(\frac{R_0}{R}\right)^{3(1+\omega )},$$
(2.7)
with the subscript $`0`$ taken to always mean current value. As to the varying cosmological term $`\mathrm{\Lambda }`$, we generally assume $`\kappa ^2\rho _\mathrm{\Lambda }=\mathrm{\Lambda }=\mathrm{\Lambda }_0(R_0/R)^n`$, where $`n`$ can, in principle, take on the values 1, 2 and 3. However, only the value $`n=1`$ corresponds to a model with additional dynamical content relative to the constant-$`\omega `$ usual models, the cases $`n=2`$ and $`n=3`$ just reducing to the Pietro-Demaret model with the cosmological dimensionless parameter (see below) $`\mathrm{\Omega }_k`$ replaced for $`\mathrm{\Omega }_k+\mathrm{\Omega }_\mathrm{\Lambda }`$ and $`\mathrm{\Omega }_M`$ replaced for $`\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }`$, respectively. We then take for the most general conservation law for $`\mathrm{\Lambda }`$
$$\kappa ^2\rho _\mathrm{\Lambda }\mathrm{\Lambda }=\mathrm{\Lambda }_0\left(\frac{R_0}{R}\right).$$
(2.8)
In any event, however, $`\mathrm{\Lambda }`$ can be taken to represent the energy density of the field $`\varphi `$ corresponding to a particular value of parameter $`\omega `$ ($`\omega =2/3`$ in the chosen conservation law (2.8)), so that if we would allow both signs for $`\mathrm{\Lambda }`$ and $`\rho _\varphi `$ then keeping simultaneously $`\mathrm{\Lambda }`$ and $`\rho _\varphi `$ in the field equations would just be redundant. Nevertheless, one still can consistently consider field equations with $`\mathrm{\Lambda }`$ and $`\rho _\varphi `$ simultaneously provided that we restrict their values to be either (i) $`\mathrm{\Lambda }>0`$, $`\rho _\varphi <0`$, or (ii) $`\mathrm{\Lambda }<0`$, $`\rho _\varphi >0`$. In what follows, we shall confine ourselves to just case (i), looking at the quantity $`\rho _v\rho _\mathrm{\Lambda }+\rho _\varphi `$ as the total vacuum energy which will always be taken to be $`\rho _v0`$. Bearing in mind such a restriction, we choose Eqs. (2.7) and (2.8) as our conservation laws and introduce then the following dimensionless cosmological parameters
$$\mathrm{\Omega }_k=\frac{k}{R_0^2H_0^2},\mathrm{\Omega }_\varphi =\frac{\kappa ^2\rho _{\varphi 0}}{3H_0^2},\mathrm{\Omega }_M=\frac{\kappa ^2\rho _{M0}}{3H_0^2}$$
(2.9)
$$\mathrm{\Omega }_\mathrm{\Lambda }=\frac{\mathrm{\Lambda }_0}{3H_0^2},$$
(2.10)
which should satisfy the quadrilateral constraint
$$\mathrm{\Omega }_0=1\mathrm{\Omega }_k=\mathrm{\Omega }_M+\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda },$$
(2.11)
rather than the constraint implied by the usual cosmological triangle. From the conservation laws and the first field equation, we obtain a differential constraint on the scale factor
$`\dot{R}^2=R_0^2H_0^2[\mathrm{\Omega }_M{\displaystyle \frac{R_0}{R}}`$
$$+\mathrm{\Omega }_k+\mathrm{\Omega }_\varphi \left(\frac{R_0}{R}\right)^{1+3\omega }+\mathrm{\Omega }_\mathrm{\Lambda }\frac{R}{R_0}],$$
(2.12)
and from the relations between $`R`$ and $`V`$, and $`\dot{R}`$ and $`V^{}`$, derived by di Prieto and Demaret , the following generalized constraint on the scalar quintessence potential
$`\left({\displaystyle \frac{V^{}}{V_0^{}}}\right)^2=\mathrm{\Omega }_M\left({\displaystyle \frac{V}{V_0}}\right)^{\frac{\omega +2}{\omega +1}}`$
$$+\mathrm{\Omega }_\varphi \left(\frac{V}{V_0}\right)^2+\mathrm{\Omega }_k\left(\frac{V}{V_0}\right)^{\frac{3\omega +5}{3(\omega +1)}}+\mathrm{\Omega }_\mathrm{\Lambda }\left(\frac{V}{V_0}\right)^{\frac{3\omega +4}{3(\omega +1)}},$$
(2.13)
which has been obtained assuming that $`\omega 1`$, with
$$V_0=\frac{3}{2}(1\omega )H_0^2\mathrm{\Omega }_\varphi ,V_0^{}=\pm H_0\sqrt{\frac{1}{2}\left(1\omega ^2\right)V_0}.$$
(2.14)
Constraint (2.13) is of course a generalization from constraint (1.1) and reduces to this when we set $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$.
Although again the exponential and cosine potentials cannot satisfy the constraint (2.13) even if we relax the condition of nonclosedness implied by nucleosynthesis and supernova observations, we note that there exist some inverse-power law potentials, similar to those used in the literature , which satisfy our field equations and conservation laws. Thus, if we set $`\mathrm{\Omega }_\varphi +\mathrm{\Omega }_k=\mathrm{\Omega }_\mathrm{\Lambda }=0`$, and $`\mathrm{\Omega }_M=1`$, we have as a solution to constraint (2.13)
$`V=V_0\left({\displaystyle \frac{\varphi _0}{\varphi }}\right)^4,`$
for $`\omega =1/3`$. On the other hand, setting $`\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }=\mathrm{\Omega }_M=0`$, we obtain another solution
$`V=V_0\left({\displaystyle \frac{\varphi _0}{\varphi }}\right)^2,`$
for $`\omega =2/3`$. Even though they correspond to particular constant values of $`\omega `$, these two potentials could still be implemented in the realm of high energy physics. They have the same general form as the Ratra-Peebles potential , though this does not actually require that $`\omega `$ be assumed constant. Indeed, the above two potentials may be regarded to belong to a potential family $`V=V_0(\varphi _0/\varphi )^{6(1+\omega )}`$ which can be related to the Ratra-Peebles potential by the field transformation $`\varphi _{RP}\alpha \varphi ^{13\omega }`$, with $`\alpha `$ a suitable dimensional constant. In addition, there are other allowable potentials for the field $`\varphi `$ which are solution to Eq. (2.13) for particular values of the quintessence parameter $`\omega `$, without imposing any restriction on the cosmological parameters $`\mathrm{\Omega }_j`$. Thus, for $`\omega =1/3`$, we get a family of generalized quintessence potentials which are given in terms of the Jacobian elliptic functions, $`J_e`$ ,
$$V(\varphi )=V_0\{J_e\left[\beta _0(\varphi \varphi _0)\right],m\}^4,$$
(2.15)
where $`\beta _0\beta _0(\mathrm{\Omega }_i)`$ and $`mm(\mathrm{\Omega }_i)1`$ (with $`i`$ some elements of the set $`\{\varphi ,\mathrm{\Lambda },M,k\}`$) are dimensionless quantities whose explicit shape will depend on the particular function $`J_e`$ being considered. For reasons which will become clear in the next section, of particular interest for reproducing a suitable accelerating model of the universe are the elliptic functions $`J_e=\mathrm{cn}`$, for which
$`\beta _0={\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_\varphi }}\left(1+{\displaystyle \frac{\mathrm{\Omega }_\mathrm{\Lambda }}{1\mathrm{\Omega }_M}}\right)},m={\displaystyle \frac{\mathrm{\Omega }_\mathrm{\Lambda }}{1+\mathrm{\Omega }_\mathrm{\Lambda }\mathrm{\Omega }_M}},`$
$`J_e=\mathrm{nc}`$, for which
$`\beta _0={\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_\varphi }}\left(1+{\displaystyle \frac{\mathrm{\Omega }_\mathrm{\Lambda }}{1\mathrm{\Omega }_M}}\right)},m={\displaystyle \frac{1\mathrm{\Omega }_M}{1+\mathrm{\Omega }_\mathrm{\Lambda }\mathrm{\Omega }_M}},`$
and $`J_e=\mathrm{sd}`$, for which
$`\beta _0={\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_\varphi }}},m={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1\mathrm{\Omega }_\mathrm{\Lambda }}{\mathrm{\Omega }_M}}\right).`$
Note, furthermore, that in the limiting case that $`m1`$, the $`\mathrm{cn}`$-solution becomes a $`\mathrm{sech}`$-solution for an open universe, the $`\mathrm{nc}`$-solution becomes a $`\mathrm{cosh}`$-solution for $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, whereas the $`\mathrm{sd}`$-solution reduces to a $`\mathrm{sinh}`$-solution for $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ (i.e. the potential first found by Di Prieto and Demaret , as it should be expected.)
On the other hand, for $`\omega =2/3`$, we obtain the potentials satisfying constraint (2.13):
$$V(\varphi )=V_0\left[A_\pm \mathrm{sinh}(\varphi \varphi _0)+B\mathrm{cosh}(\varphi \varphi _0)C\right]^1,$$
(2.16)
where
$$A_\pm =\pm \sqrt{\frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }}}\mathrm{cosh}\delta _0^{}\frac{\mathrm{\Omega }_k^2e^{\sigma _0^{}}}{4\mathrm{\Omega }_M\left(\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }\right)}$$
(2.17)
$$B=\sqrt{\frac{\mathrm{\Omega }_M}{\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }}}\mathrm{sinh}\delta _0^{}+\frac{\mathrm{\Omega }_k^2e^{\sigma _0^{}}}{4\mathrm{\Omega }_M\left(\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }\right)}$$
(2.18)
and
$$C=\frac{1}{2}\frac{\mathrm{\Omega }_k}{\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }},$$
(2.19)
with $`\delta _0^{}`$ and $`\sigma _0^{}`$ as given by $`\delta _0`$ and $`\sigma _0`$ in Eq. (1.3), but with $`\mathrm{\Omega }_\varphi `$ replaced for $`\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }`$. Clearly, for $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, solutions (2.16) reduce to solutions (1.3).
Finally, we can also have solutions to (2.13) for any $`\omega `$, satisfying $`1<\omega 0`$. These solution are in turn obtained for particular values of the cosmological parameters $`\mathrm{\Omega }_j`$. Thus, setting $`\mathrm{\Omega }_k=\mathrm{\Omega }_\mathrm{\Lambda }=0`$, we have
$$V=V_0\mathrm{sinh}^{\frac{2(\omega +1)}{\omega }}\left[\pm \left(\frac{\sqrt{\mathrm{\Omega }_M}\omega }{2(\omega +1)}\right)(\varphi \varphi _0)\right],$$
(2.20)
and for $`\mathrm{\Omega }_k=\mathrm{\Omega }_M=0`$,
$$V=V_0\mathrm{sinh}^{\frac{6(\omega +1)}{3\omega +2}}\left[\pm \sqrt{\frac{3\omega +2}{2}}(\varphi \varphi _0)\right].$$
(2.21)
Potentials (2.20) and (2.21) and at least some of the potentials in the family (2.15) can actually be regarded as generalizations from inverse-power law potentials which hold only as approximations for small values of $`\varphi \varphi _0`$, at large values of the redshift (see Sec. III). Let us for example consider the case $`Je=sd`$ in the family (2.15). For most of its cosmological evolution $`\varphi \varphi _0`$ remains very small at large values of the redshift, so that we can approximate $`V(\varphi \varphi _0)^4`$, except when the potential approaches current values. Moreover, though for quintessence the most interesting models are those where $`\omega `$ is not constant and, in particular, the tracker models , one can see that at least some of the good properties of these models may be somehow shared by the potential considered in this paper. Since at least some of our potentials can be approximated as inverse-power law functions of the field containing at least one free parameter, along their primordial evolution these potentials can be implemented in the realm of high energy physics and linked to particle models with dynamical symmetry breaking or nonperturbative effects . On the other hand, it appears that such potentials can also help solving the cosmic coincidence problem . Taking again as an illustrative example the solution $`Je=sd`$ in Eq. (2.15) we see (see Eq. (3.14)) that if one sets the initial conditions inmediately after inflation, i.e. at a redshift $`z10^{28}`$, then $`\varphi \varphi _0`$ initially, and $`\varphi \varphi _01`$ only now, so explaining why the quintessence field begins to dominate now. Tracker models also seem to improve the fine-tuning problem; we hope this to be the case with some of our potentials as well, in particular in those models where the total vacuum energy, $`\mathrm{\Omega }\varphi +\mathrm{\Omega }_\mathrm{\Lambda }`$, is zero, or generally for the reasons discussed in Sec. IV.
At first glance studying the cases $`\omega =1/6,1/3`$ could seem without any physical motivation, as several authors have already shown that quintessence models based on such values should be ruled out. However, our quintessence approach is based on the idea that the vacuum field is splitted into two parts, one manifested as a varying cosmological constant with positive energy density, and the other, the quintessence field, always having negative energy density. This splitting appears to enlarge the allowed domain of $`\omega `$-values which are physically relevant and reproduce the wanted accelerating expansion of the universe (see Sec. III). This translates, in particular, in a generalized expression for the decelaration parameter (Eq. (3.7)), according to which no value of $`\omega `$ can be ruled out from the onset.
## III Predicting cosmological models from quintessence potentials
In this section we are going to check whether the considered solutions satisfying the constraints on the scalar-field potentials (1.1) and (2.13) are suitable potentials to fullfil the quintessence’s aim, that is, whether such solutions are able to predict cosmological models matching observations from supernova Ia. More precisely, we want to compute the luminosity distance($`D_L`$)-redshift($`z`$) relations from the quintessence potentials satisfying the constraints. Bringing then that relation, $`D_L(z)`$, into the magnitude-redshift relation ,
$$m_B^{\mathrm{eff}}=\stackrel{ˇ}{M}_B+5\mathrm{log}[D_L(z)],$$
(3.1)
where $`\stackrel{ˇ}{M}_B=M_B5\mathrm{log}H_0+25`$ is the Hubble-constant-free $`B`$-band absolute magnitude at the maximum of a Ia supernova whose values have to be suitably corrected , we could thus directly compare the predictions from the cosmological quintessence potentials with the conveniently corrected observed magnitude. The luminosity distance $`D_L`$ depends only on the theory we are working on, and can be written as
$`D_LH_0={\displaystyle \frac{(1+z)}{\sqrt{|\mathrm{\Omega }_k|}}}S\{\sqrt{|\mathrm{\Omega }_k|}\times `$
$$_0^zdz^{}[\underset{j}{}\mathrm{\Omega }_j(1+z^{})^{3(1+\alpha _j)}+\mathrm{\Omega }_k(1+z^{})^2]^{\frac{1}{2}}\},$$
(3.2)
in which $`S\{x\}=\mathrm{sin}x`$ for $`k=+1`$, $`S\{x\}=x`$ for $`k=0`$ and $`S\{x\}=\mathrm{sinh}x`$ for $`k=1`$, and the parameter $`\alpha _i`$ is defined from the energy density so that
$`\rho _j\left({\displaystyle \frac{R_0}{R}}\right)^{3(1+\alpha _j)},`$
with the subscript $`j`$ labelling the distinct nongeometrical contributions, namely $`j=M,\varphi `$ and $`\mathrm{\Lambda }`$, in such a way that $`\alpha _M=0`$, $`\alpha _\varphi =\omega `$ and $`\alpha _\mathrm{\Lambda }=2/3`$. Using the definition of the redshift in terms of the scale factor $`R`$, we then attain for the argument of the squared root in the integrand of Eq. (3.2) in case of the generalized quintessence model,
$`\mathrm{\Pi }\mathrm{\Omega }_M\left({\displaystyle \frac{R_0}{R}}\right)^3`$
$$+\mathrm{\Omega }_\varphi \left(\frac{R_0}{R}\right)^{3(1+\omega )}+\mathrm{\Omega }_\mathrm{\Lambda }\left(\frac{R_0}{R}\right)+\mathrm{\Omega }_k\left(\frac{R_0}{R}\right)^2.$$
(3.3)
Inserting now the relation $`V/V_0=(R_0/R)^{3(1+\omega )}`$, we obtain for $`\mathrm{\Pi }`$:
$`\mathrm{\Pi }\mathrm{\Omega }_M\left({\displaystyle \frac{V}{V_0}}\right)^{\frac{1}{\omega +1}}`$
$$+\mathrm{\Omega }_\varphi \left(\frac{V}{V_0}\right)+\mathrm{\Omega }_\mathrm{\Lambda }\left(\frac{V}{V_0}\right)^{\frac{1}{3(\omega +1)}}+\mathrm{\Omega }_k\left(\frac{V}{V_0}\right)^{\frac{2}{3(\omega +1)}}.$$
(3.4)
It is now readily realized that $`\mathrm{\Pi }^{1/2}`$ is the same as $`V_0^{}V^{1/2}/V^{}V_0^{1/2}`$ as obtained from the constraint (2.13). Hence,
$`{\displaystyle _0^z}{\displaystyle \frac{dz^{}}{\sqrt{\mathrm{\Pi }}}}=`$
$$\frac{V_0^{}}{3(\omega +1)V_0^{\frac{1}{2}+\frac{1}{3(\omega +1)}}}_{\varphi (0)}^{\varphi (z)}\frac{d\varphi }{V(\varphi )^{\frac{3\omega +1}{6(\omega +1)}}},$$
(3.5)
and therefore we have the following relation between the luminosity distance and the quintessence potential
$`D_LH_0=`$
$$\frac{(1+z)}{\sqrt{|\mathrm{\Omega }_k|}}S\left\{\frac{V_0^{}\sqrt{|\mathrm{\Omega }_k|}}{3(\omega +1)V_0^{\frac{1}{2}+\frac{1}{3(\omega +1)}}}_{\varphi (0)}^{\varphi (z)}\frac{d\varphi }{V(\varphi )^{\frac{3\omega +1}{6(\omega +1)}}}\right\}.$$
(3.6)
On the other hand, the deceleration parameter $`q_0`$ can also be expressed in terms of the quintessence parameter $`\omega `$ as follows
$`q_0={\displaystyle \frac{1}{2}}{\displaystyle \underset{j}{}}\mathrm{\Omega }_j(1+3\alpha _j)=`$
$$\frac{1}{2}\left[\mathrm{\Omega }_M+\mathrm{\Omega }_\varphi (1+3\omega )\mathrm{\Omega }_\mathrm{\Lambda }\right].$$
(3.7)
Thus, in order to reproduce the wanted accelerating behaviour of the universe we must have a suitable combination for the values of the parameters $`\mathrm{\Omega }_j`$ and $`\omega `$, such that the resulting value of $`q_0`$ be negative. We note that for pure quintessence models with $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, any scalar-field potentials defined for $`\omega >1/3`$ could only give a topological accelerating behaviour if $`\mathrm{\Omega }_\varphi <0`$. Clearly, for $`\omega =1/3`$, irrespective of the value of $`\mathrm{\Omega }_\varphi `$, we always have $`q_0>0`$, provided $`\mathrm{\Omega }_M>0`$. In generalized quintessence models with $`\mathrm{\Omega }_\mathrm{\Lambda }>0`$, all the above situations predicting deceleration could still predict acceleration for sufficiently high positive values of $`\mathrm{\Omega }_\mathrm{\Lambda }`$.
Let us consider in what follows the cosmological predictions from the permissible quintessence potentials dealt with in Secs. I and II for the case that $`\mathrm{\Omega }_M>0`$. We shall start with the family of potentials (1.4) for $`\omega =1/6`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, in which case Eqs. (3.5) and (3.6) become
$`D_LH_0=`$
$$\frac{1+z}{\sqrt{|\mathrm{\Omega }_k|}}S\left\{\frac{\sqrt{14|\mathrm{\Omega }_k|}}{5\sqrt{\mathrm{\Omega }_\varphi }}_{\varphi (0)}^{\varphi (z)}𝑑\varphi J_e[\alpha _0(\varphi \varphi _0),m]\right\},$$
(3.8)
where we have used the definitions (2.14), and
$$q_0=\frac{1}{2}\left(\mathrm{\Omega }_M+\frac{1}{2}\mathrm{\Omega }_\varphi \right).$$
(3.9)
Substituting the different Jacobian elliptic functions in the luminosity distance expression (3.7) and integrating the resulting expression, it turns out that the only of such functions for which we obtain a consistent, real $`D_Lz`$ relation predicting a nonclosed universe with positive vacuum energy density which is dynamically accelerating is $`J_e=\mathrm{cn}`$. In this case the integral in Eq. (3.8) gives
$`{\displaystyle \frac{1}{\alpha _0\sqrt{m}}}\mathrm{arccos}\sqrt{1m+{\displaystyle \frac{m}{\sqrt{1+z^{}}}}}|_0^z,`$
where
$`\alpha _0=\sqrt{{\displaystyle \frac{7\mathrm{\Omega }_k}{(1m)\mathrm{\Omega }_\varphi }}}`$
$`m={\displaystyle \frac{\mathrm{\Omega }_k}{\mathrm{\Omega }_\varphi +2\mathrm{\Omega }_k}}.`$
For whichever combinations of values of the cosmological parameters $`\mathrm{\Omega }_M`$, $`\mathrm{\Omega }_k`$ and $`\mathrm{\Omega }_\varphi `$ satisfying the triangular constraint $`\mathrm{\Omega }_k+\mathrm{\Omega }_M+\mathrm{\Omega }_\varphi =1`$ we then obtain a dynamically accelerating universe which, according to Eq. (3.8), is however topologically decelerating if the vacuum energy density is positive. On the other hand, although all possible resulting $`5\mathrm{log}D_LH_0z`$ plots give a nearly straight line between $`z0.01`$ and $`z0.5`$ which appear to slightly accelerate thereafter, the full $`5\mathrm{log}D_LH_0`$-interval that corresponds to the $`z`$-interval of available Ia supernova observations ($`(0.011)`$) is always around 6, quite smaller than the observed value $`m_B^{\mathrm{eff}}13`$ . Thus, the scalar-field potentials (1.4) cannot conform to the observations on supernovae Ia.
We consider next the potentials for $`\omega =1/3`$ which are given by the general expression (2.15) for $`\mathrm{\Omega }_M>0`$ and total vacuum energy density $`\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }0`$. In these cases, Eqs. (3.5) and (3.6) give
$$D_LH_0=\frac{1+z}{\sqrt{|\mathrm{\Omega }_k|}}S\left\{\sqrt{\frac{|\mathrm{\Omega }_k|}{\mathrm{\Omega }_\varphi }}\varphi (z^{})|_0^z\right\}$$
(3.10)
$$q_0=\frac{1}{2}\left(\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda }\right).$$
(3.11)
For the 12 different Jacobian elliptic functions involved in solutions (2.15) we have derived the expressions of the scalar field in terms of the redshift, $`\varphi (z)`$, in the form of elliptic integrals of the first kind . It turns out that only the elliptic functions $`\mathrm{cn},\mathrm{nc}`$ and $`\mathrm{sd}`$ can generate non closed universes which are both topologically and dynamically accelerating. In the case that $`J_e=\mathrm{cn}`$, we have for the luminosity distance
$`D_LH_0={\displaystyle \frac{1+z}{\sqrt{|\mathrm{\Omega }_k|}}}\times `$
$$S\left\{2\sqrt{\frac{|\mathrm{\Omega }_k|}{\mathrm{\Omega }_M\left(1+\frac{\mathrm{\Omega }_k}{1\mathrm{\Omega }_M}\right)}}F[\mathrm{arcsin}\sqrt{\frac{z^{}}{1+z^{}}},m]|_0^z\right\},$$
(3.12)
with
$`m={\displaystyle \frac{\mathrm{\Omega }_\mathrm{\Lambda }}{1\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }}},`$
and for $`J_e=\mathrm{nc}`$,
$`D_LH_0={\displaystyle \frac{1+z}{\sqrt{|\mathrm{\Omega }_k|}}}\times `$
$$S\left\{2\sqrt{\frac{|\mathrm{\Omega }_k|}{\mathrm{\Omega }_M\left(1+\frac{\mathrm{\Omega }_k}{1\mathrm{\Omega }_M}\right)}}F[\mathrm{arcsin}(i\sqrt{z^{}}),m]|_0^z\right\},$$
(3.13)
with
$`m={\displaystyle \frac{1\mathrm{\Omega }_M}{1\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }}}.`$
In expressions (3.11) and (3.12) the symbol $`F`$ denotes elliptic integral of the first kind . These expressions give $`5\mathrm{log}D_LH_0z`$ plots for different combinations of cosmological parameters satisfying the quadrilateral constraint $`\mathrm{\Omega }_k+\mathrm{\Omega }_M+\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }=1`$ and conditions $`\mathrm{\Omega }_k,\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }0`$ which represent accelerating expansion with suitably slight deviations from straight lines occurring at $`z0.5`$ only. However, again as for the case $`\omega =1/6`$, the full variations of $`5\mathrm{log}D_LH_0`$ along the $`z`$-observation interval $`6`$, are quite smaller than the corresponding value observed in supernovae.
If we take for $`J_e`$ the function $`\mathrm{sd}`$, then the $`z`$-dependence of the scalar field can be expressed in the form
$$\varphi (z^{})=\varphi _0+2\sqrt{\frac{\mathrm{\Omega }_\varphi }{\mathrm{\Omega }_M}}F[\mathrm{arcsin}\frac{1}{\sqrt{1+m+z^{}}},m],$$
(3.14)
where
$`m={\displaystyle \frac{1}{2}}\left(\sqrt{1+{\displaystyle \frac{4}{\mathrm{\Omega }_M}}}1\right).`$
Inserting the scalar field (3.13) in Eq. (3.10) we get the wanted $`D_Lz`$ relation. In this case, this relation gives plots which show the required accelerating expansion after $`z0.6`$, both for flat and open universes, with a full variation of $`5\mathrm{log}D_LH_0`$ along the observed $`z`$-interval of the order 13, fitting well with the observations , along all available $`z`$-values.
Finally, for the case $`\omega =2/3`$, the relation (3.5) reduces to
$`D_LH_0={\displaystyle \frac{1+z}{\sqrt{|\mathrm{\Omega }_k|}}}S\{\sqrt{{\displaystyle \frac{5|\mathrm{\Omega }_k|}{\mathrm{\Omega }_\varphi }}}\times `$
$$_{\varphi (0)}^{\varphi (z)}\frac{d\varphi }{\sqrt{A_\pm \mathrm{sinh}(\varphi \varphi _0)B\mathrm{cosh}(\varphi \varphi _0)C}}\},$$
(3.15)
with the constants $`A_\pm `$, $`B`$ and $`C`$ as defined by expressions (2.17)-(2.19), and the deceleration parameter given now by
$$q_0=\frac{1}{2}\left[\mathrm{\Omega }_M\left(\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }\right)\right].$$
(3.16)
In order to obtain cosmological models described by real $`D_Lz`$ relations the conditions of integration in expression (3.14) must be such that the cosmological parameters satisfy the following conditions
$$\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }+\frac{1}{2}\left(\mathrm{\Omega }_k\pm \sqrt{2\mathrm{\Omega }_k}\right)=0$$
(3.17)
$$Q^44(4\mathrm{\Omega }_M)\left(\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }\right)=\mathrm{\Omega }_k^2$$
(3.18)
and, either
$$\mathrm{\Omega }_k=1\pm \sqrt{1+\frac{1}{4}\left(\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }\right)\left(\mathrm{\Omega }_M16\right)},$$
(3.19)
or,
$$\mathrm{\Omega }_k2\left(\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }\right)=0.$$
(3.20)
The parameter $`Q`$ in Eq. (3.18) has been introduced to simplify the equation. It is defined in terms of the $`\mathrm{\Omega }`$’s only, as:
$$Q^2=4\left(\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }\right)+\left[2\left(\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }\right)+\mathrm{\Omega }_k\right]^2.$$
(3.21)
There are two cosmological models which satisfy all these conditions. They are: Model I, a flat universe defined by the parameters $`\mathrm{\Omega }_M=1`$, $`\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }=\mathrm{\Omega }_k=0`$, and Model II, an open universe defined by the parameters $`\mathrm{\Omega }_M=\mathrm{\Omega }_\varphi +\mathrm{\Omega }_\mathrm{\Lambda }=1/4`$, $`\mathrm{\Omega }_k=1/2`$. In both cases $`B=0`$ and $`A_\pm A`$ and $`C`$ become indeterminate and real. For Model I we obtain $`q_0=+1/2`$ and, from expression (3.15),
$`D_LH_0=\sqrt{{\displaystyle \frac{2}{|A|}}}(1+z)\sqrt{|A|^2\left(|C|^2{\displaystyle \frac{1}{1+z}}\right)^2}\times `$
$$F[\mathrm{arcsin}\sqrt{\frac{|A|+|C|\frac{1}{1+z^{}}}{|A|+|C|}},\sqrt{\frac{|A|+|C|}{2|A|}}]|_0^z,$$
(3.22)
which always gives rise to a topologically and dynamically decelerating universe for any combinations of constants $`A`$ and $`C`$ satisfying $`|A|+|C|=1`$ and the integration condition $`|A|>|C|>0`$.
More interesting is Model II, for which one obtains a topologically uniform expansion, $`q_0=0`$, and again the integration condition $`|A|>|C|>0`$. Taking e.g. $`|A|=2`$ and $`|C|=1`$, it follows from expression (3.15)
$`D_LH_0=\sqrt{2}(1+z)\mathrm{sin}\{{\displaystyle \frac{(z+1)\sqrt{2\left[4\left(1+\frac{1}{1+z}\right)^2\right]}}{z+2}}\times `$
$$F[\mathrm{arcsin}\sqrt{2\left(1+\frac{1}{1+z}\right)},\frac{1}{2}]\}.$$
(3.23)
Eq. (3.23) gives rise to a $`D_Lz`$ plot which, in spite of being associated with a topologically uniform universe, starts accelerating after $`z0.5`$ in a way that matches the behaviour observed in distant supernovas. That plot, on the other hand, consistently shows a variation $`(5\mathrm{log}D_LH_0)13`$ along the observed $`z`$-interval, fitting well with the observations at all available values of the redshift. Therefore, it could be thought that Model II and the solution in terms of the Jacobian elliptic function $`\mathrm{sd}`$ for $`\omega =1/3`$ dealt with above, correspond to ”good” quintessence potentials.
It appears also relevant to perform a similar computation for the inverse-power law potentials obtained in Sec. II, which are of the type already proposed in the literature and that might be justified from high energy physics. For the first of these potentials, $`V=V_0(\varphi _0/\varphi )^4`$, one obtains
$$D_LH_0=(1+z)\mathrm{sinh}\left[\pm \sqrt{\frac{1}{2\mathrm{\Omega }_\varphi }}\left(1\frac{1}{\sqrt{1+z}}\right)\right],$$
(3.24)
and for $`V=V_0(\varphi _0/\varphi )^2`$,
$$D_LH_0=(1+z)\mathrm{sin}\left[\pm \sqrt{\frac{1}{\mathrm{\Omega }_\varphi }}\mathrm{ln}\left(\frac{1}{\sqrt{1+z}}\right)\right];$$
(3.25)
these two equations are, of course, not valid for $`k=0`$. It is interesting to note that in the two cases, we reproduce a $`D_Lz`$ plot which nearly matches the observed results, producing a distinguishable accelerating pattern starting at an expected $`z0.5`$, and a variation $`\mathrm{}(5\mathrm{log}D_LH_0)`$ between 11 and 12, only slightly smaller than what has been measured, along the observed $`z`$-interval. Moreover, our analytical formulae for the luminosity distance-redshift relation can also be applied to potentials which are defined for any value of $`\omega `$, as those given in Eqs. (2.20) and (2.21), for sufficiently large values of the redshift. Thus, for potential (2.20) we obtain
$$D_LH_02(1+z)\sqrt{\frac{1+\omega }{3\mathrm{\Omega }_\varphi \mathrm{\Omega }_M}}\left(1\frac{1}{\sqrt{1+z}}\right),$$
(3.26)
at large $`z`$, and for the potential (2.21),
$$D_LH_0(1+z)\sqrt{\frac{2(3\omega +2)}{3(\omega +1)\mathrm{\Omega }\varphi }}\left(\sqrt{1+z}1\right),$$
(3.27)
at large $`z`$ for $`\omega <2/3`$. It can be checked that these functions give $`D_Lz`$ plots which show a nearly uniform expansion at the allowed sufficiently large values of the redshift.
It is worth noticing that for the limiting expressions from solutions (2.15) and (2.16), obtained by restricting $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ and $`\mathrm{\Omega }_\varphi 0`$, we either cannot even obtain a topologically accelerating universe ($`\omega =1/3`$), or have no consistent integration procedure along the complete range of allowed $`z`$-values that leads to a definite real luminosity distance ($`\omega =2/3`$). Thus, at least for the particular potentials considered in this work, if we want to consistently predict cosmological models compatible with observations on Ia supernovae at large and moderate redshifts, it appears that quintessence should be generalized in a way that allows for a more complicate vacuum structure made up of (i) a time-dependent, ”axionic” (as it is pure imaginary classically) scalar field, $`\varphi (t)`$, with positive pressure and negative energy density, and (ii) a time-varying positive cosmological term, $`\mathrm{\Lambda }(t)`$, whose current value $`\mathrm{\Lambda }_0`$ can be quite smaller than that with which it started the cosmological evolution, in such a way that the full vacuum energy density is restricted to be $`\rho _\varphi +\rho _\mathrm{\Lambda }0`$.
## IV Summary and discussion
In this paper we have considered the problem of the quintessence potential, restricting ourselves to a constant equation of state, that is: what are the permissible potentials for a vacuum, time-dependent scalar field predicting cosmological models that conform to recent observations and, at the same time, satisfy the constraints imposed by the field equations and conservation laws, discussing their physical relevance in the case that the quintessence field corresponds to a non-tracking constant equation of state. We have generalized the usual quintessence model, introducing a positive cosmological varying term, while restricting the scalar-field energy density to be definite negative and its pressure definite positive in such a way that the overall vacuum energy density is necessarily positive or vanishing. Quintessence potentials that satisfy the above-alluded constraint, both for the usual models and for models with a varying cosmological term, have been obtained for particular values of the constant parameter defining the state equation for the scalar field in the two kinds of models. None of these potentials have hitherto been used in quintessence, except those which are given as an inverse-power law. We also obtain potentials which are given in terms of either hyperbolic functions or Jacobian elliptic functions, the latter generally reducing to the former in the limit when the cosmological term tends to zero. We have also obtained a general expression relating the luminosity distance with the quintessence potential, and this has been integrated for all particular solutions expressed in terms of the redshift. It turned out that only some of such solutions with nonzero cosmological term are able to produce cosmological models that conform to an acceleratingly expanding universe and agree with recent observations on Ia supernovae.
The cosmological term $`\mathrm{\Lambda }`$ we have used in our generalized quintessence model depends on the cosmological time through a linear dependence on the scale factor $`R`$. This varying character of $`\mathrm{\Lambda }`$ could a priori be an useful property to help solving the known cosmological constant problem, though not by itself only. Actually, one could not justify how an initial vacuum energy density of the order $`M_p^4`$ may be lowered down to a value smaller than $`10^{47}`$ GeV invoking such a dependence; instead, if the $`R`$-dependence of $`\mathrm{\Lambda }`$ would be assumed to be the same along the entire cosmological evolution, then one would need a conservation law for $`\mathrm{\Lambda }`$ given by $`\mathrm{\Lambda }=\mathrm{\Lambda }_0(R_0/R)^\gamma `$, with $`\gamma 123/50`$. However, the conservation law chosen in this paper, $`\mathrm{\Lambda }=\mathrm{\Lambda }_0(R_0/R)`$, is assumed to hold only in the late classical cosmological regime; it could well be that during primordial expansion $`\gamma `$ had taken on values larger than $`123/50`$. For example, taking $`\gamma 3`$ during the primeval expansion up to $`R10^4`$ cm, and $`\gamma =1`$ thereafter, would solve the cosmological constant problem. The price one would pay to get such a big reward would just be the allowance for the dynamical content of the quintessence field to be that of the conventional models, with $`\gamma =3`$ and $`\gamma =2`$ (see Sec. II) during its early evolution.
The conclusions obtained in this work are not general. They just refer to the solutions that correspond to the particular values of the quintessence parameter $`\omega =1/6,1/3,2/3`$ and -1, and some special cases for any $`\omega `$. Possibly there will be other potentials corresponding to different, intermediate values of $`\omega `$ that also reproduce the observed cosmological expansion within the generalized quintessence model. We do not believe however this to be the case in the realm of the conventional quintessence model, though more work is obviously needed to reach a final verdit on this.
###### Acknowledgements.
For helpful comments and a careful reading of the manuscript, the author thanks C. Sigüenza. Thanks are also due to M. Moles for enlightening conversations. This research was supported by DGICYT under Research Project No. PB97-1218.
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# A Brief Note on Discrete Space Effects
## Abstract
We discuss briefly discrete space effects and show that this leads to space reflection asymmetry and also a minor modification of Einstein’s energy-mass formula.
Snyder had shown how it is possible to consider discrete space time consistently with the Lorentz transformation. Discrete space time continues to receive attention over the years from several scholars (Cf. for example for details). Indeed in the Dirac’s relativistic theory of the electron, this discretization is evident - averages over the Compton scale are required to eliminate Zitterbewegung effects and non Hermitian position operators and to recover meaningful physics \- and that includes special relativity. It is in this light that a recent formulation of an electron in terms of a Kerr-Newman metric becomes meaningful and further this leads to a meaningful, if phenomenological mass spectrum .
All this pleasingly dovetails with the fact that if the minimum space time cut offs are taken at the Compton scale, then we have a non commutative geometry viz.,
$$[x,y]=0(l^2)$$
(1)
and similar equations, and which further, leads directly to the Dirac equation (Cf. for details). This can be easily seen from the fact that, given (1), the usual infinitessimal coordinate shift in Minkowski space, is,
$$\psi ^{}(x_j)=[1+ıϵ(ıϵ_{ljk}x_k\frac{}{x_j})+0(ϵ^2)]\psi (x_j)$$
(2)
The choice
$$t=\left(\begin{array}{c}10\hfill \\ 01\hfill \end{array}\right),\stackrel{}{x}=\left(\begin{array}{c}0\stackrel{}{\sigma }\hfill \\ \stackrel{}{\sigma }0\hfill \end{array}\right)$$
provides a representation for the coordinates, as can be easily verified and then from (2) we recover the Dirac equation.
If on the other hand terms $`l^2`$ are neglected, then (1) gives the usual commutation relations of Quantum Theory. Discrete space time is therefore a higher order correction to usual Quantum Theory.
In this context it has been pointed out that the discrete time provides an explanation for the puzzling Kaon decay which violates time reversal symmetry and also the decay of a pion into an electron and a positron.
We now observe that from an intuitive point of view space or time reversal symmetries based on space time points theory cannot be taken for granted if space time is discrete. This can immediately be seen from (1): If we retain terms $`l^2`$, then there is no invariance under space reflections.
Indeed in the same vein, as discussed earlier, the fact that the Compton wavelength of the nearly massless neutrino is very large provides an explanation of its handedness.
We finally point out another $`0(l^2)`$ effect, which can be demonstrated in a simple way by invoking the derivation of the wave equation of a particle by replacing the continuum by a set of ”lattice points”. In this case we have an equation like
$$Ea(x_n)=E_0a(x_n)Aa(x_n+b)Aa(x_nb).$$
where $`b`$ is the distance between successive ”lattice points”. This leads to
$$E=E_02Acoskb.$$
(3)
We can choose the zero of energy in such a way that when $`b0`$, we have $`E=2A`$ which is then identified with the rest energy $`mc^2`$ of the particle(Cf. for details).
However if we do not neglect terms $`b^2=l^2`$, then we have from (3)
$$\left|\frac{E}{mc^2}1\right|0(l^2)$$
It must be reiterated that when terms $`l^2`$ are neglected, we recover the usual theory. Finally, the above conclusions are true with minor modifications in case the minimum cut off is non-zero but not the Compton wavelength.
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# Effects of uniaxial strain in 𝐋𝐚𝐌𝐧𝐎_𝟑
## I Introduction
The ‘colossal’ magnetoresistive manganese perovskites have been a focus of recent attention. Since most of the technological applications require thin films on substrates, it is important to understand the effects of strains induced by substrates. Because Mn $`e_g`$ electrons, which determine important physical properties of these materials, are coupled to the lattice degrees of freedom through the Jahn-Teller (JT) coupling, it is expected that uniaxial or biaxial strains are important and that even relatively small strains may result in observable effects on the electronic properties of these materials. Recently, the effects of substrate-induced strains on the properties of thin manganate films have been studied experimentally. It is indeed found that the ferromagnetic and metal-insulator transition temperature, $`T_c`$, depends sensitively on the volume-preserving uniaxial strains, as do the magnetic anisotropy, magnetoresistance, and charge ordering transition.
In this work, we study the effects of uniaxial strains in $`\mathrm{LaMnO}_3`$, which is the parent compound of the doped manganese perovskites. Our motivations are to further clarify the properties of this interesting compound and to test techniques and fix parameters so the more complicated behaviors of the doped compounds may be studied. At very high temperatures, bulk $`\mathrm{LaMnO}_3`$ exists in cubic perovskite structure, but at T $`<`$ 750 K, it has a static $`(\pi ,\pi ,0)`$ $`3x^2r^2/3y^2r^2`$ type Jahn-Teller distortion. It also has a uniform tetragonal distortion, which originates from the coupling of the staggered and uniform distortions by an anharmonic elastic energy. This coupling implies that the substrate-induced strain may affect the orbital ordering.
In Ref. , following Kanamori, a model with harmonic Mn-O and Mn-Mn elastic forces and a local anharmonic energy term was used to study lattice distortions in this material. It predicted $`(\pi ,\pi )`$ type ordering in $`xy`$ plane as observed in $`\mathrm{LaMnO}_3`$. However, according to this model, $`(\pi ,\pi ,0)`$ and $`(\pi ,\pi ,\pi )`$ type orderings have the same energy. The reason why $`(\pi ,\pi ,0)`$ ordering is favored has not been understood so far. Bulk $`\mathrm{LaMnO}_3`$ has so called A-type antiferromagnetic (AF) ordering below 140 K, in which spins align parallel in $`xy`$ plane and antiparallel along $`z`$ direction. This peculiar spin ordering pattern is the result of the exchange interaction between Mn ions which depends on the $`e_g`$ orbital ordering as explained in Ref. . Indications of a coupling between orbital ordering and magnetic ordering have also been found in a recent X-ray resonant scattering experiment.
In this paper, we present more general expressions for the elastic energy than those in Ref. . We use these to study the ground state energy and distortions in bulk state, and the effects of uniaxial strains in thin films. We also present a simple model of the magnetic interaction depending explicitly on the orbital states, and we use this to study the magnetostriction effects in bulk state and the change of the magnetic interaction due to the strains in thin films. We examine the changes in the band structure and optical conductivity due to the strains, using a nearest-neighbor tight binding approximation. We also compare the energies of $`(\pi ,\pi ,0)`$ and $`(\pi ,\pi ,\pi )`$ type orderings, and examine why the observed $`(\pi ,\pi ,0)`$ distortion is favored over the $`(\pi ,\pi ,\pi )`$ type distortion. Our calculations suggest that 2 % tensile strain can change the observed A-type(layered) antiferromagnetic state into a purely antiferromagnetic state. This change would induce large changes in band structure and optical conductivity spectrum, which we calculate. Finally, we show that the magnetostriction effect at the Neel transition is large.
The rest of the paper is organized in the following way. Section II introduces a model of elastic energy, Sect. III a model of the magnetic interaction and magnetostriction effect, and Sect. IV a tight binding model for band structure and optical conductivity. Section V compares the ground state energies of $`(\pi ,\pi ,0)`$ and $`(\pi ,\pi ,\pi )`$ type distortions. Section VI presents the results. Section VII summarizes the main conclusions. In Appendices A and B, we show how we determine general expressions of the elastic energy, and how we determine the parameters, respectively.
## II Model of lattice energy
### A Overview
To study the effects of uniaxial strains in thin films, we first need to understand the strains already present in bulk state. The elastic energy depends on three important variables: the oxygen displacement along Mn-Mn direction, the three-dimensional Mn ion displacement, and the Mn $`e_g`$ electron orbital state.$`^\text{,}`$ More precisely, the degrees of freedom we will consider are defined in the following way: In the ideal cubic perovskite structure with lattice constant $`a_0`$, the Mn ions are located at $`a_0\stackrel{}{i}`$, and oxygen ions at $`a_0(\stackrel{}{i}+\widehat{a}/2)`$, where $`i_x`$, $`i_y`$, and $`i_z`$ are integers and $`\widehat{a}`$=$`\widehat{x}`$, $`\widehat{y}`$, and $`\widehat{z}`$. We write the displacement of Mn at $`a_0\stackrel{}{i}`$ as $`a_0(\stackrel{}{e}\stackrel{}{i}+\stackrel{}{\delta }_\stackrel{}{i})`$, where $`\stackrel{}{\delta }_\stackrel{}{i}`$ represents nonzero-wavevector Mn-ion displacement, $`\stackrel{}{e}=e_{xx}\widehat{x}+e_{yy}\widehat{y}+e_{zz}\widehat{z}`$, and $`e_{ij}`$ is the conventional strain tensor referred to the ideal cubic perovskite lattice (we need only the diagonal components). The displacement of oxygen at $`a_0(\stackrel{}{i}+\widehat{a}/2)`$ along Mn-Mn axis is $`a_0[\stackrel{}{e}(\stackrel{}{i}+\widehat{a}/2)+u_\stackrel{}{i}^a\widehat{a}]`$, where $`u_\stackrel{}{i}^a`$ represents O-ion displacement with nonzero wavevector. Figure 1 shows $`\stackrel{}{\delta }_\stackrel{}{i}`$ and $`u_\stackrel{}{i}^{x,y,z}`$. We assume that we have already minimized the elastic energy over the displacements of O ions perpendicular to the Mn-Mn axis and the displacement of La ions. Therefore, even though we do not consider the buckling of the Mn-O-Mn bond explicitly, its effect is implicitly included in our harmonic and anharmonic elastic constants below.
We treat the elastic energy due to these strains in up to cubic anharmonic terms, and the coupling of the strain to the Mn $`e_g`$ orbital state by the Jahn-Teller coupling. In this paper, instead of representing the elastic energy in terms of spring constants between different ions as done in Refs. and , we will introduce a more general and, we hope, more useful formulation in which long wavelength lattice distortions are treated via conventional elastic theory, while the short wavelength atomic motions are treated explicitly.
### B General energy expressions for $`(\pi ,\pi ,0)`$ and $`(\pi ,\pi ,\pi )`$ distortions
The energy per Mn ion due to uniform strains, $`e_{xx}`$, $`e_{yy}`$, and $`e_{zz}`$, is most conveniently written in terms of the following combinations:
$`Q_{1u}`$ $`=`$ $`{\displaystyle \frac{a_0}{\sqrt{3}}}(e_{xx}+e_{yy}+e_{zz}),`$ (1)
$`Q_{2u}`$ $`=`$ $`{\displaystyle \frac{a_0}{\sqrt{2}}}(e_{xx}e_{yy}),`$ (2)
$`Q_{3u}`$ $`=`$ $`{\displaystyle \frac{a_0}{\sqrt{6}}}(2e_{zz}e_{xx}e_{yy}).`$ (3)
We have
$$\frac{E_u}{N_{Mn}}=\frac{1}{2}K_BQ_{1u}^2+\frac{1}{2}K^{}(Q_{2u}^2+Q_{3u}^2),$$
(4)
where the bulk modulus $`K_B=a_0(c_{11}+2c_{12})=3a_0c_B`$, the Jahn-Teller shear modulus $`K^{}=a_0(c_{11}c_{12})=2a_0c^{}`$, and $`c_{ij}`$ are the usual elastic constants.
In addition to these uniform strains, we consider staggered distortions with wave vector $`\stackrel{}{k}=(\pi ,\pi ,0)`$ or $`(\pi ,\pi ,\pi ).`$ We represent the amplitudes of the distortions by $`a_0\stackrel{}{\delta }_\stackrel{}{k}`$ for Mn ions and $`a_0u_\stackrel{}{k}^{x,y,z}`$ for O ions. Translational symmetry implies that the uniform distortion and the staggered distortion do not couple with each other up to the second order. Therefore, the harmonic elastic energy due to the staggered distortions simply adds to the harmonic uniform strain energy. In Appendix A, we use symmetry arguments to obtain the general forms of the elastic energies due to $`\stackrel{}{\delta }_\stackrel{}{k}`$ and $`u_\stackrel{}{k}^{x,y,z}`$ for $`(\pi ,\pi ,0)`$ and $`(\pi ,\pi ,\pi )`$ distortions. In the Jahn-Teller coupling energy $`E_{JT}`$, which we will introduce later, only the lattice distortions which have even parity about Mn ions appear. Therefore, $`E_{JT}`$ does not depend on the Mn-distortions $`\delta _\stackrel{}{k}^{x,y,z}`$ for either $`\stackrel{}{k}=(\pi ,\pi ,0)`$ or $`\stackrel{}{k}=(\pi ,\pi ,\pi )`$, or the z-direction oxygen distortion $`u_\stackrel{}{k}^z`$ for $`\stackrel{}{k}=(\pi ,\pi ,0)`$. The energy cost of the relevant distortions is most conveniently written in terms of
$`Q_{2s}`$ $`=`$ $`{\displaystyle \frac{a_0}{\sqrt{2}}}(v_{sx}v_{sy}),`$ (5)
$`Q_{3s}`$ $`=`$ $`{\displaystyle \frac{a_0}{\sqrt{6}}}(2v_{sz}v_{sx}v_{sy}),`$ (6)
where $`v_{sa}`$ is the $`(\pi \pi 0)`$ or $`(\pi \pi \pi )`$ amplitude of $`v_\stackrel{}{i}^a`$ $`=u_\stackrel{}{i}^au_{\stackrel{}{i}\widehat{a}}^a`$. The energy of a staggered distortion, $`E_s`$, depends upon the ordering wavevector and, for the distortions we consider, is
$`{\displaystyle \frac{E_s(\pi ,\pi ,0)}{N_{Mn}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}K_{2s}Q_{2s}^2+{\displaystyle \frac{1}{2}}K_{3s}Q_{3s}^2,`$ (7)
$`{\displaystyle \frac{E_s(\pi ,\pi ,\pi )}{N_{Mn}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}K_s\left(Q_{2s}^2+Q_{3s}^2\right),`$ (8)
where the $`K_{2s\text{}}K_{3s\text{}}`$and $`K_s`$ are the elastic constants defined in Appendix A, and arise mainly from the Mn-O bond stretching mode.
We compare the sizes of $`K^{}`$, $`K_{2s}`$, $`K_{3s}`$, and $`K_s`$. This comparison is important to determine the ground state distortions. A uniform strain changes more bonds than a staggered distortion. For example, within a harmonic nearest-neighbor approximation, $`Q_{2s}`$ and $`Q_{3s}`$ modes involve only the Mn-O bond, but $`Q_{2u}`$ and $`Q_{3u}`$ modes involve both Mn-O and Mn-Mn bonds. Therefore, uniform modes have larger elastic moduli than staggered modes, which remains true when reasonable further neighbor interactions are included. Therefore, we expect $`K^{}>K_{3s},`$ $`K_{2s}`$, $`K_s`$. Our analysis of the general expression of the elastic energy given in Appendix A shows $`K_{3s}`$ $`>K_{2s}`$.
The staggered lattice distortion is caused by the Jahn-Teller coupling to the Mn $`e_g`$ orbital state. The Mn $`e_g`$ electron state on site $`a_0\stackrel{}{i}`$ is represented by
$$|\theta _\stackrel{}{i}>=\mathrm{cos}\theta _\stackrel{}{i}|3z^2r^2>+\mathrm{sin}\theta _\stackrel{}{i}|x^2y^2>,$$
(9)
where $`|3z^2r^2>`$ and $`|x^2y^2>`$ are the two linearly independent $`e_g`$ orbitals on the site. The cases of interest here are two-sublattice distortions. We represent the orbital states of $`e_g`$ electrons on these two sublattices as $`\theta _1`$ and $`\theta _2`$. Then the JT energy is$`^\text{,}`$
$$\frac{E_{JT}}{N_{Mn}}=\sqrt{\frac{3}{2}}\lambda \frac{1}{2}\left[\mathrm{cos}2\theta _1(Q_{3u}+Q_{3s})+\mathrm{sin}2\theta _1(Q_{2u}+Q_{2s})+\mathrm{cos}2\theta _2(Q_{3u}Q_{3s})+\mathrm{sin}2\theta _2(Q_{2u}Q_{2s})\right].$$
(10)
To represent the energy only in terms of lattice distortions, we minimize the above Jahn-Teller energy with respect to the orbital states $`\theta _1`$ and $`\theta _2`$. We obtain
$`\mathrm{cos}2\theta _{1,2}^{\mathrm{min}}`$ $`=`$ $`{\displaystyle \frac{Q_{3u}\pm Q_{3s}}{\sqrt{(Q_{2u}\pm Q_{2s})^2+(Q_{3u}\pm Q_{3s})^2}}},`$ (11)
$`\mathrm{sin}2\theta _{1,2}^{\mathrm{min}}`$ $`=`$ $`{\displaystyle \frac{Q_{2u}\pm Q_{2s}}{\sqrt{(Q_{2u}\pm Q_{2s})^2+(Q_{3u}\pm Q_{3s})^2}}},`$ (12)
and
$$\frac{E_{JT}}{N_{Mn}}=\frac{1}{2}\sqrt{\frac{3}{2}}\lambda \left[\sqrt{(Q_{2u}+Q_{2s})^2+(Q_{3u}+Q_{3s})^2}+\sqrt{(Q_{2u}Q_{2s})^2+(Q_{3u}Q_{3s})^2}\right].$$
(13)
We also include the largest anharmonic energy, which is the one between the nearest neighbor Mn-O pair. It is given by
$$E_{anh}=\frac{4}{\sqrt{3}}Aa_0^3\underset{i,a}{}\left(\frac{e_{aa}}{2}+u_i^a\delta _i^a\right)^3+\left(\delta _i^au_{ia}^a\frac{e_{aa}}{2}\right)^3.$$
(14)
Total energy is the sum of the terms considered so far:
$$E_{tot}^{elastic}=E_u+E_s+E_{JT}+E_{anh}.$$
(15)
By minimizing $`E_{tot}^{elastic}`$, we find the distortions induced by the JT coupling and anharmonic energy terms, which will be discussed below for $`\stackrel{}{k}=(\pi ,\pi ,0)`$ and $`\stackrel{}{k}=(\pi ,\pi ,\pi )`$ distortions separately.
### C Energy minimization for $`(\pi ,\pi ,0)`$ distortion
We minimize $`E_u+E_s+E_{JT}`$ and then treat $`E_{anh}`$ as a perturbation, since we expect and will show below that the anharmonic term is small compared to the harmonic terms. We find that in the ground state of $`E_u+E_s+E_{JT}`$, the distortion mode which has the smallest modulus among $`Q_{2u}`$, $`Q_{3u}`$, $`Q_{2s}`$, and $`Q_{3s}`$ is non-zero, and all the other distortion modes vanish. Since, we found $`K^{}>K_{3s}>K_{2s},`$ the ground state of $`E_u+E_s+E_{JT}`$ is
$`Q_{3s}=Q_{2u}=Q_{3u}`$ $`=`$ $`0,`$ (16)
$`Q_{2s}`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2}}}{\displaystyle \frac{\lambda }{K_{2s}}},`$ (17)
$`E_{\mathrm{min}}`$ $`=`$ $`{\displaystyle \frac{3}{4}}{\displaystyle \frac{\lambda ^2}{K_{2s}}},`$ (18)
$`\theta _1,\theta _2`$ $`=`$ $`\pi /4,3\pi /4.`$ (19)
It is noteworthy that in this ground state the Mn lattice itself preserves cubic symmetry, and only oxygen ions make staggered distortions.
We next study how the anharmonic energy term changes the above ground state. We represent $`E_{anh}`$ in Eq. (14) by $`\stackrel{}{\delta }_\stackrel{}{k}`$, $`u_\stackrel{}{k}^{x,y,z}`$, and $`e_{xx,yy,zz}`$. Direct expansion (or symmetry argument) shows that each of $`\stackrel{}{\delta }_\stackrel{}{k}`$ and $`u_\stackrel{}{k}^z`$ appears only as the second order in $`E_{anh}`$, which implies that these distortions remain zero unless the coupled uniform strains exceed certain values and the lattice becomes unstable. After representing $`E_{anh}(\pi \pi 0)`$ in terms of $`Q_{1u,2u,3u,2s,3s}`$, the same argument shows $`Q_{3s}=Q_{2u}=0`$. By taking the largest remaining term, we obtain
$$\frac{E_{anh}(\pi \pi 0)}{N_{Mn}}AQ_{2s}^2(Q_{1u}\frac{1}{\sqrt{2}}Q_{3u}).$$
(20)
The total energy for $`\stackrel{}{k}=(\pi ,\pi ,0)`$ distortion which we will minimize to find strains in bulk state or in thin films is
$`E_{\mathrm{tot}}(\pi \pi 0)`$ $`=`$ $`(E_u+E_s+E_{JT}+E_{anh})/N_{\text{Mn}}`$ (21)
$`=`$ $`{\displaystyle \frac{1}{2}}K_BQ_{1u}^2+{\displaystyle \frac{1}{2}}K^{}Q_{3u}^2+{\displaystyle \frac{1}{2}}K_{2s}Q_{2s}^2\sqrt{{\displaystyle \frac{3}{2}}}\lambda \sqrt{Q_{2s}^2+Q_{3u}^2}+AQ_{2s}^2\left(Q_{1u}{\displaystyle \frac{1}{\sqrt{2}}}Q_{3u}\right).`$ (22)
In bulk state, there is no external constraint. When we minimize $`E_{\mathrm{tot}}(\pi \pi 0)`$ in a leading order in $`A`$ to find the lattice distortions in the bulk state, we obtain
$$E_{tot}^{MIN}(\pi \pi 0)=\frac{3}{4}\frac{\lambda ^2}{K_{2s}}A^2\frac{9}{16}\frac{\lambda ^4}{K_{2s}^4}\frac{2K^{}2K_{2s}+K_B}{K_B(K^{}K_{2s})}+O(A^4),$$
(23)
for
$`Q_{1u}`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{\lambda ^2}{K_{2s}}}{\displaystyle \frac{A}{K_{2s}^2}}+O(A^3),`$ (24)
$`Q_{2s}`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2}}}{\displaystyle \frac{\lambda }{K_{2s}}}+O(A^2),`$ (25)
$`Q_{3s}`$ $`=`$ $`Q_{2u}=0,`$ (26)
$`Q_{3u}`$ $`=`$ $`{\displaystyle \frac{3\lambda ^2A}{2\sqrt{2}(K^{}K_{2s})K_{2s}^2}}+O(A^3).`$ (27)
The results show that the observed uniform tetragonal distortion is due to the anharmonic term which couples staggered and uniform distortions. $`\mathrm{LaMnO}_3`$ expands upon heating, which implies $`A<0`$. Therefore, above result indicates $`Q_{3u}<0`$, which is consistent with the observed distortion in bulk $`\mathrm{LaMnO}_3`$. Since $`K^{}K_{2s}`$ is order of magnitude smaller than $`K_{2s}`$, $`Q_{3u}`$ is order of magnitude larger than $`Q_{1u}`$. Up to order of $`A`$, the anharmonic term does not change $`Q_{2s}`$.
We note “$`Q_{3u}`$ mode softening” : When we have lattice distortion $`Q_{2s}=\sqrt{3/2}\frac{\lambda }{K_{2s}}`$, the JT coupling term in Eq. (22), expanded about small $`Q_{3u}`$, effectively reduces $`Q_{3u}`$ mode modulus by $`K_{2s}`$. Therefore, when the anharmonic term induces $`Q_{3u}`$, the restoring spring constant is $`K^{}K_{2s}`$ rather than $`K^{}`$. This is the reason why we have $`K^{}K_{2s}`$ in the denominator of Eq. (27) and have a relatively large $`Q_{3u}`$($``$ 35 % of $`Q_{2s}`$ from crystallography data). Thus we expect the shear modulus corresponding to the $`Q_{3u}`$ distortion to be much smaller in $`\mathrm{LaMnO}_3`$ than in doped cubic manganates.
### D Energy minimization for $`(\pi ,\pi ,\pi )`$ distortion
By applying similar considerations to $`(\pi ,\pi ,\pi )`$ distortion, and using the condition $`K^{}>K_s`$, we find that the degenerate ground states of $`E_u+E_s(\pi \pi \pi )+E_{JT}(\pi \pi \pi )`$ are
$`Q_{2u}=Q_{3u}`$ $`=`$ $`0,`$ (28)
$`Q_{2s}`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2}}}{\displaystyle \frac{\lambda }{K_s}}\mathrm{sin}2\theta _1,`$ (29)
$`Q_{3s}`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2}}}{\displaystyle \frac{\lambda }{K_s}}\mathrm{cos}2\theta _1,`$ (30)
$`E_{\mathrm{min}}`$ $`=`$ $`{\displaystyle \frac{3}{4}}{\displaystyle \frac{\lambda ^2}{K_s}},`$ (31)
$`\theta _2`$ $`=`$ $`\pi \theta _1,`$ (32)
where $`\theta _1`$ is an arbitrary angle between 0 and $`\pi `$. After we include the same anharmonic energy and apply the same arguments used for the $`(\pi ,\pi ,0)`$ distortion, we find the total energy expression for the $`(\pi ,\pi ,\pi )`$ distortion which we will minimize further is
$`E_{tot}(\pi \pi \pi )`$ $`=`$ $`[E_u+E_s(\pi \pi \pi )+E_{JT}(\pi \pi \pi )+E_{anh}(\pi \pi \pi )]/N_{Mn}`$ (33)
$`=`$ $`{\displaystyle \frac{1}{2}}K_BQ_{1u}^2+{\displaystyle \frac{1}{2}}K^{}(Q_{2u}^2+Q_{3u}^2)+{\displaystyle \frac{1}{2}}K_s(Q_{2s}^2+Q_{3s}^2)`$ (36)
$`\sqrt{{\displaystyle \frac{3}{2}}}\lambda {\displaystyle \frac{1}{2}}\left[\sqrt{(Q_{2u}+Q_{2s})^2+(Q_{3u}+Q_{3s})^2}+\sqrt{(Q_{2u}Q_{2s})^2+(Q_{3u}Q_{3s})^2}\right]`$
$`+A\left[Q_{2s}^2(Q_{1u}{\displaystyle \frac{1}{\sqrt{2}}}Q_{3u})+Q_{3s}^2(Q_{1u}+{\displaystyle \frac{1}{\sqrt{2}}}Q_{3u})\sqrt{2}Q_{2s}Q_{3s}Q_{2u}\right].`$
When we minimize $`E_{tot}(\pi \pi \pi )`$, we obtain
$$E_{tot}^{MIN}(\pi \pi \pi )=\frac{3}{4}\frac{\lambda ^2}{K_s}A^2\frac{9}{16}\frac{\lambda ^4}{K_s^4}\frac{2K^{}2K_s+K_B}{K_B(K^{}K_s)}+O(A^4)$$
(37)
for
$`Q_{1u}`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{\lambda ^2}{K_s}}{\displaystyle \frac{A}{K_s^2}}+O(A^3),`$ (38)
$`Q_{2s}`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2}}}{\displaystyle \frac{\lambda }{K_s}}+O(A^2),`$ (39)
$`Q_{3s}`$ $`=`$ $`Q_{2u}=0,`$ (40)
$`Q_{3u}`$ $`=`$ $`{\displaystyle \frac{3\lambda ^2A}{2\sqrt{2}(K^{}K_s)K_s^2}}+O(A^3).`$ (41)
The other two physically equivalent distortions obtained by permuting $`x`$, $`y`$, and $`z`$ from above results are also degenerate ground states. The ground state energies for the ($`\pi \pi 0`$) and ($`\pi \pi \pi `$) distortions will be compared in Sect. V.
### E Energy minimization for strained films
In this paper, by uniaxial strain, we mean a tetragonal strain with axis perpendicular to film plane. This strain can be applied by growing epitaxial films on square lattice substrates with lattice parameters different from the $`xy`$ plane lattice parameter for bulk $`\mathrm{LaMnO}_3`$. For thin films with uniaxial strains, we calculate the lattice distortions in the following way: We assume that ($`\pi \pi 0`$) distortion pattern is favored even in strained films. We also assume perfect epitaxy; therefore $`e_{xx}`$ and $`e_{yy}`$ are determined by substrates. Then $`Q_{2s}`$ and $`e_{zz}`$ for given strains can be found by minimizing $`E_{tot}`$ about these distortions, which also give $`Q_{1u}`$ and $`Q_{3u}`$. From $`Q_{2s}`$ and $`Q_{3u}`$, we find the $`e_g`$ orbital states, $`\theta _1`$ and $`\theta _2`$, by
$`\mathrm{cos}(2\theta _1)`$ $`=`$ $`{\displaystyle \frac{Q_{3u}}{\sqrt{Q_{2s}^2+Q_{3u}^2}}},`$ (42)
$`\theta _2`$ $`=`$ $`\pi \theta _1.`$ (43)
Since the external strain changes $`Q_{2s}`$ and $`Q_{3u}`$, it also changes the JT splitting, given by $`\delta \mathrm{\Delta }E_{JT}=\sqrt{6}\lambda \sqrt{Q_{2s}^2+Q_{3u}^2}\sqrt{6}\lambda \sqrt{(Q_{2s}^{eq})^2+(Q_{3u}^{eq})^2}.`$ Parameters of the model, i.e., $`K_B`$, $`K^{}`$, $`K_{2s}`$, $`\lambda `$, and $`A,`$ are determined from experiments, as explained in Appendix B. The determined parameter values are shown in Table I.
## III Model of magnetic interaction and magnetostriction effects
### A Magnetic interaction
In this section, we present a model describing magnetism in bulk and strained films of LaMnO<sub>3</sub>. The superexchange interaction between Mn ions in $`\mathrm{LaMnO}_3`$ depends on the electron orbital overlap, particularly the overlap between Mn $`e_g`$ orbitals and O $`p`$ orbitals. It is also argued that $`t_{2g}`$ electrons also contribute antiferromagnetic interaction. We build a simple model below, which incorporates these two contributions of exchange interactions.
We calculate the superexchange interaction due to the $`e_g`$ electrons using a similar method as in Ref. . We assume $`e_g`$ spin is always parallel to $`t_{2g}`$ spin at each site due to the strong Hund’s coupling. For the two Mn ions, one at $`\stackrel{}{i}`$ and the other at $`\stackrel{}{i}+\widehat{z}`$, we find
$$J_{e_g}(\theta _1,\theta _2)=J_F\left(\mathrm{sin}^2\theta _1\mathrm{cos}^2\theta _2+\mathrm{cos}^2\theta _1\mathrm{sin}^2\theta _2\right),$$
(44)
where $`J_F`$ is a positive parameter of the model. This is, in fact, equivalent to the model in Ref. , if we assume that the state with two holes on the intermediate oxygen ion (which was considered in Ref. ) requires an infinite energy. Equation (44) shows that when one of $`\theta _1`$ and $`\theta _2`$ is zero (i.e., $`|3z^2r^2>`$ state) and the other is $`\pi /2`$ (i.e., $`|x^2y^2>`$ state), $`J_{e_g}(\theta _1,\theta _2)`$ is most ferromagnetic due to the maximized hopping between the filled orbital on one site and the empty orbital on the other site. When $`\theta _1=\theta _2=\pi /2`$, or $`\theta _1=\theta _2=0`$, $`J_{e_g}(\theta _1,\theta _2)`$ is least ferromagnetic, since the hopping between the filled and empty orbitals on the two sites vanishes. The $`t_{2g}`$ superexchange is expected to be independent of $`\theta _1`$ and $`\theta _2`$, and always antiferromagnetic, so we set $`J_{t_{2g}}=J_{AF}`$. Therefore, along the $`z`$ direction the total exchange interaction is
$`J_z`$ $`=`$ $`J_{t_{2g}}+J_{e_g}`$ (45)
$`=`$ $`J_{AF}J_F\left(\mathrm{sin}^2\theta _1\mathrm{cos}^2\theta _2+\mathrm{cos}^2\theta _1\mathrm{sin}^2\theta _2\right).`$ (46)
The sign of the total superexchange is determined by the competition between the $`e_g`$ ferromagnetism and $`t_{2g}`$ antiferromagnetism. Along $`x`$ and $`y`$ directions, proper rotations result in
$`J_x`$ $`=`$ $`J_{AF}J_F\left(\mathrm{sin}^2(\theta _1{\displaystyle \frac{2\pi }{3}})\mathrm{cos}^2(\theta _2{\displaystyle \frac{2\pi }{3}})+\mathrm{cos}^2(\theta _1{\displaystyle \frac{2\pi }{3}})\mathrm{sin}^2(\theta _2{\displaystyle \frac{2\pi }{3}})\right),`$ (47)
$`J_y`$ $`=`$ $`J_{AF}J_F\left(\mathrm{sin}^2(\theta _1+{\displaystyle \frac{2\pi }{3}})\mathrm{cos}^2(\theta _2+{\displaystyle \frac{2\pi }{3}})+\mathrm{cos}^2(\theta _1+{\displaystyle \frac{2\pi }{3}})\mathrm{sin}^2(\theta _2+{\displaystyle \frac{2\pi }{3}})\right).`$ (48)
Using the condition $`\theta _2=\pi \theta _1`$ obtained before, we find
$`J_x`$ $`=`$ $`J_yJ_{xy}=J_{xy}^0+J^{}\mathrm{cos}4\theta _1,`$ (49)
$`J_z`$ $`=`$ $`J_z^0+J^{}\mathrm{cos}4\theta _1,`$ (50)
where $`J_{xy}^0=J_{AF}5J_F/8`$, $`J_z^0=J_{AF}J_F/4`$, and $`J^{}=J_F/4`$. It explicitly shows how the magnetic coupling depends on the orbital states.
### B Magnetostriction effects
In this section, we present our model of magnetostriction effect in bulk LaMnO<sub>3</sub>. We expect magnetostriction effect for the following reason: Since $`J_x`$, $`J_y`$, and $`J_z`$ depend on the orbital states $`\theta _1`$ and $`\theta _2`$, once certain magnetic ordering occurs, the magnetic ordering, in turn, will change $`\theta _1`$ and $`\theta _2`$ to gain further magnetic energy. Through the JT coupling, this change in the orbital states can cause the change in the JT strains. Magnetism and lattice distortion are coupled through the Mn $`e_g`$ orbital degree of freedom.
To estimate the size of the magnetostriction effects, we add to Eq. (22) the term
$$\frac{E_{mag}}{N_{Mn}}=4\left(2|J_{xy}(\theta _1)|+|J_z(\theta _1)|\right),$$
(51)
which represents the mean-field T=0 magnetic energy, and find the changes in the orbital states. For A-type antiferromagnetic state, after dropping constant terms, we obtain
$$\frac{E_{mag}}{N_{Mn}}=4J^{}\mathrm{cos}4\theta _1.$$
(52)
By adding $`E_{mag}/N_{Mn}`$ to the total energy, we can find the extra structure at T=0 due to the magnetic order.
Using the Landau free energy method, we examine whether the magnetostriction effect makes the phase transition first order or not. From the energy gain due to the magnetic order at T=0, and the mean field estimate of $`T_c`$, we obtain the Landau free energy per site, $`f`$ :
$$f(\theta _1,Q_{1u},Q_{2s},Q_{3u},m)=2[TT_c(\theta _1)]m^2+T_c(\theta _1)m^4+E_{elastic}(\theta _1,Q_{1u},Q_{2s},Q_{3u}),$$
(53)
where $`m=M/M(T=0)`$ is the normalized magnetization and $`E_{elastic}(\theta _1,Q_{1u},Q_{2s},Q_{3u})`$ is the total elastic energy obtained in Sect. II. Magnetic ordering temperature, $`T_c(\theta _1)`$, can be obtained by molecular field theory:
$$T_c(\theta _1)=2\left(4|J_{xy}(\theta _1)|+2|J_z(\theta _1)|\right).$$
(54)
We can find the order of the phase transition and the size of the magnetostriction effect in the following approximation: Since $`Q_{1u}`$ is not directly coupled to $`\theta _1`$, we expect the change in $`Q_{1u}`$ is small, which can be seen from the numerical results in Table II. Therefore, we neglect $`Q_{1u}`$ dependence and obtain $`E_{elastic}(\theta _1,Q_{2s},Q_{3u})`$, which we minimize further about $`Q_{2s}`$ and $`Q_{3u}`$ to obtain $`E_{elastic}(\theta _1)`$. By expanding $`E_{elastic}(\theta _1)`$ about the minimum energy orbital state for $`m=0`$, $`\theta _1(m=0)`$, we obtain
$$E_{elastic}(\theta _1)\frac{1}{2}K_\theta \left(\theta _1\theta _1(m=0)\right)^2,$$
(55)
where $`K_\theta =12\lambda ^2(K^{}K_{2s})/(K_{2s}K^{})`$, and $`\theta _1(m=0)=\frac{1}{2}\mathrm{cos}^1[\sqrt{3}\lambda A/(K_{2s}(K^{}K_{2s}))]`$. By substituting $`T_c(\theta _1)T_c^0[1+\alpha (\theta _1\theta _1(m=0))]`$ and $`E_{elastic}(\theta _1)`$ in Eq. (53), and minimizing $`f(\theta _1,m)`$ about $`\theta _1`$, we obtain
$$\theta _1(m)\theta _1(m=0)=\frac{1}{K_\theta }\left[2T_c^0\alpha m^2T_c^0\alpha m^4\right],$$
(56)
$$f(m)=2(TT_c^0)m^2+\frac{T_c^0}{K_\theta }(K_\theta 2\alpha ^2T_c^0)m^4+O(m^6).$$
(57)
Therefore, we will have the second order phase transition when $`K_\theta >2\alpha ^2T_c^0`$, and the first order transition when $`K_\theta <2\alpha ^2T_c^0`$. This result implies that when the lattice is soft (small $`K_\theta `$) or the coupling between magnetic interaction and lattice is large (large $`\alpha `$), the transition becomes first-order. The estimate of the parameters for $`\mathrm{LaMnO}_3`$ in Sect. VI.C predicts the second-order phase transition. Above equation for $`\theta _1(m)`$ also gives the net change of $`\theta _1`$, $`\theta _1(T=0)\theta _1(T>T_c)=\alpha T_c^0/K_\theta `$.
## IV Model of band structure and optical conductivity
We use a tight-binding approximation to calculate band structure and optical conductivity for strained films. This method is explained in detail in Ref. , and summarized in this section. According to band theory calculations the conduction band is derived mainly from the Mn $`e_g`$ symmetric d-orbitals and is well separated from other bands. Therefore, we only consider Mn $`e_g`$ levels. Kinetic energy and chemical potential terms are
$$H_{\mathrm{KE}}+H_\mu =\frac{1}{2}\underset{\stackrel{}{i},\stackrel{}{\delta },a,b,\alpha }{}t_\stackrel{}{\delta }^{ab}d_{\stackrel{}{i}a\alpha }^{}d_{\stackrel{}{i}+\stackrel{}{\delta }b\alpha }+H.c.\mu \underset{\stackrel{}{i},a,\alpha }{}d_{\stackrel{}{i}a\alpha }^{}d_{\stackrel{}{i}a\alpha }.$$
(58)
Here $`\stackrel{}{i}`$ represents the coordinates of Mn sites, $`\delta (=\pm x,y,z)`$ labels the nearest neighbors of Mn sites, $`a`$ and $`b`$ represent the two degenerate Mn $`e_g`$ orbitals on a site, $`\alpha `$ denotes the spin state, and $`t_\stackrel{}{\delta }^{ab}`$ is the hopping amplitude between orbital $`a`$ on site $`\stackrel{}{i}`$ and $`b`$ on site $`\stackrel{}{i}+\stackrel{}{\delta }`$. We choose $`|\psi _1>=|3z^2r^2>`$ and $`|\psi _2>=|x^2y^2>`$ as the two linearly independent $`e_g`$ orbitals on each site as before. The hopping matrix $`t_\stackrel{}{\delta }^{ab}`$ has a special form: For hopping along $`z`$ direction, it connects only the two $`|3z^2r^2>`$ states, thus $`t_z^{ab}=t_z^{ab}=t_o`$ for a=b=1, and zero otherwise. The hopping matrices in other bond directions are obtained by appropriate rotations. The Jahn-Teller coupling for uniform $`Q_3`$ distortion and staggered $`Q_2`$ distortion with wave vector $`\stackrel{}{K}_{lattice}=(\pi ,\pi ,0)`$ or $`(\pi ,\pi ,\pi )`$ is
$$H_{\mathrm{JT}}=\sqrt{\frac{3}{2}}\lambda \underset{\stackrel{}{i},\alpha }{}\left(\begin{array}{c}d_{1,\stackrel{}{i},\alpha }^{}\\ d_{2,\stackrel{}{i},\alpha }^{}\end{array}\right)^T\left(\begin{array}{cc}Q_{3u}& \mathrm{exp}(i\stackrel{}{K}_{lattice}\stackrel{}{i})Q_{2s}\\ \mathrm{exp}(i\stackrel{}{K}_{lattice}\stackrel{}{i})Q_{2s}& Q_{3u}\end{array}\right)\left(\begin{array}{c}d_{1,\stackrel{}{i},\alpha }\\ d_{2,\stackrel{}{i},\alpha }\end{array}\right).$$
(59)
The Hund’s coupling for antiferromagnetic core spin configuration with wave vector $`\stackrel{}{K}_{spin}`$ ($`(0,0,\pi )`$ for A type AF, $`(\pi ,\pi ,\pi )`$ for purely AF) is
$$H_{\mathrm{Hund}}=J_\mathrm{H}S_\mathrm{c}\underset{\stackrel{}{i},a}{}\left[\left(1\mathrm{exp}(i\stackrel{}{K_{spin}}\stackrel{}{i})\right)d_{\stackrel{}{i},a,}^{}d_{\stackrel{}{i},a,}+\left(1+\mathrm{exp}(i\stackrel{}{K_{spin}}\stackrel{}{i})\right)d_{\stackrel{}{i},a,}^{}d_{\stackrel{}{i},a,}\right].$$
(60)
The total Hamiltonian is the sum of the terms considered so far. By diagonalizing this in $`k`$ space, we can find the energy levels and eigenstates.
Optical conductivity per volume due to the transitions between Mn $`e_g`$ levels, $`\sigma `$, can be calculated using the eigenstates and energy levels found from the above Hamiltonian. Using the standard linear response theory,$`^\text{,}`$ optical conductivity is given by
$$\sigma _p^{\lambda \nu }=\frac{1}{i\omega N_{Mn}a_0^3}\underset{n}{}\frac{<0|J_{p\lambda }^{}|n><n|J_{p\nu }|0>}{\mathrm{}\omega (E_nE_0)+iϵ},$$
(61)
where $`ϵ`$ is an infinitesimal and $`J_p`$ is given by $`\widehat{\stackrel{}{J_p}}=\frac{iea_0}{2\mathrm{}}_{\stackrel{}{i},\stackrel{}{\delta },a,b,\alpha }t_\stackrel{}{\delta }^{ab}\stackrel{}{\delta }(d_{\stackrel{}{i}a\alpha }^{}d_{\stackrel{}{i}+\stackrel{}{\delta }b\alpha }H.c.)`$.
## V Comparison of $`(\pi ,\pi ,0)`$ and $`(\pi ,\pi ,\pi )`$ type distortions
In the purely local model considered in Ref. , for which only nearest-neighbor Mn-O and Mn-Mn springs are considered, we have $`K_s`$ = $`K_{2s}`$ and $`E_{tot}^{MIN}(\pi \pi \pi )=E_{tot}^{MIN}(\pi \pi 0)`$ up to order $`A^2`$ according to Eqs. (23) and (37). However, in real situation for which farther ion-ion elastic energies exist, we have $`K_s`$ $``$ $`K_{2s}`$ and the two distortions have different energies. Equations (23) and (37) shows that if $`K_{2s}<K_s`$, then the leading order term stabilizes the $`(\pi ,\pi ,0)`$ distortion over the $`(\pi ,\pi ,\pi )`$ distortion. If $`K_s`$ and $`K_{2s}`$ are very close to $`K^{}`$, the $`A^2`$ term will be smaller for the larger of $`K_s`$ and $`K_{2s}`$, opposite trend of the leading order term. For parameter values determined in Appendix B, we find that $`K_{2s}<K_s`$ is required to make $`(\pi \pi 0)`$ distortion favored. Optic phonon spectrum along $`\stackrel{}{k}=\kappa (\pi ,\pi ,0)`$ and $`\kappa (\pi ,\pi ,\pi )`$ ($`0<\kappa <1`$ ) would be useful to check this condition. At $`\kappa =0`$, these two modes have the same energy. As $`\kappa `$ approaches to 1, if the $`(\pi ,\pi ,\pi )`$ mode has a higher energy than the $`(\pi ,\pi ,0)`$ mode, it would be an indication that $`K_s>K_{2s}`$. So far, the phonon spectrum for $`\mathrm{LaMnO}_3`$ has not been calculated. Ghosez et al. have calculated phonon spectra for similar compounds, $`\mathrm{BaTiO}_3`$, $`\mathrm{PbTiO}_3`$, and $`\mathrm{PbZrO}_3`$. These results show that $`(\pi ,\pi ,0)`$ mode ($`M_\text{2}`$ point in Fig. 1 of Ref. ) has a higher energy than $`(\pi ,\pi ,\pi )`$ mode ($`R_{\text{12’}}`$ point in Fig.1. of Ref. ) by 11 %, 7%, and 2 % respectively, contrary to our expectation for $`\mathrm{LaMnO}_3`$. Since the energies of these modes depend sensitively on transition metal elements, a phonon spectrum calculation for $`\mathrm{LaMnO}_3`$ is necessary.
We examine the possibility that the two distortions have different band energies. For this purpose, we use a tight binding Hamiltonian introduced in the previous section, and calculate the band structures and total band energies for the two distortion patterns. On the first Brillouin zone boundary (i.e., $`|k_z|=\pi /2`$ or $`|k_x|+|k_y|=\pi `$ planes) and on the planes satisfying $`|k_x|=|k_y|`$, we find that the band structures for $`(\pi ,\pi ,0)`$ and $`(\pi ,\pi ,\pi )`$ distortions are identical. Between these planes, when a band is well separated from other bands, it has a similar band structure for the $`(\pi ,\pi ,0)`$ and $`(\pi ,\pi ,\pi )`$ distortions. Since the filled bands are well separated from the empty bands by the Jahn-Teller splitting, the results show that the change of the filled bands are negligible. We find that the total band energy per Mn ion changes less than 1 meV between the $`(\pi ,\pi ,0)`$ and $`(\pi ,\pi ,\pi )`$ distortions. We therefore conclude that the strain (lattice restoring force) effects are crucial.
In our calculation for the strained films, we assume that the sign of the energy difference between the two types of ordering is not changed by applied strains, and consider the $`(\pi ,\pi ,0)`$ ordering only.
## VI Results
### A Lattice and orbital states in strained films
In this section, we present our calculations for strained films. We use $`e_{||}`$ to denote the substrate induced additional strain, $`e_{xx}e_{xx}^{bulk}`$. We examine the strain in the range of -2 % $`<`$ $`e_{||}`$ $`<`$ 2 %. Since $`a_04.03\stackrel{}{\text{A}}`$, the difference of the $`xy`$ plane lattice parameter between bulk $`\mathrm{LaMnO}_3`$ and substrate is between $`0.08`$ $`\stackrel{}{\text{A}}`$ and 0.08 $`\stackrel{}{\text{A}}`$. We represent the changes in $`Q_{3u}`$, $`Q_{2s}`$, and $`e_{zz}`$ by writing $`\delta `$ in front. For example, $`\delta Q_{3u}=Q_{3u}Q_{3u}^{bulk}`$. Our calculation shows that $`\delta e_{zz}`$ versus $`e_{||}`$ in this range is close to linear. We obtain $`\delta e_{zz}`$/$`e_{||}`$ $``$ -1.8 for the parameter set obtained from Ref. , and $`\delta e_{zz}`$/$`e_{||}`$ $``$ -1.4 for the parameter set obtained from Ref. . These ratios are about 2 times larger than the ratio for $`\mathrm{La}_{0.7}\mathrm{Ca}_{0.3}\mathrm{MnO}_3`$ film. This is due to the softening of the $`Q_{3u}`$ mode in $`\mathrm{LaMnO}_3`$.
From $`\delta e_{zz}`$ versus $`e_{||}`$, we can also find $`\delta Q_{1u}`$ and $`\delta Q_{3u}`$ versus $`e_{||}`$. The results are shown in Fig. 2. To understand these results, we find the leading terms of $`\delta Q_{1u}`$/($`a_0e_{||}`$), $`\delta Q_{2s}`$/($`a_0e_{||}`$), and $`\delta Q_{3u}`$/($`a_0e_{||}`$). They are
$`{\displaystyle \frac{\delta Q_{1u}}{a_0e_{||}}}`$ $`=`$ $`2\sqrt{3}{\displaystyle \frac{K^{}K_{2s}}{K_B+2\left(K^{}K_{2s}\right)}}+O(A^2),`$ (62)
$`{\displaystyle \frac{\delta Q_{2s}}{a_0e_{||}}}`$ $`=`$ $`{\displaystyle \frac{3}{\sqrt{2}}}{\displaystyle \frac{K_B\left(2K^{}3K_{2s}\right)+4\left(K^{}K_{2s}\right)^2}{\left(K_B+2K^{}2K_{2s}\right)\left(K^{}K_{2s}\right)K_{2s}^2}}\lambda A+O(A^3),`$ (63)
$`{\displaystyle \frac{\delta Q_{3u}}{a_0e_{||}}}`$ $`=`$ $`\sqrt{6}{\displaystyle \frac{K_B}{K_B+2\left(K^{}K_{2s}\right)}}+O(A^2).`$ (64)
This shows that $`Q_{2s}`$ changes more slowly than $`Q_{3u}`$, because the staggered distortion is coupled to the uniform strain only through the anharmonic term. Since $`K^{}K_{2s}`$ is almost one order smaller than $`K_B`$, $`Q_{1u}`$ also changes slowly compared with $`Q_{3u}`$. Therefore, the main effect of the uniaxial strain is the change in the uniform tetragonal distortion without much change in staggered distortion or volume. Figure 2 shows that $`\delta Q_{2s}`$/($`a_0e_{||}`$) can be either positive or negative depending on the parameter values, whereas $`\delta Q_{1u}`$/($`a_0e_{||}`$) $`>`$ 0 and $`\delta Q_{3u}`$/($`a_0e_{||}`$) $`<`$ 0 always.
We obtain $`\theta _1`$ 54.8 $`\stackrel{}{\text{A}}`$ for the bulk state from Eq. (42). $`\theta _1`$ versus $`e_{||}`$ is shown in Fig. 3. The change in $`\theta _1`$ is about $`\pm 5^o15^o`$. For tensile strains, $`\theta _1`$ and $`\theta _2`$ approach towards $`90^o`$, which corresponds to $`|x^2y^2>`$. For compressive strains, $`\theta _1`$ and $`\theta _2`$ approach to $`0^o`$ and $`180^o`$, which corresponds to $`|3z^2r^2>`$. This can be understood from the fact that in $`\theta _1`$,$`\theta _2`$=$`|x^2y^2>`$ state, the x-y plane Mn-O-Mn distance tends to be farthest, and in $`\theta _1`$, $`\theta _2`$ = $`|3z^2r^2>`$, shortest due to the electron distribution in the $`xy`$ plane.
The change in the Jahn-Teller splitting, $`\delta \mathrm{\Delta }E_{JT}`$, is shown in Fig. 4, which is about $`\pm `$ 0.02-0.2 eV for $`\pm `$ 2 % strain.
### B Magnetic property in strained films
According to Ref. , $`J_{xy}`$ = -1.66 meV, $`J_z`$=1.16 meV for bulk $`\mathrm{LaMnO}_3`$. Since $`\theta _1=54.8^o`$ for bulk, we obtain $`J_F=7.52`$ meV and $`J_{AF}=4.50`$ meV, which corresponds to $`J_{xy}^0=0.22`$ meV, $`J_z^0=2.62`$ meV, and $`J^{}=1.88`$ meV. Using these parameter values, we find $`J_{xy}`$ and $`J_z`$ versus $`\theta _1`$, which are plotted in Fig. 5 (a). As the two orbital states approach to $`\pi /4`$ and $`3\pi /4`$, i.e., $`|z^2x^2>`$ and $`|z^2y^2>`$, both $`J_{xy}`$ and $`J_z`$ become more ferromagnetic, and as they approach to 0 and $`\pi `$, i.e., $`|3z^2r^2>`$ and $`|3z^2r^2>`$, or to $`\pi /2`$ and $`\pi /2`$, i.e., $`|x^2y^2>`$ and $`|x^2y^2>`$, $`J_{xy}`$ and $`J_z`$ become less ferromagnetic. This is due to the orbital-state-dependent hopping between the filled and empty orbitals, which mediates ferromagnetic interaction. When $`20^o<\theta _1<70^0`$, the magnetic ground state remains A-type antiferromagnetic. Outside this range, both $`J_{xy}`$ and $`J_z`$ become positive, and purely antiferromagnetic state is the ground state. In fact, Fig. 3 shows that 2 % tensile strain can change $`\theta _1`$ close to $`70^o`$, and turn the material into a purely antiferromagnetic state.
The mean field estimates of $`T_c`$ are shown in Fig. 5 (b). $`T_c`$ for bulk state is about 210 K, somewhat larger than the measured $`T_c`$ = 140 K. It shows that $`\pm `$ 2 % strain changes $`T_c`$ by about $`\pm `$ 50 K. Relatively large change in $`T_c`$ and the possible change into purely AF state are due to the strong dependence of magnetism on the $`e_g`$ orbital state and the strong Jahn-Teller coupling between the $`e_g`$ orbital state and the lattice distortion.
### C Magnetostriction effects in bulk state
In this subsection we present the magnetostriction effects calculated by the model in Section III. B. We use the two sets of parameter values in Table I and $`J^{}=`$ 1.88 meV. The results obtained by numerical minimization of the Landau free energy \[Eq. (53)\] are shown in Table II. The change in the JT strain, $`\delta ϵ^{}`$, is about 0.003 $``$ 0.01. $`\delta ϵ^{}/[ϵ^{}(T>T_c)]`$ is about -0.08 for the parameter set from Ref. , and -0.31 for that from Ref. . These results show that when the effective JT modulus of $`Q_{3u}`$ mode, $`K^{}K_{2s}`$, is smaller, the magnetostriction effect is larger.
We obtain $`\alpha `$2, and $`2\alpha ^2T_c^0`$ 0.1 eV for both parameter sets from Refs. and . We obtain $`K_\theta `$ =1.2 eV for Ref. , and $`K_\theta `$ =0.3 eV for Ref. . Therefore, the transition will be of the second order. However, if we have a softer $`K_\theta `$ (a third or a tenth), or a stronger magnetostriction coupling $`\alpha `$ (twice or three times), then we will have a first-order phase transition. Our numerical minimization of the free energy confirms these results.
Recently, the orbital ordering in $`\mathrm{LaMnO}_3`$ has been directly observed using a resonant X-ray scattering technique. In this result, the orbital ordering versus temperature curve has a change of the slope at $`T=T_N`$. The sign of the change indicates that the orbital states change away from $`3x^2r^2/3y^2r^2`$ ($`\theta _1`$=$`60^o`$) as $`T0`$ below $`T_N`$, which is consistent with our calculation. Recent neutron diffraction study measured $`\mathrm{sin}\theta _1`$ versus $`T`$. Our results predict that $`\mathrm{sin}\theta _1`$ (which corresponds to $`c_2`$ in Fig. 3 in Ref. ) changes by 0.01 for parameters from Ref. and 0.04 for parameters from Ref. between T=0 and T $`>T_N`$. However, Ref. shows negligible change in $`\mathrm{sin}\theta _1`$ between T=0 and T $`>T_N`$, which indicates a smaller magnetostriction coupling $`\alpha `$, or a larger elastic modulus $`K_\theta `$ than the values obtained above.
### D Band structures in strained films
First, we summarize the band structure in bulk state A-type antiferromagnetic $`\mathrm{LaMnO}_3,`$ which is explained in detail in Ref. . Crudely speaking, the bands fall into 4 pairs, which may be understood by setting $`t_o=0`$ \[as occurs at $`\stackrel{}{k}=(\pi /2,\pi /2,\pi /2)`$\]; in this case we have four separate energy levels on each site, which are $`E_{1,2}=\sqrt{\frac{3}{2}}\lambda \sqrt{Q_{2s}^2+Q_{3u}^2}`$, $`E_{3,4}=\sqrt{\frac{3}{2}}\lambda \sqrt{Q_{2s}^2+Q_{3u}^2}`$, $`E_{5,6}=2J_\mathrm{H}S_\mathrm{c}\sqrt{\frac{3}{2}}\lambda \sqrt{Q_{2s}^2+Q_{3u}^2}`$, and $`E_{7,8}=2J_\mathrm{H}S_\mathrm{c}+\sqrt{\frac{3}{2}}\lambda \sqrt{Q_{2s}^2+Q_{3u}^2}`$. To find the three parameter values of our model Hamiltonian, i.e., $`t_0`$, $`\lambda `$, and $`J_\mathrm{H}S_\mathrm{c}`$, we fit our band structure calculation to the LDA (local density-functional approximation) band calculation for the JT distorted $`\mathrm{LaMnO}_3`$ in Ref. at high symmetry points in reciprocal space. The standard deviation is $``$ 0.2 eV. The determined parameter values are $`t_o`$=0.622 eV, $`\lambda `$=1.38 eV/$`\stackrel{}{\text{A}}`$, and $`2J_\mathrm{H}S_\mathrm{c}`$=2.47 eV. The fitted band structure is shown in Fig. 1 in Ref. .
When the strain does not change A-type antiferromagnetic core spin configuration, the main effect of the strain is the change in the band width. The results are shown in Fig. 6. Solid lines are for the compressive strain, and dashed lines are for the tensile strain. Bulk band structure can be approximately obtained by taking the average of the two band structures. At $`(\pi /2,\pi /2,\pi /2)`$, where the effective hopping vanishes, the energy level change is $`\delta \mathrm{\Delta }E_{JT}`$ obtained before. From this point, dispersions along $`(\pi /2,\pi /2,0)`$ and along $`(\pi ,0,\pi /2)`$ represent the hoppings in $`z`$ direction and in $`xy`$ plane, respectively. Between $`(\pi /2,\pi /2,\pi /2)`$ and $`(\pi /2,\pi /2,0)`$, the widths of the lower JT bands, $`E_{1,2}`$ and $`E_{5,6}`$, are increased for compressive strains, whereas the widths of the upper JT bands, $`E_{3,4}`$ and $`E_{7,8}`$, are decreased. This is related to the changes in $`\theta _1`$ and $`\theta _2`$ due to the strains: As $`\theta _1`$ and $`\theta _2`$ approach to 0 and $`\pi `$, the lower JT level approaches to $`3z^2r^2`$ state which has a large hopping along $`z`$ direction, whereas the upper JT level approaches to $`x^2y^2`$ which has no hopping along $`z`$ direction. But between $`(\pi /2,\pi /2,\pi /2)`$ and $`(\pi ,0,\pi /2)`$, the dispersion does not change much, indicating that the average hopping is not changed much in $`xy`$ plane due to the alternating orbital pattern in $`xy`$ plane.
As pointed out in the previous section, 2 % tensile strain may induce purely antiferromagnetic ground state. Due to the strong Hund’s coupling, this results in substantial changes in band structure and optical conductivity as shown below and in the next section, respectively. Band structure for the purely AF state can be obtained by using $`\stackrel{}{K}_{spin}=(\pi ,\pi ,\pi )`$ in the model described in Sect. IV. The results are shown in Fig. 7 for the same lattice distortions and parameters used for the band structure shown as dotted lines in Fig. 6. Between $`(\pi /2,\pi /2,\pi /2)`$ and $`(\pi /2,\pi /2,0)`$, the two band structures are identical since it involves only $`z`$ direction hopping. Between $`(\pi /2,\pi /2,\pi /2)`$ and $`(\pi ,0,\pi /2)`$, where only $`xy`$ directional hopping is involved, for A-type AF state the different JT levels repel each other, while for purely AF state the different Hund’s levels repel each other. This represents different mixings for different spin configurations: the mixing for A-type AF state is mainly between the different JT states, while the mixing for purely AF state is mainly between the different Hund’s states. the band structure between $`(\pi ,0,0)(0,0,0)(\pi /2,\pi /2,\pi /2)`$ shows the same trend, which results in a small indirect band gap for the purely AF state. However, due to the on-site Coulomb repulsion of about 1-2 eV neglected in the above calculation it is unlikely to have the insulator-to-semimetal transition in this material.
### E Optical conductivities in strained films
From the band structure, we have calculated optical conductivities for strained films in A-type AF ground state. Results are shown in Fig. 8. $`\sigma _{xx}`$ or $`\sigma _{yy}`$ (solid lines) shows a relatively small changes by strains. The spectral weight of the Hund’s peak in $`\sigma _{zz}`$ (dotted lines) at around 2.5 eV is increased (decreased) about 20 % as we apply 2 % compressive (tensile) strain. This seems consistent with the change of the average hopping along $`z`$ direction by strains mentioned in the previous section.
Optical conductivities for the purely AF state are calculated in the same way. Figure 9 shows the results for $`\sigma _{xx}`$ (solid line) and $`\sigma _{zz}`$ (dotted line). First, the sharp Hund’s peak in $`\sigma _{zz}`$ is disappeared. This can be understood from the band structure, particularly between $`(\pi /2,\pi /2,\pi /2)`$ and $`(\pi ,0,\pi /2)`$. The sharp Hund’s peak in $`\sigma _{zz}`$ for A-type AF state originates from the two parallel bands split by the Hund’s coupling, 2$`J_\text{H}S_\text{c}`$. However, this structure disappears when we have purely AF core spin as seen in Fig. 7. Comparison of Fig. 9 (a) (2 % strain) and (b) (bulk) shows that as orbital state is changed toward $`x^2y^2`$ by the tensile strain, which has zero hopping along z-direction, the spectral weight of the Hund’s peak in $`\sigma _{zz}`$ decreases, as observed for the A-type AF state in Fig. 8. $`\sigma _{xx}`$ shows prominent peaks at the Hund’s splitting (2.5-3.5 eV) and at the Hund-plus-JT splitting (4-5 eV) due to purely AF spin state in contrast to the A-type AF state (see Fig. 8). The JT peaks at around 1 eV in both $`\sigma _{xx}`$ and $`\sigma _{zz}`$ are due to the strong hybridization between major and minor spin states. If we increase $`J_HS_c`$ and reduce this hybridization, the JT peaks decrease, as can be seen by comparing Figs. 9(a) and 9(c). Above results show that changes in $`xy`$ plane spin configuration make differences in $`\sigma _{zz}`$ due to the changes in hybridization.
For the calculations so far, we have assumed zero on-site Coulomb repulsion $`U`$. However, as explained in Ref. , in this material there exists $`U1.6`$ eV. Therefore actual peak positions will be higher by $``$ 1.6 eV and the spectral weight of each peak will be reduced inversely proportional to the peak energy.
## VII Conclusion
In summary, we developed a model of elastic energy for $`\mathrm{LaMnO}_3`$ and solved for uniaxial strains in thin films. We found that $`\pm 2`$ % strain can change the uniform tetragonal strain and $`e_g`$ orbital states without much change in the staggered distortion or volume. We found that 2 % tensile strain can change the magnetic ground state into purely antiferromagnetic state, inducing dramatic changes in band structure and optical conductivity. Magnetostriction effect at $`T_\text{N}`$ in bulk state is found to be large. We examined the possibility that the lattice energy will favor ($`\pi \pi 0`$) ordering over ($`\pi \pi \pi `$) ordering. We also noted “$`Q_{3u}`$ mode softening” in $`\mathrm{LaMnO}_3`$. The results presented in this paper for $`\mathrm{LaMnO}_3`$ show the strong coupling between lattice, Mn $`e_g`$ orbital state, and exchange interaction, which lies at the root of the novel properties of doped manganese perovskites.
This work is supported by NSF-DMR-9705182 and the University of Maryland MRSEC.
## A General expression of elastic energy for staggered distortion
In this Appendix, we show how we can get the general expression of the elastic energy due to distortions with wave vector $`\stackrel{}{k}`$ in perovskite structure. We again consider three dimensional displacement of Mn ion, $`\stackrel{}{\delta }`$, and the displacements of oxygen ions along Mn-Mn axis, $`u^x`$, $`u^y`$, and $`u^z`$ as defined in the text and shown in Fig. 1. We consider the displacements with wave vector $`\stackrel{}{k}`$, $`\stackrel{}{\delta }_\stackrel{}{i}=\stackrel{}{\delta }_\stackrel{}{k}e^{i\stackrel{}{k}\stackrel{}{i}}`$ and $`u_\stackrel{}{i}^{x,y,z}=u_\stackrel{}{k}^{x,y,z}e^{i\stackrel{}{k}\stackrel{}{i}}`$. When we define $`d_{1,2,3}(\stackrel{}{k})=\delta _\stackrel{}{k}^{x,y,z}`$ and $`d_{4,5,6}(\stackrel{}{k})=u_\stackrel{}{k}^{x,y,z}`$, the energy due to these strains, $`E_s(\stackrel{}{k})`$, is given by
$$E_s(\stackrel{}{k})=\underset{i,j}{}d_i(\stackrel{}{k})D_{ij}(\stackrel{}{k})d_j(\stackrel{}{k}),$$
(A1)
where $`D_{ij}(\stackrel{}{k})=D_{ji}(\stackrel{}{k})`$.
In special cases, symmetry arguments make certain terms vanish or equal. First, when $`k_x=\pi `$, mirror symmetry operation about $`x=0`$ plane changes $`k_x`$ to $`k_x`$. Since $`k_x=\pi `$ and $`k_x^{}=\pi `$ are equivalent, the symmetry operation changes $`\stackrel{}{k}`$ back to $`\stackrel{}{k}`$. Since this operation changes $`\delta _\stackrel{}{k}^x`$ to $`\delta _\stackrel{}{k}^x`$, odd order terms of $`\delta _\stackrel{}{k}^x`$ should vanish, and therefore, $`D_{12}(\stackrel{}{k})=D_{13}(\stackrel{}{k})=D_{14}(\stackrel{}{k})=D_{15}(\stackrel{}{k})=D_{16}(\stackrel{}{k})=0`$. Second, when $`k_x=0`$, mirror symmetry operation about $`x=0`$ plane changes $`\stackrel{}{k}`$ back to $`\stackrel{}{k}`$, $`\delta _\stackrel{}{k}^x`$ to $`\delta _\stackrel{}{k}^x`$, and $`u_\stackrel{}{k}^x`$ to $`u_\stackrel{}{k}^x`$. Therefore, odd order terms of $`\delta _\stackrel{}{k}^x`$ and $`u_\stackrel{}{k}^x`$ vanish. Therefore, $`D_{12}(\stackrel{}{k})=D_{13}(\stackrel{}{k})=D_{15}(\stackrel{}{k})=D_{16}(\stackrel{}{k})=0`$ and $`D_{42}(\stackrel{}{k})=D_{43}(\stackrel{}{k})=D_{45}(\stackrel{}{k})=D_{46}(\stackrel{}{k})=0`$. Third, when $`k_x=\pm k_y`$, mirror operation interchanging $`x`$ axis and $`\pm y`$ axis changes $`\stackrel{}{k}`$ back to $`\stackrel{}{k}`$, and $`\delta _x`$ to $`\pm \delta _y`$, $`u_x`$ to $`\pm u_y`$. Therefore, $`D_{11}(\stackrel{}{k})=D_{22}(\stackrel{}{k})`$, $`D_{44}(\stackrel{}{k})=D_{55}(\stackrel{}{k})`$, $`D_{14}(\stackrel{}{k})=D_{25}(\stackrel{}{k})`$, $`D_{15}(\stackrel{}{k})=D_{24}(\stackrel{}{k})`$, $`D_{1j}(\stackrel{}{k})=\pm D_{2j}(\stackrel{}{k})`$, $`D_{4j}(\stackrel{}{k})=\pm D_{5j}(\stackrel{}{k})`$, where $`j=3,6`$
If we apply these rules to $`\stackrel{}{k}=(\pi ,\pi ,0)`$, and $`(\pi ,\pi ,\pi )`$ distortions, we obtain the following expressions.
$`{\displaystyle \frac{E_s[\stackrel{}{k}=(\pi ,\pi ,0)]}{N_{Mn}}}`$ $`=`$ $`D_{11}(\stackrel{}{k})(\delta _{\stackrel{}{k}}^{x}{}_{}{}^{2}+\delta _{\stackrel{}{k}}^{y}{}_{}{}^{2})+D_{33}(\stackrel{}{k})\delta _{\stackrel{}{k}}^{z}{}_{}{}^{2}+D_{44}(\stackrel{}{k})(u_{\stackrel{}{k}}^{x}{}_{}{}^{2}+u_{\stackrel{}{k}}^{y}{}_{}{}^{2})`$ (A3)
$`+D_{66}(\stackrel{}{k})u_{\stackrel{}{k}}^{z}{}_{}{}^{2}+2D_{45}(\stackrel{}{k})u_\stackrel{}{k}^xu_\stackrel{}{k}^y+2D_{36}(\stackrel{}{k})\delta _\stackrel{}{k}^zu_\stackrel{}{k}^z.`$
$`{\displaystyle \frac{E_s[\stackrel{}{k}=(\pi ,\pi ,\pi )]}{N_{Mn}}}`$ $`=`$ $`D_{11}(\stackrel{}{k})(\delta _{\stackrel{}{k}}^{x}{}_{}{}^{2}+\delta _{\stackrel{}{k}}^{y}{}_{}{}^{2}+\delta _{\stackrel{}{k}}^{z}{}_{}{}^{2})+D_{44}(\stackrel{}{k})(u_{\stackrel{}{k}}^{x}{}_{}{}^{2}+u_{\stackrel{}{k}}^{y}{}_{}{}^{2}+u_{\stackrel{}{k}}^{z}{}_{}{}^{2})`$ (A5)
$`+2D_{45}(\stackrel{}{k})(u_\stackrel{}{k}^xu_\stackrel{}{k}^y+u_\stackrel{}{k}^yu_\stackrel{}{k}^z+u_\stackrel{}{k}^zu_\stackrel{}{k}^x).`$
$`K_{2s}`$, $`K_{3s}`$, and $`K_s`$ in the text are defined as $`K_{2s}=[D_{44}(\pi \pi 0)D_{45}(\pi \pi 0)]/2`$, $`K_{3s}=3[D_{44}(\pi \pi 0)+D_{45}(\pi \pi 0)]/2`$, and $`K_s=[D_{44}(\pi \pi \pi )D_{45}(\pi \pi \pi )]/2`$. Therefore, the condition for $`K_{3s}>K_{2s}`$ is $`D_{44}(\pi \pi 0)>2D_{45}(\pi \pi 0)`$. We expect this condition is well-satisfied for the following reasons: First, the stability of lattice implies $`D_{44}(\pi \pi 0)>0`$. Second, the Coulomb repulsion between oxygen will favor similar O-O distances, which implies $`D_{45}(\pi \pi 0)>0`$.
## B Parameters for the model of elastic energy
Optical absorption experiment for $`\mathrm{LaMnO}_3`$ shows Mn-O bond stretching mode peak at 70.3 meV. From this we can find the effective Mn-O spring constant $`K_1`$=7.36 eV/$`\stackrel{2}{\stackrel{}{\text{A}}}`$ (Ref. ). If we assume that $`K_1`$ is the main contribution to $`K_{2s}`$, we obtain $`K_{2s}K_1`$/2 = 3.68 eV/$`\stackrel{2}{\stackrel{}{\text{A}}}`$. In Ref. , we used this value of $`K_1`$ to estimate the Jahn-Teller coupling constant $`\lambda `$, the value of which is consistent with LDA band calculation.
For $`\mathrm{La}_{0.83}\mathrm{Sr}_{0.17}\mathrm{MnO}_3`$, $`c_{11}`$ and $`c_{12}`$ have been measured in an ultrasound experiment. We believe the bulk modulus $`c_B=(c_{11}+2c_{12})/3`$ =143 GPa($`K_B`$=10.8 eV/$`\stackrel{2}{\stackrel{}{\text{A}}}`$) of $`\mathrm{La}_{0.83}\mathrm{Sr}_{0.17}\mathrm{MnO}_3`$ at 200 K (orthorhombic structure) can be used as an approximate value for the bulk modulus of $`\mathrm{LaMnO}_3`$. The value of $`c_B`$ at 310 K (rhombohedral phase) is about 176 GPa, which gives a rough estimate of uncertainty in $`c_B`$ of about 20 %.
Since $`\mathrm{La}_{0.83}\mathrm{Sr}_{0.17}\mathrm{MnO}_3`$ is in the doping region where a structural change happens, and $`c^{}`$ is sensitive to the structural transition unlike bulk modulus, we do not believe the measured $`c^{}`$ of this material is a good estimate of $`c^{}`$ for $`\mathrm{LaMnO}_3`$. Indeed, $`c^{}`$ for $`\mathrm{La}_{0.83}\mathrm{Sr}_{0.17}\mathrm{MnO}_3`$ is much smaller than those for other typical perovskite oxides. For example, the LDA calculation in Ref. predicted $`c_B=`$ 199 GPa, $`c^{}=`$ 142 GPa, and $`c^{}/c_B=`$ 0.71 for $`\mathrm{SrTiO}_3`$, which is close to the measurements for $`\mathrm{SrTiO}_3`$, $`c_B`$ =179 GPa, $`c^{}`$=115 GPa, and $`c^{}/c_B`$ = 0.64. The same LDA calculation results showed that other perovskite oxides, such as BaTiO<sub>3</sub>, CaTiO<sub>3</sub>, KNbO<sub>3</sub>, NaNbO<sub>3</sub>, PbTiO<sub>3</sub>, PbZrO<sub>3</sub>, BaZrO<sub>3</sub>However, for $`\mathrm{La}_{0.83}\mathrm{Sr}_{0.17}\mathrm{MnO}_3`$, $`c_B`$=143 GPa, $`c^{}`$=48 GPa , and $`c^{}/c_B`$=0.34 at 200 K and $`c_B`$=176 GPa, $`c^{}`$=35 GPa , and $`c^{}/c_B`$=0.20 at 300 K. Therefore $`c^{}/c_B`$ ratio is less than half of the values for typical perovskites.
Therefore, we instead use the $`c_{12}/c_{11}`$ ratio measured in $`\mathrm{La}_{0.7}\mathrm{Ca}_{0.3}\mathrm{MnO}_3`$ thin film (Ref. ) and $`\mathrm{La}_{0.6}\mathrm{Sr}_{0.4}\mathrm{MnO}_3`$ thin film (Ref. ), since their doping ranges are relatively far from the structural phase transition doping ratios. The results are $`c_{12}/c_{11}`$=0.312 for Ref. and $`c_{12}/c_{11}`$=0.374 for Ref. . Using
$$\frac{c^{}}{c_B}=\frac{3}{2}\frac{1c_{12}/c_{11}}{1+2c_{12}/c_{11}},$$
(B1)
we obtain $`c^{}`$ and $`K^{}`$ shown in Table I.
We obtain $`Q_{2s}^{eq}`$ and $`Q_{3u}^{eq}`$ for bulk $`\mathrm{LaMnO}_3`$ from crystallography data: $`Q_{2s}^{eq}`$=0.398 $`\stackrel{}{\text{A}}`$, $`Q_{3u}^{eq}`$=-0.142 $`\stackrel{}{\text{A}}`$ (Ref. ). Therefore, the three unknown quantities of the model, $`Q_{1u}^{eq}`$, $`\lambda `$, and $`A`$, are determined from the equilibrium condition:
$$\frac{E}{Q_{1u}}|_{eq}=0,\frac{E}{Q_{2s}}|_{eq}=0,\frac{E}{Q_{3u}}|_{eq}=0.$$
(B2)
Obtained parameter values are shown in Table I. The values of $`A`$ are small enough to justify our approximation of anharmonic terms. For example, the largest dropped term, $`AQ_{3u}^3/(3\sqrt{2})`$, is 0.5-2 % of $`K^{}Q_{3u}^2/2`$ for parameters in Table I. We obtain $`Q_{1u}^{\text{eq}}=`$ 0.024 $`\stackrel{}{\text{A}}`$ and 0.005 $`\stackrel{}{\text{A}}`$ for the parameters from Refs. and , respectively, which shows that the average bond length does not change much. Therefore, it is reasonable to approximate $`a_0`$ by the average Mn-Mn distance observed in bulk $`\mathrm{LaMnO}_3`$, 4.03 $`\stackrel{}{\text{A}}`$. The $`\lambda `$ values are close to the independent estimate $`\lambda `$=1.38 eV/$`\stackrel{}{\text{A}}`$ obtained from band structure fitting. Negative sign of $`A`$ is consistent with thermal expansion. It is noteworthy that the size of $`A`$ is largely different for the two parameter sets, which is the consequence of the $`Q_{3u}`$ mode softening: a small difference in $`K^{}`$ gives a quite large difference in $`K^{}K_{2s}`$ (0.89 eV/$`\stackrel{2}{\stackrel{}{\text{A}}}`$and 0.19 eV/$`\stackrel{2}{\stackrel{}{\text{A}}}`$for the parameters from Refs. and , respectively), which results in a large difference in the estimate of $`A`$.
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# Photon acceleration in variable ultra-relativistic outflows and high-energy spectra of Gamma-Ray Bursts.
## 1 Introduction
The most common model of gamma-ray burst (GRB) sources involves a relativistic outflow in which shocks occur and radiate away a fraction of the bulk kinetic energy. For typical model parameters the synchrotron emission peaks around 1 keV in the comoving frame, or $`100`$ keV in the observer frame. Generally, this is considered to be the primary spectrum observed. However, in internal shocks in the neighborhood of the flow photosphere, and also in external shocks in some cases (e.g. Madau, Blandford, & Rees 2000), the shocks can have a non-negligible Thomson scattering depth, which results in upscattering of these primary photons. Single scattering on individual shock-accelerated electrons with Lorentz factors $`\gamma _e300`$ produces photons with energies $`\gamma _e^2`$ keV in the comoving, or $`\mathrm{\Gamma }\gamma _e^210\mathrm{\Gamma }_2`$ GeV in the observer frame, where $`\mathrm{\Gamma }=100\mathrm{\Gamma }_2`$ is the bulk Lorentz factor.
Here we concentrate on a different, multiple scattering component. This is associated with mildly relativistic motions of different ejecta shells or turbulent cells resulting from multiple shock interactions, which contain more energy than the shock-accelerated highly relativistic electrons. Multiple interacting shells are naturally expected in internal shocks, and also in external shocks when a longer lasting modulated outflow runs into the first decelerated shell (e.g. Fenimore & Ramirez-Ruiz 2000; Kumar & Piran 2000).
Repeated scatterings using the energy of these bulk motions boost the photon energy through the equivalent of the Fermi acceleration mechanism of particles (Blandford & Payne 1981). This is related to Thompson’s (1994) photon scattering off Alvfén waves, but our mechanism relies instead on relative bulk motions. It differs also in using synchrotron photons instead of thermal photons as its source term, and hence leads to different characteristic energies. This bulk Comptonization results in a spectrum extending at least up to $``$ MeV in the comoving frame and $`100\mathrm{\Gamma }_2`$ MeV in the observer frame. The spectral power or luminosity per decade can increase as steeply as linearly in the photon energy. This provides a natural explanation for those GRB spectra (e.g. Preece et al 1999) which show a positive $`\nu F_\nu `$ slope above the MeV range (generally $`+1`$), which cannot be explained by direct synchrotron radiation from Fermi shock-accelerated electrons.
The component made up of bulk-scattered photons can extend up to a maximum observed energy $`100\mathrm{\Gamma }_2`$ MeV. Beyond this energy, Klein-Nishina and electron recoil effects set in, and the spectrum reverts to being dominated by the seed spectrum (with negative or flat power law slope) of the unscattered photons above the synchrotron peak.
In this work we adopt a test photon model, which assumes no back reaction on the plasma. Several potentially important effects are neglected (upscattered photons can heat electrons in the colliding shells and produce pairs, light pressure ensures that the total comoving energy of the scattered photons cannot exceed the total kinetic energy of the relative shell motions). These effects will be investigated elswehere (Gruzinov, Mészáros , & Rees 2000). The simplified approach that we use here retains the essential properties of the bulk motion comptonization phenomenon, and allows us to explore the main qualitative features it introduces in the spectra.
In §2 we specify the GRB model, in §3 we give an analytical model of Fermi acceleration of photons by colliding shell, and in §4 we describe our Monte Carlo simulations. The results are discussed and related to current and future observations in §5.
## 2 The GRB internal shock model
As a specific example to illustrate the effect, we restrict ourselves here to the internal shock model. We use a standard set of parameters for the ultra-relativistic outflow model of GRBs (e.g. Mészáros & Rees 2000), i.e. luminosity $`L=10^{52}L_{52}`$ erg/s, $`L_{52}1`$; terminal Lorentz factor $`\mathrm{\Gamma }100`$; variability time scale $`t_v=10^3t_{ms}`$ s, with $`t_{ms}1`$. With these parameters the shock between shells ejected at typical time intervals $`t_v`$ with $`\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Gamma }`$ occurs at a radius $`R\mathrm{\Gamma }^2ct_v`$. To determine the optical depth to Thomson scattering, we estimate the proper density of the wind as
$$n\frac{L}{4\pi R^2\mathrm{\Gamma }^2m_pc^3}.$$
(1)
The proper radial width of colliding shells is $`R/\mathrm{\Gamma }`$, and the optical depth is
$$\tau n\sigma _TR/\mathrm{\Gamma }\frac{L_{52}}{t_{ms}}\left(\frac{200}{\mathrm{\Gamma }}\right)^5.$$
(2)
We see that optical depths of $`\text{ }>0.1`$ are natural. In fact, values of $`\tau 1`$ are not purely coincidental, since in the standard GRB model the parameters are chosen so as to make $`\tau <1`$ enabling a non-thermal spectrum, and the “preferred” model uses smaller values of $`\mathrm{\Gamma }`$ (which do not require overly low baryon loads), thus making $`\tau \text{ }<1`$. The observational fits to the models provide reasonable support for such a choice of parameters (e.g. van Paradijs, Kouveliotou & Wijers, 2000).
Most of the seed photons are emitted at the synchrotron peak. Assuming mildly relativistic internal shocks, and magnetic field and electron energies equal to a fraction $`\xi _B`$ and $`\xi _e`$ of their equipartition values, the synchrotron peak frequency in the comoving frame is
$$h\nu \left(\frac{\xi _e}{0.1}\right)^2\left(\frac{\xi _B}{0.1}\right)^{1/2}\frac{L_{52}^{1/2}}{t_{ms}}\left(\frac{200}{\mathrm{\Gamma }}\right)^3\mathrm{keV}.$$
(3)
## 3 Analytical model of photon acceleration by colliding shells
The basic features of the bulk Comptonization process can be understood by means of a simple analytical model, which reproduces the essence of the Monte Carlo simulation discussed in §4. This model is sufficient to show that the high-energy slope of $`\nu F_\nu `$ can be positive, and gives a value for the cut off energy.
Consider a photon between two shells of equal optical depth $`\tau 1`$ (Figure 1). The shells collide with an ultra-relativistic relative Lorentz factor $`\gamma 1`$ <sup>1</sup><sup>1</sup>1$`\gamma =\frac{1}{2}\left(\frac{\mathrm{\Gamma }_A}{\mathrm{\Gamma }_B}+\frac{\mathrm{\Gamma }_B}{\mathrm{\Gamma }_A}\right)`$. The case of $`\gamma 1`$ requires a more cumbersome analysis. At our accuracy, we just assume that the ultra-relativistic solution is a useful approximation in the mildly relativistic regime.. Initially, in the frame comoving with shell B, the photon frequency is $`\nu `$ and the incidence angle is 0. We assume the Thomson regime, $`h\nu m_ec^2`$. The probability that the photon passes through shell B without scattering is $`1\tau `$. The photon scatters forward with probability $`\tau /2`$, and it scatters backward with probability $`\tau /2`$. In the frame comoving with shell A, the backward scattered photon has a small incidence angle. The average (over the scattering kernel) of the Lorentz-transformed frequency is
$$\nu =\frac{_{\pi /2}^\pi 𝑑\theta \mathrm{sin}\theta (1+\mathrm{cos}^2\theta )(1\mathrm{cos}\theta )}{_{\pi /2}^\pi 𝑑\theta \mathrm{sin}\theta (1+\mathrm{cos}^2\theta )}\gamma \nu =\frac{25}{16}\gamma \nu .$$
(4)
Thus, after one scattering the total energy of photons is changed by a factor of $`(25/32)\tau \gamma `$, and the frequency is changed by a factor of $`(25/16)\gamma `$. Then the observed spectrum is a power law with the luminosity per frequency octave $`\nu F_\nu \nu ^\beta `$, with
$$\beta =\frac{\mathrm{log}(25\tau \gamma /32)}{\mathrm{log}(25\gamma /16)}.$$
(5)
The slope is positive if $`\tau \gamma >1.28`$. The steepest slope is $`\beta =1`$, which is achieved in the limit $`\gamma \tau ^1`$. This model is similar to that for the ultrarelativistic isotropic multiple scattering (c.f. Rybicki & Lightman 1979, p. 212) where the slope is $`\beta =1+\mathrm{log}(\tau ^{})/\mathrm{log}(\gamma ^2)`$. In the last equation we should use $`\tau ^{}\tau `$ and $`\gamma ^{}=\sqrt{(1+\gamma )/2}`$, which is the Lorentz factor of the colliding shells in the center of mass frame.
A finite convergence time of the two geometrically thin shells does not introduce a limit on the maximum photon energy, as can be seen from the equivalent of Zeno’s paradox (since the photons always travel faster than the shells approach, the number of photon reflections can in principle grow arbitrarily large for a shrinking separation). However, the effective number of scatterings should be limited, introducing a corresponding spectral cut off at high energies, due to (i) a gradual Klein-Nishina decline of the cross section; (ii) a sharp cut-off of the energy of scattered photons due to electron recoil (the maximal energy of a photon backward scattered from shell B observed in frame A is $`\gamma m_ec^2`$); (iii) radiation pressure back reaction effects (if $`\beta >0`$, the total energy of the scattered photons builds up and the photon pressure can prevent the collision of the shells).
## 4 Monte Carlo simulations
Monte Carlo simulations of the photon acceleration model sketched in Figure 1 were performed in order to check the simple estimates of §3. These confirm that the analytical results, including equation (5), are a useful approximation in the mildly relativistic regime. They also allow us to find in more detail the shape of the Compton recoil / Klein-Nishina cut off (see Figure 2). The Monte Carlo simulations were done using two approximations, which simplify the problem without affecting significantly the result: polarization effects are neglected, and multiple scattering during one passage through a shell is neglected.
Initially, a seed synchrotron spectrum is released. The seed spectrum is isotropic in the center of mass frame, rising at low energies ($`\beta =4/3`$), with a smooth transition to a spectrum decreasing at high energies (where we take as an example $`\beta =0.5`$).
In Figure 2 we show a spectral calculation for two shells similar to those of Figure 1, using a relative Lorentz factor $`\gamma =5`$. The unscattered seed synchrotron spectrum is marked “$`\tau =0`$”. The spectrum of all photons that finally escape outwards (to the right) through the shell B, as measured in the (comoving) frame B, is shown in Figure 2 for different values of $`\tau `$. The Compton recoil / Klein-Nishina effects are noticeable at about 1 MeV comoving energies, and there is a sharp Compton recoil cut off at $`h\nu /(m_ec^2)=\gamma `$. In the observer frame this spectrum would be blueshifted by the factor $`\mathrm{\Gamma }_B`$.
The dependence on the relative Lorentz factor is illustrated in Figure 3, showing the effect for a fiducial relative Lorentz factor $`\gamma =2`$, which may also be characteristic of the reverse shock in an external deceleration shock scenario. As expected, the steepening remains present but becomes less strong as weaker shocks are considered. The cut-off, however, is inherent to the scattering physics, and although it has a tendency to decrease somewhat with decreasing $`\gamma `$, its order of magnitude remains in the neighborhood of $`h\nu m_ec^21`$ MeV (comoving frame).
## 5 Discussion
The Fermi acceleration of photons in the standard GRB fireball model has several consequences of theoretical and observational interest. The first is that it can naturally produce $`\nu F_\nu `$ spectra which are harder than the input synchrotron spectrum, including the possibility of an increasing $`\nu F_\nu `$ above the usual break found in the Band parameterization of spectra (Band et al 1999, Preece et al 1999). The latter is a property that is hard to obtain with a synchrotron model. Such rising $`\nu F_\nu \nu ^\beta `$ are expected, in this model, to have $`\beta \text{ }<1`$, and this implies that in some GRB most of the energy is at energies $`h\nu 100(\mathrm{\Gamma }/10^2)`$ MeV, well above the BATSE instrument band on the Compton Gamma Ray Observatory (CGRO). The current observational situation is that BATSE finds approximately 16 % of the spectra to be rising at $`\text{ }>1`$ MeV, and the observed rise is not faster than $`\beta =1`$ (Preece et al 1999). These fits generally cut off above 1.8 MeV, and the break is usually not much lower than this, so there is some uncertainty. The COMPTEL instrument on CGRO, sensitive up to 30 MeV, has analyzed $`30`$ bursts (Schoenfelder et al 2000), and an analysis of the slopes indicates in several cases $`\nu F_\nu `$ slopes $`\beta 0`$, with one burst of $`\beta 0.5`$ (Kippen et al 1999). The EGRET experiment on CGRO has detected $`30`$ bursts with the scintillation counters in the 1-200 MeV range, and $``$ 7 bursts with the spark chambers in the 100 MeV-30 GeV range. In this range the spectra are largely noise-dominated (Schaefer et al 1998, Bromm & Schaefer 1999). However, the scintillation spectral slopes (Catelli, Dingus & Schneid, 1997) are compatible with $`\beta \text{ }<0`$, although some could be positive and others negative, which is compatible with the analysis of spark chamber data (e.g. Sommers, 1994; Hurley et al, 1994). Large area detectors such as GLAST should be able to obtain more definite answers in the 20 MeV-300 GeV range.
Another implication of the photon acceleration described here is that it provides a natural mechanism to increase the efficiency of conversion of baryon bulk motion into photon energy. This is of interest since in general the internal shock synchrotron efficiency for radiating in the BATSE band is limited to 1-10% (Kumar 2000; Spada, Panaitescu & Mészáros 2000; see however Fenimore & Ramirez-Ruiz 2000, and observation-based estimates by Freedman & Waxman 2000). The increase in the radiative efficiency is simply given by the increase in the value of $`\nu F_\nu `$ at different energies, or by the integral $`F_\nu 𝑑\nu `$ in the range of interest. In the generic examples shown, this increase is substantial. Of course, the spectra shown in Figures 2 and 3 are test photon spectra, which do not take into account the back-reaction of radiation. The latter can become important when a substantial fraction of the bulk energy has been converted to radiation through Fermi acceleration, and the natural limit for this effect can be estimated as $`50\%`$ radiative efficiency.
The calculations presented here are meant to illustrate the consequences of photon acceleration by bulk motions. A self-consistent calculation is needed in order to explore the back-reaction of the radiation pressure and pair formation on the shell dynamics. Pair formation has an angle averaged cross section $`\sigma _{\gamma \gamma }\sigma _T/8`$ at a threshold which is similar to that for Klein-Nishina effects. It is not very important in low compactness situations (shocks at large radii), while in high compactness cases the exponential tail of the photon distribution would lead to a pair cascade, which lowers the cut off in the comoving spectra of Figures 2,3 to $`\text{ }<`$ 0.5 MeV, possibly with some pile-up of photons at this energy (Gruzinov, Mészáros & Rees, 2000).
###### Acknowledgements.
We thank V. Connaughton, B. Dingus, P. Kumar and M.J. Rees for valuable comments. AG was supported by the W. M. Keck Foundation and NSF PHY-9513835. PM was supported by NASA NAG5-2857, the Guggenheim Foundation and the Institute for Advanced Study.
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# 1 Introduction
## 1 Introduction
Once the quantum algebras and quantum groups had been introduced , , , it was obvious to apply the same techniques also in the super case. Indeed, at present there exist a considerable number of papers on quantum supergroups. Most of them consider deformations of the general linear supergroup $`\mathrm{GL}(m|n)`$ , , but more recently there also appeared a paper by Lee and Zhang , in which a deformation of $`\mathrm{OSP}(1|2n)`$ is discussed, and a paper by Zhang , dealing with a deformation of $`\mathrm{OSP}(2|2n)`$.
Quantum supergroups and quantum superalgebras are dual objects. This implies that there are (at least) two methods to introduce a quantum supergroup. One may regard the quantum supergroup as the primary object and define it by means of the $`R`$–matrix approach, thus generalizing the methods of Ref. to the super case. On the other hand, one may also start from the quantum superalgebra and construct the (algebra of functions on the) quantum supergroup as a sub–Hopf–superalgebra of the finite dual of the former. In Manin’s paper the first approach is used, while Zhang , , prefers to use the second. It is believed that both methods should lead to the same objects, however, since it is notoriously difficult to put the duality of quantum (super)groups and quantum (super)algebras on a firm basis (including a proof that the dual pairing is non–degenerate(!)), there is still a lot to be done to substantiate this belief (for example, see Ref. ; this paper contains an extensive list of references on quantum supergroups and quantum superalgebras).
The present paper is a direct sequel to Ref. , in which I have calculated the $`R`$–matrix of the quantum superalgebra $`U_q(𝔰𝔭𝔬(2n|2m))`$ in the vector representation. Here, I am going to use this $`R`$–matrix to introduce the corresponding quantum supergroup $`\text{SPO}_q(2n|2m)`$ and to construct some of its comodule superalgebras, among which is an $`\text{SPO}_q(2n|2m)`$–covariant quantum Weyl superalgebra. This will be done by means of techniques which (in the non–super case) have been described in Ref. (i.e., we use the first of the two approaches described above). However, as we are going to see, in the super case some amendments are necessary in order to cope with the nototious problem that the representations are not necessarily completely reducible. The reason for considering the $`\text{SPO}_q(2n|2m)`$ quantum supergroup rather than $`\mathrm{OSP}_q(2m|2n)`$ has been explained in Ref. .
The present work falls into two parts. The first half (Sections 2 to 4) contains results which are applicable to arbitrary quantum supergroups (not just the symplecto–orthogonal ones). In fact, these results are formulated in the general graded framework, where $`\mathrm{\Gamma }`$ in an arbitrary abelian group and $`\sigma `$ is a commutation factor on $`\mathrm{\Gamma }`$. In the second half (Sections 5 and 6), the special case of the $`\text{SPO}_q(2n|2m)`$ quantum supergroups is considered.
More precisely, this paper is set up as follows. In Section 2 we are going to recall some elementary facts about matrix elements of graded representations of a $`\sigma `$–bialgebra or $`\sigma `$–Hopf algebra. These facts will be used to motivate the definitions and constructions of the subsequent sections. In particular, they serve to find the appropriate sign factors (i.e., commutation factors) by which the formulae of the present paper differ from those in the non–graded setting.
Usually, a quantum (super)group is constructed in two steps. First, one introduces the corresponding bi–(super)algebra $`A(R)`$, then one amends this definition to obtain the quantum (super)group itself, which is a Hopf (super)algebra. Basic to the construction of $`A(R)`$ is a finite–dimensional (graded) vector space $`V`$ and some (even) linear operator $`R`$ in $`VV`$ which, in the more advanced part of the theory, is assumed to be invertible and to satisfy the (graded) Yang–Baxter equation. The $`\sigma `$–bialgebras $`A(R)`$ are discussed to some extent in Section 3. In particular, we prove (under the foregoing additional assumptions on $`R`$) that $`A(R)`$ is coquasitriangular.
Section 4 is devoted to the discussion of graded comodules and comodule algebras over $`A(R)`$. We concentrate on those of these structures that are derived from the vector comodule or its dual.
In Section 5 we introduce the quantum supergroup $`\text{SPO}_q(2n|2m)`$ and derive some of its properties. In particular, we show that $`\text{SPO}_q(2n|2m)`$ is coquasitriangular, and that there exists a natural dual pairing of the Hopf superalgebras $`\text{SPO}_q(2n|2m)`$ and $`U_q(𝔰𝔭𝔬(2n|2m))`$. (According to the general folklore, both of these results are expected to be true.)
Section 6 contains a discussion of some graded $`\text{SPO}_q(2n|2m)`$–comodule superalgebras. After describing the general setting, we concentrate on the construction of the $`\text{SPO}_q(2n|2m)`$–covariant quantum Weyl superalgebras.
The final Section 7 contains a discussion of our results and some hints to further research.
We close this introduction by some comments on our conventions and notation. Throughout the present work, the base field will be an arbitrary field $`𝕂`$ of characteristic zero. The pairing of a vector space and its dual will be denoted by pointed brackets, and the action of an algebra $`A`$ on an $`A`$–module will be denoted by a dot. It should be obvious that we are going to work in the general setting of graded algebraic structures (see Ref. ). To remind the reader of this fact, the graded tensor product of linear maps and of graded algebras will be denoted by $`\overline{}`$. However, it might be worthwhile to stress that even if the overbar seems to be missing this does not indicate that we have returned to the non–graded setting.
## 2 Some elementary properties of matrix elements of representations
In the present section we recall some elementary properties of the matrix elements of a graded representation of a $`\sigma `$–Hopf algebra (or $`\sigma `$–bialgebra). These properties will serve to motivate the definitions and constructions given later.
First of all, we introduce some notation. Throughout the present section, $`\mathrm{\Gamma }`$ will be an abelian group, and $`\sigma `$ will be a commutation factor on $`\mathrm{\Gamma }`$ with values in $`𝕂`$. All gradations appearing in this section are understood to be $`\mathrm{\Gamma }`$–gradations. In the following, we shall freely use the notation and results of Ref. .
Let $`V`$ be a finite–dimensional graded vector space, let $`V^{\mathrm{gr}}`$ be the graded dual of $`V`$, and let $`\text{Lgr}(V)`$ be the space of all linear maps of $`V`$ into itself that can be written as a sum of homogeneous linear maps. In Ref. , we have denoted the latter space by $`\text{Lgr}(V,V)`$. Since $`V`$ is finite–dimensional, it is known that, regarded as vector spaces, $`V^{\mathrm{gr}}`$ is equal to $`V^{}`$, the dual of $`V`$, and $`\text{Lgr}(V)`$ is equal to the space of all linear maps of $`V`$ into itself. We use the present notation in order to stress that $`V^{\mathrm{gr}}`$ is a graded vector space, and that $`\text{Lgr}(V)`$ is a graded associative algebra.
Suppose now that $`(e_i)_{iI}`$ is a homogeneous basis of $`V`$, where $`I`$ is some index set, and that $`e_i`$ is homogeneous of degree $`\eta _i\mathrm{\Gamma }`$. Let $`(e_i^{})_{iI}`$ be the basis of $`V^{\mathrm{gr}}`$ that is dual to $`(e_i)_{iI}`$, and let $`(E_{ij})_{i,jI}`$ be the basis of $`\text{Lgr}(V)`$ that canonically corresponds to $`(e_i)_{iI}`$:
$$E_{ij}(e_k)=\delta _{jk}e_i\text{for all }i,j,kI.$$
Then $`e_i^{}`$ is homogeneous of degree $`\eta _i`$, and $`E_{ij}`$ is homogeneous of degree $`\eta _i\eta _j`$.
It is well-known that the class of $`\mathrm{\Gamma }`$–graded vector spaces, endowed with the usual tensor product of graded vector spaces and with the twist maps defined by means of $`\sigma `$, forms a tensor category (see Ref. for details). A (generalized) Hopf algebra $`H`$ living in this category is called a $`\sigma `$–Hopf algebra. More explicitly, $`H`$ is a $`\mathrm{\Gamma }`$–graded associative algebra with a unit element, and it is endowed with a coproduct $`\mathrm{\Delta }`$, a counit $`\epsilon `$, and an antipode $`S`$, which satisfy the obvious axioms (in the category). In particular, this implies that the structure maps $`\mathrm{\Delta }`$, $`\epsilon `$, and $`S`$ are homogeneous of degree zero. Needless to say, $`\sigma `$–bialgebras are introduced similarly. (For generalized Hopf algebras living in more general categories, see Ref. .)
Let $`H`$ be a $`\sigma `$–Hopf algebra, and let $`H^{}`$ be the set of all linear forms on $`H`$ that vanish on a graded two–sided ideal of $`H`$ of finite codimension. This definition generalizes the definition given in Sweedler’s book to the general graded case considered here.
It is easy to see that the properties derived in Ref. can immediately be extended to the present setting. In particular, this is true for the criteria characterizing the elements of $`H^{}`$. We don’t want to go into detail here, but mention some properties not given there.
First of all we note that in the preceding definition the graded two–sided ideals can be replaced by graded left or right ideals (indeed, every graded left or right ideal of finite codimension contains a graded two–sided ideal of finite codimension). Secondly, the elements of $`H^{}`$ can also be characterized by the following property:
A linear form $`f`$ on $`H`$ belongs to $`H^{}`$ if and only if there exists a finite–dimensional graded (left) $`H`$–module $`V`$, an element $`xV`$, and a linear form $`x^{}V^{\mathrm{gr}}`$ such that
$$f(h)=x^{},hx\text{for all }hH.$$
This property shows very clearly that $`H^{}`$ is closely related to the theory of finite–dimensional graded representations of $`H`$.
Obviously, $`H^{}`$ is a graded subspace of $`H^{\mathrm{gr}}`$, the graded dual of $`H`$. Actually, much more is true, for the structure of $`H`$ as a $`\sigma `$–Hopf algebra leads to a similar structure on $`H^{}`$. More precisely, the transpose of $`\mathrm{\Delta }`$ induces the multiplication in $`H^{}`$, the counit $`\epsilon `$ is the unit element of $`H^{}`$, the transpose of the product map of $`H`$ induces the coproduct $`\mathrm{\Delta }_{}`$ of $`H^{}`$, the evaluation at the unit element of $`H`$ is the counit $`\epsilon _{}`$ of $`H^{}`$, and the transpose of $`S`$ induces the antipode $`S_{}`$ of $`H^{}`$. This $`\sigma `$–Hopf algebra $`H^{}`$ is called the finite (or continuous) dual of $`H`$.
Using the notation introduced at the beginning of this section, let us now assume in addition that $`V`$ is a graded $`H`$–module, and that $`\pi :H\text{Lgr}(V)`$ is the graded representation afforded by it. If $`i,jI`$, we define the linear form $`\pi _{ij}`$ on $`H`$ by
$$\pi _{ij}(h)=e_i^{},\pi (h)e_j\text{for all }hH.$$
Obviously, $`\pi _{ij}`$ is an element of $`H^{}`$, and is homogeneous of degree $`\eta _j\eta _i`$. Moreover, we have
$$\pi (h)=\underset{i,jI}{}\pi _{ij}(h)E_{ij}\text{for all }hH.$$
Then our first important observation is that
$$\mathrm{\Delta }_{}(\pi _{ij})=\underset{kI}{}\sigma (\eta _k\eta _i,\eta _k\eta _j)\pi _{ik}\pi _{kj}$$
(2.1)
for all $`i,jI`$. By definition of the antipode, this implies that
$$\underset{kI}{}\sigma (\eta _k\eta _i,\eta _k\eta _j)S_{}(\pi _{ik})\pi _{kj}=\underset{kI}{}\sigma (\eta _k\eta _i,\eta _k\eta _j)\pi _{ik}S_{}(\pi _{kj})=\delta _{i,j}1_H^{},$$
(2.2)
for all $`i,jI`$, where we have used the obvious fact that
$$\epsilon _{}(\pi _{ij})=\delta _{ij}\text{for all }i,jI.$$
Let us next recall that $`V`$ has a canonical structure of a graded right $`H^{}`$–comodule. It is uniquely fixed by the requirement that it induces, in a canonical way, the original structure of a graded left $`H`$–module on $`V`$. More precisely, let
$$\delta :VVH^{}$$
be the structure map of $`V`$, regarded as a graded right $`H^{}`$–comodule, let $`xV`$ be homogeneous of degree $`\xi `$, and set
$$\delta (x)=\underset{^r}{}x_r^0x_r^1,$$
with $`x_r^0V`$ and $`x_r^1H^{}`$. Then, if $`hH`$ is homogeneous of degree $`\eta `$, we have
$$hx=\pi (h)x=\sigma (\eta ,\xi )\underset{^r}{}x_r^0x_r^1,h,$$
where $`,`$ denotes the dual pairing of $`H^{}`$ and $`H`$. Explicitly, we find
$$\delta (e_i)=\underset{jI}{}\sigma (\eta _i,\eta _j\eta _i)e_j\pi _{ji},$$
(2.3)
for all $`iI`$.
Similarly, let $`V^{}`$ be the graded right $`H`$–module that is dual to $`V`$. Regarded as a graded vector space, $`V^{}`$ is equal to $`V^{\mathrm{gr}}`$, and the module structure is given by
$$x^{}h,x=x^{},hx,$$
for all $`hH`$, $`xV`$, and $`x^{}V^{}`$. (Of course, this time $`,`$ denotes the dual pairing of $`V^{\mathrm{gr}}`$ and $`V`$.) More explicitly, we have
$$e_i^{}h=\underset{jI}{}\pi _{ij}(h)e_j^{},$$
for all $`hH`$ and all $`iI`$. Then $`V^{}`$ has a canonical structure of a graded left $`H^{}`$–comodule. It is uniquely fixed by the requirement that it induces, in a canonical way, the original structure of a graded right $`H`$–module on $`V^{}`$. More precisely, let
$$\delta ^{}:V^{}H^{}V^{}$$
be the structure map of $`V^{}`$, regarded as a graded left $`H^{}`$–comodule, let $`zV^{}`$, and set
$$\delta ^{}(z)=\underset{^r}{}z_r^1z_r^0,$$
where $`z_r^1H^{}`$, and where $`z_r^0V^{}`$ is homogeneous of degree $`\zeta _r^0`$. Then, if $`hH`$ is homogeneous of degree $`\eta `$, we have
$$zh=\underset{^r}{}\sigma (\zeta _r^0,\eta )z_r^1,hz_r^0.$$
Explicitly, we find
$$\delta ^{}(e_i^{})=\underset{jI}{}\sigma (\eta _j,\eta _i\eta _j)\pi _{ij}e_j^{},$$
(2.4)
for all $`iI`$.
Let us next consider two cases where our assumptions are more restrictive. In the first case, we suppose that on $`V`$ there exists a non–degenerate invariant bilinear form $`b`$ which is homogeneous of degree zero (see Appendix A of Ref. for some comments on invariant bilinear forms). Define a linear map
$$f_r:VV^{\mathrm{gr}}$$
by
$$f_r(y),x=\sigma (\eta ,\xi )b(x,y),$$
where $`x`$ and $`y`$ are elements of $`V`$ which are homogeneous of degrees $`\xi `$ and $`\eta `$, respectively. Then the assumption that $`b`$ is invariant is equivalent to the requirement that
$$\pi (S(h))=f_r^1\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(\pi (h))^\mathrm{T}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}f_r\text{for all }hH,$$
where <sup>T</sup> denotes the $`\sigma `$–transpose of a linear map (see Eqn. (A.5) of Ref. ). In terms of matrix elements, this condition takes the form
$$\pi _{ij}(S(h))=\underset{k,\mathrm{}}{}(F_r^1)_{ik}(\pi (h)^\mathrm{T})_k\mathrm{}(F_r)_\mathrm{}j\text{for all }i,jI\text{ and all }hH,$$
(2.5)
where $`F_r`$ is the matrix of $`f_r`$:
$$f_r(e_j)=\underset{iI}{}(F_r)_{ij}e_i^{}\text{for all }jI.$$
To understand the meaning of Eqn. (2.5) we note that
$$\pi _{ij}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}S=S_{}(\pi _{ij}),$$
(2.6)
moreover, we recall that the matrix elements of $`\pi (h)^\mathrm{T}`$ are given by
$$(\pi (h)^\mathrm{T})_k\mathrm{}=\sigma (\eta _{\mathrm{}},\eta _{\mathrm{}}\eta _k)\pi _\mathrm{}k(h).$$
(2.7)
Thus Eqn. (2.5) expresses $`S_{}(\pi _{ij})`$ as a linear combination of the matrix elements $`\pi _\mathrm{}k`$.
For convenience, we also give a formula for the matrix $`F_r`$. Setting
$$C_{ij}=b(e_i,e_j)$$
and
$$\sigma _{ij}=\sigma (\eta _i,\eta _j),\sigma _i=\sigma (\eta _i,\eta _i),$$
for all $`i,jI`$, we obtain
$$(F_r)_{ij}=\sigma _{ji}C_{ij}=\sigma _iC_{ij}=C_{ij}\sigma _j\text{for all }i,jI.$$
(2.8)
In the second case, we assume in addition that $`H`$ is quasitriangular, i.e., that we are given a universal $`R`$–matrix $``$ for $`H`$. By definition, $``$ is an invertible element of $`H\overline{}H`$ and is homogeneous of degree zero. Moreover, it satisfies the following relations:
$$\mathrm{\Delta }(h)=(P\mathrm{\Delta }(h))\text{for all }hH,$$
(2.9)
$$(\mathrm{\Delta }\text{id}_H)()=_{13}_{23},(\text{id}_H\mathrm{\Delta })()=_{13}_{12}.$$
(2.10)
Our notation is standard: $`P:H\overline{}HH\overline{}H`$ is the twist map (in the graded sense), and the elements $`_{12}`$, $`_{13}`$, and $`_{23}`$ of $`H\overline{}H\overline{}H`$ are defined by
$$_{12}=1,_{23}=1$$
$$_{13}=(P\text{id}_H)(_{23})=(\text{id}_HP)(_{12}).$$
Let us define the $`R`$–matrix $`R`$ (more precisely, the $`R`$–operator $`R`$) acting in $`VV`$ by
$$R=(\pi \pi )().$$
Using the notation
$$=\underset{^r}{}R_r^1R_r^2,$$
(2.11)
we have
$$R=\underset{ijk\mathrm{}I}{}R_{ij,k\mathrm{}}E_{ik}\overline{}E_j\mathrm{},$$
with
$$R_{ij,k\mathrm{}}=\underset{^r}{}\pi _{ik}(R_r^1)\pi _j\mathrm{}(R_r^2).$$
Acting with $`\pi \pi `$ on both sides of Eqn. (2.9), we obtain
$$\underset{a,bI}{}\sigma (\eta _j\eta _{\mathrm{}},\eta _a\eta _k)R_{ij,ab}\pi _{ak}\pi _b\mathrm{}=\underset{a,bI}{}\sigma (\eta _j\eta _b,\eta _a\eta _k)\pi _{jb}\pi _{ia}R_{ab,k\mathrm{}},$$
(2.12)
for all $`i,j,k,\mathrm{}I`$. These equations generalize the famous RTT– relations. By using the following technical device, they can be written in the familiar form.
Consider the algebra
$$\text{Lgr}(V)\overline{}\text{Lgr}(V)\overline{}H^{}\text{Lgr}(VV)\overline{}H^{}$$
and define the following elements therein:
$$\pi _1=\underset{i,jI}{}E_{ij}\text{id}_V\pi _{ij},\pi _2=\underset{i,jI}{}\text{id}_VE_{ij}\pi _{ij}.$$
(2.13)
Then the Eqns. (2.12) take the form
$$(R\epsilon )\pi _1\pi _2=\pi _2\pi _1(R\epsilon )$$
(2.14)
(recall that $`\epsilon `$ is the unit element of $`H^{}`$).
Remark 2.1. Some comments on the notation are in order. If $`A`$ and $`B`$ are two graded associative algebras, we denote their graded tensor product by $`A\overline{}B`$ and the decomposable tensors in $`A\overline{}B`$ by $`ab`$ (where $`aA`$ and $`bB`$, see the appendix of Ref. ). Let us next recall that there exists a canonical isomorphism of graded algebras
$$\text{Lgr}(V)\overline{}\text{Lgr}(V)\text{Lgr}(VV);$$
for all $`g,g^{}\text{Lgr}(V)`$, it maps the tensor $`gg^{}`$ onto the graded tensor product $`g\overline{}g^{}`$ of the linear maps $`g`$ and $`g^{}`$ (see the Eqns. (3.4) and (3.5) of Ref. ). Originally, $`R`$ is an element of $`\text{Lgr}(VV)`$. However, occasionally (as above) it is useful to regard it also as an element of $`\text{Lgr}(V)\overline{}\text{Lgr}(V)`$. As such, it will be written in the form $`R=_{ijk\mathrm{}I}R_{ij,k\mathrm{}}E_{ik}E_j\mathrm{}`$. The reader should be careful not to confound $`E_{ik}E_j\mathrm{}`$ with the usual (non–graded) tensor product of the linear maps $`E_{ik}`$ and $`E_j\mathrm{}`$ (which has been used in Ref. , but will never be used in the present work).
Now let $`A`$ be a graded sub–$`\sigma `$–bialgebra of $`H^{}`$. Then the universal $`R`$–matrix yields a bilinear form $`\varrho `$ on $`A`$ that is defined by
$$\varrho (a,b)=ab,\text{for all }a,bA.$$
Explicitly, we have
$$\varrho (a,b)=\underset{^r}{}\sigma (\rho _r^1,\rho _r^2)a,R_r^1b,R_r^2\text{for all }a,bA.$$
(We use the notation introduced in Eqn. (2.11) and assume that $`R_r^1`$ and $`R_r^2`$ are homogeneous of degrees $`\rho _r^1`$ and $`\rho _r^2`$, respectively.) In particular, if the linear forms $`\pi _{ij}`$ belong to $`A`$, we have
$$\varrho (\pi _{ik},\pi _j\mathrm{})=\sigma (\eta _i\eta _k,\eta _j\eta _{\mathrm{}})R_{ij,k\mathrm{}}\text{for all }i,j,k,\mathrm{}I.$$
(2.15)
Moreover, the defining properties of a universal $`R`$–matrix imply the following properties of $`\varrho `$. First of all, the bilinear form $`\varrho `$ is homogeneous of degree zero. Secondly, $`\varrho `$ is convolution invertible.
To understand this statement, we recall that $`AA`$ has a canonical structure of a graded coalgebra. Correspondingly, its graded dual $`(AA)^{\mathrm{gr}}`$ is a graded algebra (whose multiplication is called convolution). Since homogeneous bilinear forms on $`A`$ can be canonically identified with homogeneous linear forms on $`AA`$, it makes sense to require that $`\varrho `$ is invertible in $`(AA)^{\mathrm{gr}}`$, and this is the meaning of the statement above. More explicitly, this is to say that there exists a (necessarily unique) bilinear form $`\varrho ^{}`$ on $`A`$, which is homogeneous of degree zero and satisfies
$$\underset{^{rs}}{}\sigma (\alpha _r^2,\beta _s^1)\varrho (a_r^1,b_s^1)\varrho ^{}(a_r^2,b_s^2)=\underset{^{rs}}{}\sigma (\alpha _r^2,\beta _s^1)\varrho ^{}(a_r^1,b_s^1)\varrho (a_r^2,b_s^2)=\epsilon (a)\epsilon (b)$$
(2.16)
for all $`a,bA`$. In these equations, we have set
$$\mathrm{\Delta }(a)=\underset{^r}{}a_r^1a_r^2,\mathrm{\Delta }(b)=\underset{^s}{}b_s^1b_s^2,$$
and the elements $`a_r^1,a_r^2,b_s^1,b_s^2A`$ are supposed to be homogeneous of the degrees $`\alpha _r^1,\alpha _r^2,\beta _s^1,\beta _s^2`$, respectively.
Using the same notation, the relation (2.9) implies that
$$\underset{^{rs}}{}\sigma (\alpha _r^2,\beta _s^1)\varrho (a_r^1,b_s^1)a_r^2b_s^2=\underset{^{rs}}{}\sigma (\alpha _r^1+\alpha _r^2,\beta _s^1)b_s^1a_r^1\varrho (a_r^2,b_s^2)$$
(2.17)
for all $`a,bA`$, and the relations (2.10) imply that
$$\varrho (aa^{},b)=\underset{^s}{}\sigma (\alpha ^{},\beta _s^1)\varrho (a,b_s^1)\varrho (a^{},b_s^2)$$
(2.18)
$$\varrho (a,bb^{})=\underset{^r}{}\varrho (a_r^1,b^{})\varrho (a_r^2,b)$$
(2.19)
for all $`a,a^{},b,b^{}A`$, where $`a^{}`$ is supposed to be homogeneous of degree $`\alpha ^{}`$.
A bilinear form $`\varrho `$ on a $`\sigma `$–bialgebra $`A`$ that has the properties given above is called a universal $`r`$–form for $`A`$, and a $`\sigma `$–bialgebra (or $`\sigma `$–Hopf algebra) endowed with a universal $`r`$–form is said to be coquasitriangular.
If $`\varrho `$ is a universal $`r`$–form for a bialgebra $`A`$, then so is $`\varrho ^{}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}P`$, where $`P:A\overline{}AA\overline{}A`$ denotes the twist. The conditions above imply that
$$\varrho (a,1)=\varrho (1,a)=\epsilon (a)\text{for all }aA.$$
(2.20)
(To prove this one has to use the fact that $`\varrho `$ is invertible.) Moreover, if $`A`$ is a $`\sigma `$–Hopf algebra, it follows that
$$\varrho (S(a),b)=\varrho ^{}(a,b),\varrho ^{}(a,S(b))=\varrho (a,b),$$
(2.21)
for all $`a,bA`$.
Let $`A^{\mathrm{cop}}`$ (resp. $`A^{\mathrm{aop}}`$) be the $`\sigma `$–bialgebra which, regarded as a graded algebra (resp. as a graded coalgebra) is equal to $`A`$, but whose coalgebra (resp. algebra) structure is opposite (in the graded sense) to that of $`A`$. Then the Eqns. (2.18), (2.19), and (2.20) show that $`\varrho `$ is a dual pairing (possibly degenerate) of the $`\sigma `$–bialgebras $`A^{\mathrm{cop}}`$ and $`A`$, and also of the $`\sigma `$–bialgebras $`A`$ and $`A^{\mathrm{aop}}`$.
## 3 The $`\sigma `$–bialgebras $`A(R)`$
We keep the notation introduced at the beginning of Section 2 and start by recalling a simple fact (see Ref. for the non–graded case). Let $`C`$ be a graded coalgebra. Equivalently, this means that $`C`$ is a graded vector space and a coalgebra, and that the structure maps of the latter are homogeneous of degree zero. Consider the tensor algebra $`T(C)`$ of the graded vector space $`C`$. It is well–known that $`T(C)`$ is a $`\times \mathrm{\Gamma }`$–graded algebra. Then there exists a unique coalgebra structure on $`T(C)`$ that converts $`T(C)`$ into a $`\sigma `$–bialgebra, and such that the canonical injection $`\iota :CT(C)`$ is a homomorphism of graded coalgebras. The $`\sigma `$–bialgebra $`T(C)`$ and the injection $`\iota `$ have the following universal property:
If $`B`$ is any $`\sigma `$–bialgebra, and if $`\phi :CB`$ is a homomorphism of graded coalgebras, there exists a unique $`\sigma `$–bialgebra homomorphism $`\overline{\phi }:T(C)B`$ such that $`\phi =\overline{\phi }\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}\iota `$.
We apply the foregoing construction to the following special case. Let the index set $`I`$ and the family of degrees $`\eta =(\eta _i)_{iI}`$ be as at the beginning of Section 2. We consider the following graded coalgebra $`C(\eta )`$ : It has a basis $`(X_{ij})_{i,jI}`$, such that $`X_{ij}`$ is homogeneous of degree $`\eta _j\eta _i`$, and such that the coproduct $`\mathrm{\Delta }`$ and the counit $`\epsilon `$ are given by
$$\mathrm{\Delta }(X_{ij})=\underset{kI}{}\sigma (\eta _k\eta _i,\eta _k\eta _j)X_{ik}X_{kj}$$
$$\epsilon (X_{ij})=\delta _{ij},$$
for all $`i,jI`$. Obviously, according to Eqn. (2.1), the graded coalgebra $`C(\eta )`$ is (isomorphic to) the dual of the graded algebra $`\text{Lgr}(V)`$. We stress that $`T(C(\eta ))`$ is a free algebra with free generators $`X_{ij}`$; $`i,jI`$.
Now let
$$R:VVVV$$
be a linear map which is homogeneous of degree zero. As before, we use the notation
$$R=\underset{ijk\mathrm{}I}{}R_{ij,k\mathrm{}}E_{ik}\overline{}E_j\mathrm{}.$$
(3.1)
The homogeneity condition is equivalent to the requirement that
$$R_{ij,k\mathrm{}}=0\text{if }\eta _i+\eta _j\eta _k+\eta _{\mathrm{}}.$$
Consider the graded algebra
$$\text{Lgr}(V)\overline{}\text{Lgr}(V)\overline{}T(C(\eta ))\text{Lgr}(VV)\overline{}T(C(\eta ))$$
and introduce the following elements therein:
$$X_1=\underset{i,jI}{}E_{ij}\text{id}_VX_{ij},X_2=\underset{i,jI}{}\text{id}_VE_{ij}X_{ij}$$
(3.2)
(see the Eqns. (2.13)). Define the elements $`X_{ij,k\mathrm{}}T(C(\eta ))`$ ; $`i,j,k,\mathrm{}I`$, through the following equation:
$$(R1)X_1X_2X_2X_1(R1)=\underset{ijk\mathrm{}I}{}E_{ik}E_j\mathrm{}X_{ij,k\mathrm{}}$$
(see Eqn. (2.14)). According to the definition, $`X_{ij,k\mathrm{}}`$ is homogeneous of degree $`\eta _k\eta _i+\eta _{\mathrm{}}\eta _j`$ with respect to the $`\mathrm{\Gamma }`$–gradation and homogeneous of degree 2 with respect to the $``$–gradation. Explicitly, we obtain
$$X_{ij,k\mathrm{}}=\underset{a,bI}{}\sigma (\eta _j\eta _{\mathrm{}},\eta _a\eta _k)R_{ij,ab}X_{ak}X_b\mathrm{}\underset{a,bI}{}\sigma (\eta _j\eta _b,\eta _a\eta _k)X_{jb}X_{ia}R_{ab,k\mathrm{}},$$
for all $`i,j,k,\mathrm{}I`$.
Let $`J(R)`$ be the (two–sided) ideal of $`T(C(\eta ))`$ that is generated by the elements $`X_{ij,k\mathrm{}}`$ ; $`i,j,k,\mathrm{}I`$. It is not difficult to see that $`J(R)`$ is, in fact, a $`\mathrm{\Gamma }`$–graded biideal and a $`\times \mathrm{\Gamma }`$–graded ideal of $`T(C(\eta ))`$. Consequently, the quotient
$$A(R)=T(C(\eta ))/J(R)$$
inherits from $`T(C(\eta ))`$ the structure of a $`\sigma `$–bialgebra and also of a $`\times \mathrm{\Gamma }`$–graded algebra. (Our notation is somewhat incomplete since, obviously, $`A(R)`$ also depends on $`\sigma `$ and on $`\eta =(\eta _i)_{iI}`$.) The canonical image of $`X_{ij}`$ in $`A(R)`$ will be denoted by $`t_{ij}`$, for all $`i,jI`$.
Consider the graded algebra
$$\text{Lgr}(V)\overline{}\text{Lgr}(V)\overline{}A(R)\text{Lgr}(VV)\overline{}A(R)$$
and introduce the following elements therein:
$$T_1=\underset{i,jI}{}E_{ij}\text{id}_Vt_{ij},T_2=\underset{i,jI}{}\text{id}_VE_{ij}t_{ij}.$$
(3.3)
Then $`A(R)`$ is the universal graded algebra, generated by elements $`t_{ij}`$; $`i,jI`$, which are homogeneous of degree $`\eta _j\eta _i`$ and satisfy the following RTT– relation
$$(R1)T_1T_2=T_2T_1(R1).$$
(3.4)
Explicitly, this relation is equivalent to
$$\underset{a,bI}{}\sigma (\eta _j\eta _{\mathrm{}},\eta _a\eta _k)R_{ij,ab}t_{ak}t_b\mathrm{}=\underset{a,bI}{}\sigma (\eta _j\eta _b,\eta _a\eta _k)t_{jb}t_{ia}R_{ab,k\mathrm{}},$$
(3.5)
for all $`i,j,k,\mathrm{}I`$. Moreover, the coproduct $`\mathrm{\Delta }`$ and the counit $`\epsilon `$ are uniquely fixed by the equations
$$\mathrm{\Delta }(t_{ij})=\underset{kI}{}\sigma (\eta _k\eta _i,\eta _k\eta _j)t_{ik}t_{kj}$$
(3.6)
$$\epsilon (t_{ij})=\delta _{ij},$$
(3.7)
for all $`i,jI`$.
It is well–known that there exists a useful equivalent formulation of these relations. Let $`P:VVVV`$ be the twist. Explicitly, we obtain
$$P=\underset{i,jI}{}\sigma (\eta _i,\eta _i)E_{ji}\overline{}E_{ij}$$
(the factor $`\sigma (\eta _i,\eta _i)`$ is not a misprint). If $`g`$ and $`g^{}`$ are two elements of $`\text{Lgr}(V)`$ which are homogeneous of degrees $`\gamma `$ and $`\gamma ^{}`$, respectively, we have
$$P\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(g\overline{}g^{})=\sigma (\gamma ,\gamma ^{})(g^{}\overline{}g)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}P.$$
Define
$$\widehat{R}=PR.$$
Then the relations (3.4) can also be written in the form
$$(\widehat{R}1)T_1T_2=T_1T_2(\widehat{R}1).$$
(3.8)
Up to now, $`R`$ could be an arbitrary element of $`\text{Lgr}(VV)`$ that is homogeneous of degree zero. This indicates that we haven’t yet proved anything really substantial about the $`\sigma `$–bialgebras $`A(R)`$. Presumably, the most important property of the $`A(R)`$’s is that, under certain natural additional assumptions, they are coquasitriangular. More precisely, we have the following theorem.
Theorem 1. We keep the notation used above. Let $`R`$ be an element of $`\text{Lgr}(VV)`$ that is homogeneous of degree zero. Suppose in addition that $`R`$ is invertible and satisfies the graded Yang–Baxter equation. Then there exists a unique universal $`r`$–form $`\varrho `$ on $`A(R)`$ such that
$$\varrho (t_{ik},t_j\mathrm{})=\sigma (\eta _i\eta _k,\eta _j\eta _{\mathrm{}})R_{ij,k\mathrm{}}\text{for all }i,j,k,\mathrm{}I.$$
(3.9)
(We are using the notation introduced in Eqn. (3.1).)
Sketch of the proof
We start from the following fact which, in the non-graded case, has been mentioned in Ref. . Consider two graded coalgebras $`C`$ and $`D`$ and the corresponding $`\sigma `$–bialgebras $`T(C)`$ and $`T(D)`$ introduced at the beginning of this section. If $`\psi _1:D\times C𝕂`$ is any bilinear form which is homogeneous of degree zero, there exists a unique dual pairing $`\psi `$ of the $`\sigma `$–bialgebras $`T(D)`$ and $`T(C)`$ that extends $`\psi _1`$.
Applying this result to $`C=C(\eta )`$ and $`D=C(\eta )^{\mathrm{op}}`$, we obtain a unique dual pairing $`\psi `$ of the $`\sigma `$–bialgebras $`T(C(\eta ))^{\mathrm{cop}}=T(C(\eta )^{\mathrm{op}})`$ and $`T(C(\eta ))`$ such that
$$\psi (X_{ik},X_j\mathrm{})=\sigma (\eta _i\eta _k,\eta _j\eta _{\mathrm{}})R_{ij,k\mathrm{}}\text{for all }i,j,k,\mathrm{}I.$$
Let us next investigate under which conditions this form induces a bilinear form on $`A(R)`$. For this to be the case it is necessary that
$$\psi (X_{ij,k\mathrm{}},X_{ab})=0\text{and}\psi (X_{ab},X_{ij,k\mathrm{}})=0$$
for all $`i,j,k,\mathrm{},a,bI`$. Each of these conditions is satisfied if and only if the graded Yang–Baxter equation holds for $`R`$. Conversely, if this is the case, the bilinear form $`\psi `$ indeed induces a bilinear form $`\varrho `$ on $`A(R)`$. Obviously, $`\varrho `$ is a dual pairing of the $`\sigma `$–bialgebras $`A(R)^{\mathrm{cop}}`$ and $`A(R)`$.
Our next task is to show that the generalized commutation relations (2.17) are satisfied. The basic first step of the proof consists in checking that these relations hold for $`a=t_{ik}`$ and $`b=t_j\mathrm{}`$, for all $`i,j,k,\mathrm{}I`$. But up to unimportant overall factors, these are just the defining relations for $`A(R)`$.
Finally, we have to prove that $`\varrho `$ is convolution invertible. If $`\varrho ^{}`$ is an inverse of $`\varrho `$, the Eqns. (2.16), applied to $`a=t_{ik}`$ and $`b=t_j\mathrm{}`$, show that $`R`$ is invertible: In fact, the inverse of $`R`$ is given in terms of $`\varrho ^{}`$ exactly as $`R`$ is given in terms of $`\varrho `$ (see Eqn. (3.9)). Accordingly, we now assume that $`R`$ is invertible.
It is well–known (and easy to see) that if $`R`$ satisfies the graded Yang–Baxter equation, then so does $`\stackrel{~}{R}=PR^1P`$. Moreover, $`\stackrel{~}{R}`$ and $`R`$ define the same $`\sigma `$–bialgebra: We have
$$A(\stackrel{~}{R})=A(R).$$
(3.10)
Consequently, we can apply the foregoing results with $`R`$ replaced by $`\stackrel{~}{R}`$, which gives another dual pairing $`\stackrel{~}{\rho }`$ of the $`\sigma `$–bialgebras $`A(R)^{\mathrm{cop}}`$ and $`A(R)`$. According to the preceding discussion, we expect that $`\varrho ^{}=\stackrel{~}{\rho }\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}P`$ is the inverse of $`\varrho `$, and this is indeed the case.
## 4 Graded comodules and comodule algebras <br>over $`A(R)`$
We keep the notation introduced at the beginning of Section 2 and in Section 3. In the present section we are going to discuss certain natural graded $`A(R)`$–comodules and $`A(R)`$–comodule algebras. Actually, we make a slight generalization. The point is the following. At last, we are not so much interested in the $`\sigma `$–bialgebras $`A(R)`$ themselves, but rather in certain $`\sigma `$–Hopf algebra envelopes therefrom. Typically, the latter are obtained from $`A(R)`$ by adjoining the inverses of certain elements of $`A(R)`$, or by requiring additional relations (i.e., by going to a quotient). In any case, such an envelope will be a $`\sigma `$–bialgebra $`A`$, and there will be a natural homomorphism of $`\sigma `$–bialgebras $`A(R)A`$. For convenience, the natural image of $`t_{ij}`$ in $`A`$ will also be denoted by $`t_{ij}`$. This will be the setting considered in the present section. Stated differently, $`A`$ will be a $`\sigma `$–bialgebra containing certain elements $`t_{ij}`$; $`i,jI`$, such that $`t_{ij}`$ is homogeneous of degree $`\eta _j\eta _i`$, for all $`i,jI`$, such that the relation (3.4) (or its equivalents (3.5), (3.8)) is satisfied, and such that the Eqns. (3.6) and (3.7) hold true.
To begin with, we note that $`V`$ is a graded right $`A`$–comodule in a natural way: The structure map
$$\delta _1:VVA$$
is defined by
$$\delta _1(e_i)=\underset{jI}{}\sigma (\eta _i,\eta _j\eta _i)e_jt_{ji},$$
(4.1)
for all $`iI`$ (see Eqn. (2.3)).
Similarly, $`V^{\mathrm{gr}}`$ is a graded left $`A`$–comodule: The structure map
$$\delta _1^{}:V^{\mathrm{gr}}AV^{\mathrm{gr}}$$
is defined by
$$\delta _1^{}(e_i^{})=\underset{jI}{}\sigma (\eta _j,\eta _i\eta _j)t_{ij}e_j^{},$$
(4.2)
for all $`iI`$. We note that the graded comodule $`V^{\mathrm{gr}}`$ is dual to the graded comodule $`V`$ in the sense that
$$(\text{id}_A)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(\delta _1^{}\text{id}_V)=(\text{id}_A)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(\text{id}_{V^{\mathrm{gr}}}\delta _1).$$
(4.3)
(These are maps of $`V^{\mathrm{gr}}V`$ into $`A`$, and $`:V^{\mathrm{gr}}V𝕂`$ is given by the dual pairing.) Obviously, if one of the structure maps $`\delta _1`$ and $`\delta _1^{}`$ is known, the other is uniquely fixed by this equation.
Remark 4.1. At this point we insert a remark that could have been made earlier. Looking at the Eqns. (4.1) and (4.2), we are tempted to introduce the elements
$$\stackrel{~}{t}_{ij}=\sigma (\eta _j,\eta _i\eta _j)t_{ij}\text{for all }i,jI.$$
(4.4)
Then these equations take the familiar form
$$\delta _1(e_i)=\underset{jI}{}e_j\stackrel{~}{t}_{ji},\delta _1^{}(e_i^{})=\underset{jI}{}\stackrel{~}{t}_{ij}e_j^{}$$
for all $`iI`$, which implies that
$$\mathrm{\Delta }(\stackrel{~}{t}_{ij})=\underset{kI}{}\stackrel{~}{t}_{ik}\stackrel{~}{t}_{kj}$$
for all $`i,jI`$. (Needless to say, this equation can easily be checked directly.) Summarizing, the elements $`\stackrel{~}{t}_{ij}`$ seem to have properties which are closer to what we are accustomed to from the non–graded case. However, unfortunately I do not know how to write the commutation relations for the $`\stackrel{~}{t}_{ij}`$ in a factorized form like Eqn. (3.4). My conclusion is that, depending on the problem at hand, both types of generators may have advantages. (Actually, in the proof of Theorem 1 I have mostly used the generators $`\stackrel{~}{t}_{ij}`$, whereas for the proof of Eqn. (3.10) the factorization in Eqn. (3.4) is of great help.)
Once the graded $`A`$–comodule structures of $`V`$ and $`V^{\mathrm{gr}}`$ are defined, they can be used to convert $`T(V)`$ and $`T(V^{\mathrm{gr}})`$ into graded $`A`$–comodule algebras. Thus, there exists a unique homomorphism of graded algebras
$$\delta :T(V)T(V)\overline{}A$$
such that Eqn. (4.1) (with $`\delta _1`$ replaced by $`\delta `$) is satisfied, and it is easy to see that this map converts $`T(V)`$ into a graded right $`A`$–comodule, hence into a graded right $`A`$–comodule algebra. Note that, for every integer $`n0`$, the subspace $`T_n(V)=V^n`$ is a graded subcomodule of $`T(V)`$, in fact, this is just the $`n`$th tensorial power of the graded comodule $`V`$. Its structure map is the map
$$\delta _n:T_n(V)T_n(V)A$$
(4.5)
induced by $`\delta `$. (The map $`\delta _1`$ is the map introduced in Eqn. (4.1).)
Similarly, there exists a unique homomorphism of graded algebras
$$\delta ^{}:T(V^{\mathrm{gr}})A\overline{}T(V^{\mathrm{gr}})$$
such that Eqn. (4.2) (with $`\delta _1^{}`$ replaced by $`\delta ^{}`$) is satisfied, and it is easy to see that this map converts $`T(V^{\mathrm{gr}})`$ into a graded left $`A`$–comodule algebra. For every integer $`n0`$, the subspace $`T_n(V^{\mathrm{gr}})=(V^{\mathrm{gr}})^n`$ is a graded subcomodule of $`T(V^{\mathrm{gr}})`$, in fact, this is just the $`n`$th tensorial power of the graded comodule $`V^{\mathrm{gr}}`$. Its structure map is the map
$$\delta _n^{}:T_n(V^{\mathrm{gr}})AT_n(V^{\mathrm{gr}})$$
(4.6)
induced by $`\delta ^{}`$. (The map $`\delta _1^{}`$ is the map introduced in Eqn. (4.2).)
The graded comodule $`T_n(V^{\mathrm{gr}})`$ is dual to the graded comodule $`T_n(V)`$. Indeed, since the case $`n=0`$ is trivial, we may assume that $`n1`$. Let us define a canonical pairing
$$,:T_n(V^{\mathrm{gr}})\times T_n(V)𝕂$$
by
$$x_1^{}\mathrm{}x_n^{},x_1\mathrm{}x_n=\underset{1a<bn}{}\sigma (\xi _b^{},\xi _a)\underset{1cn}{}x_c^{},x_c$$
for all $`x_1,\mathrm{},x_nV`$; $`x_1^{},\mathrm{},x_n^{}V^{\mathrm{gr}}`$, which are homogeneous of the degrees $`\xi _1,\mathrm{},\xi _n,\xi _1^{},\mathrm{},\xi _n^{}`$, respectively. Then $`\delta _n`$ and $`\delta _n^{}`$ are related by an equation analogous to Eqn. (4.3).
The algebra $`T(V)`$ is free, with free generators $`e_i`$ ; $`iI`$. In order to obtain more interesting examples of graded $`A`$–comodule algebras, we have to construct suitable quotients of $`T(V)`$. In principle, this is easily done (see Ref. ). Let $`W`$ be any graded $`A`$–subcomodule of $`T(V)`$, let $`J(W)`$ be the ideal of $`T(V)`$ that is generated by $`W`$, and let
$$B_r(W)=T(V)/J(W)$$
be the corresponding quotient. Of course, $`B_r(W)`$ is a graded algebra. By assumption, we have
$$\delta (W)WA.$$
Obviously, this implies that
$$\delta (J(W))J(W)A.$$
Consequently, the map $`\delta `$ induces a map
$$\delta :B_r(W)B_r(W)\overline{}A$$
(note the abuse of notation), and it is clear that this map converts $`B_r(W)`$ into a graded right $`A`$–comodule algebra. Let us denote the canonical image of $`e_i`$ in $`B_r(W)`$ by $`x_i`$, for all $`iI`$. Then $`B_r(W)`$ is the universal graded algebra, generated by elements $`x_i`$ ; $`iI`$, which are homogeneous of degree $`\eta _i`$ and satisfy the relations “given by $`W`$ ”. Moreover, the comodule structure of $`B_r(W)`$ is fixed by
$$\delta (x_i)=\underset{jI}{}\sigma (\eta _i,\eta _j\eta _i)x_jt_{ji}\text{for all }iI.$$
(4.7)
Instead of writing the defining relations in general, let us give a simple example. Suppose that $`W`$ is a graded subcomodule of $`T_2(V)=VV`$. Let $`(a^s)_{1sr}`$ be a family of elements of $`W`$ that generates the vector space $`W`$, and let us write
$$a^s=\underset{i,jI}{}a_{ij}^se_ie_j\text{for }1sr.$$
Then the defining relations of $`B_r(W)`$ can be written in the form
$$\underset{i,jI}{}a_{ij}^sx_ix_j=0\text{for }1sr.$$
In practice, it might be difficult to prove that a graded subspace of $`T(V)`$ is a graded subcomodule. Indeed, all we know about the algebra structure of $`A`$ are the commutation relations (3.5). It is at this point where the techniques of Ref. are of great help.
Let $`S\text{Lgr}(VV)`$ be homogeneous of degree zero. As usual, we write
$$S=\underset{ijk\mathrm{}I}{}S_{ij,k\mathrm{}}E_{ik}\overline{}E_j\mathrm{}.$$
(4.8)
We recall that the homogeneity condition is equivalent to the requirement that
$$S_{ij,k\mathrm{}}=0\text{if }\eta _i+\eta _j\eta _k+\eta _{\mathrm{}}.$$
Let us also introduce the elements
$$T_1,T_2\text{Lgr}(V)\overline{}\text{Lgr}(V)\overline{}A\text{Lgr}(VV)\overline{}A$$
as in Eqn. (3.3). Then $`S`$ is an endomorphism of the graded $`A`$–comodule $`VV`$ if and only if
$$(S1)T_1T_2=T_1T_2(S1).$$
(4.9)
Explicitly, this relation is equivalent to
$$\underset{a,bI}{}\sigma (\eta _j\eta _{\mathrm{}},\eta _a\eta _k)S_{ij,ab}t_{ak}t_b\mathrm{}=\underset{a,bI}{}\sigma (\eta _j\eta _b,\eta _i\eta _k)t_{ia}t_{jb}S_{ab,k\mathrm{}},$$
(4.10)
for all $`i,j,k,\mathrm{}I`$.
A look at Eqn. (3.8) shows that $`S=\widehat{R}`$ satisfies Eqn. (4.9). Indeed, the algebra $`A(R)`$ is defined such that $`\widehat{R}`$ is an endomorphism of the graded $`A(R)`$–comodule $`VV`$. It follows that every polynomial in $`\widehat{R}`$ satisfies this equation as well. Since the image and the kernel of every such $`S`$ are graded $`A`$–subcomodules of $`VV`$ (and hence of $`T(V)`$), we have found some natural candidates for the subcomodule $`W`$.
We can now comment on how the present formulation generalizes the approach of Ref. . In that reference, the $`A(R)`$–comodule $`VV`$ is completely reducible (provided that $`q`$ is generic), and the subcomodules can be constructed as the images of suitable polynomials in $`\widehat{R}`$. As we are going to see, for the quantum supergroups $`\text{SPO}_q(2n|2n)`$ this is not true anymore, and we are forced to use the generalized method described above.
The graded left $`A`$–comodule algebras can be constucted in a completely analogous way, and hence it should not be necessary to give all the details. In particular, for every graded subcomodule $`W^{}`$ of $`T(V^{\mathrm{gr}})`$, we obtain a graded left $`A`$-comodule algebra $`B_{\mathrm{}}(W^{})`$. However, one point should be mentioned. In this second case, we need to know the endomorphisms of the graded $`A`$–comodule $`V^{\mathrm{gr}}V^{\mathrm{gr}}`$. Every element of $`\text{Lgr}(V^{\mathrm{gr}}V^{\mathrm{gr}})`$ can be written as the graded transpose $`S^\mathrm{T}`$ of some element $`S\text{Lgr}(VV)`$, and $`S^\mathrm{T}`$ is an endomorphism of the graded $`A`$–comodule $`V^{\mathrm{gr}}V^{\mathrm{gr}}`$ if and only if $`S`$ is an endomorphism of the graded $`A`$–comodule $`VV`$, i.e., if and only if the equivalent conditions (4.9) and (4.10) are satisfied. Using the notation (4.8) and assuming that $`S`$ is homogeneous of degree zero, we have
$$S^\mathrm{T}(e_i^{}e_j^{})=\underset{k,\mathrm{}I}{}\sigma (\eta _j,\eta _k\eta _i)S_{ij,k\mathrm{}}e_k^{}e_{\mathrm{}}^{}.$$
## 5 The quantum supergroup $`\text{SPO}_q(2n|2m)`$
In the present and in the subsequent section we are going to apply the general theory to the $`R`$–matrix that has been calculated in Ref. , namely, to the $`R`$–matrix of $`U_q(𝔰𝔭𝔬(2n|2m))`$ in the vector representation. Correspondingly, we shall use the notation introduced in that reference. First of all, this implies that in the following we choose
$$\mathrm{\Gamma }=_2=\{\overline{0},\overline{1}\}$$
and define
$$\sigma (\alpha ,\beta )=(1)^{\alpha \beta }\text{for all }\alpha ,\beta _2.$$
Moreover, we set
$$r=n+m,d=nm,$$
and the index set $`I`$ is equal to
$$I=\{r,r+1,\mathrm{},2,1,1,2,\mathrm{},r1,r\}.$$
The degrees $`\eta _i`$; $`iI`$, are given by
$$\eta _i=\{\begin{array}{cc}\overline{0}\hfill & \text{for }1|i|n\hfill \\ \overline{1}\hfill & \text{for }n+1|i|r.\hfill \end{array}$$
We recall that
$$\sigma _i=\sigma (\eta _i,\eta _i)\text{and}\sigma _{i,j}=\sigma (\eta _i,\eta _j)\text{for all }i,jI.$$
In the present setting, $`V`$ is the vector module of $`U_q(𝔰𝔭𝔬(2n|2m))`$. It carries a $`U_q(𝔰𝔭𝔬(2n|2m))`$–invariant bilinear form $`b^q`$, whose matrix $`C^q=(C_{i,j}^q)`$ with respect to the basis $`(e_i)_{iI}`$ has been given at the end of Section 4 of Ref. . Finally, we note that the $`R`$–matrix $`R`$ and the braid generator $`\widehat{R}`$ are given by
$$\begin{array}{ccc}\hfill R& =& \underset{i}{}q^{\sigma _i}E_{i,i}\overline{}E_{i,i}+\underset{i}{}q^{\sigma _i}E_{i,i}\overline{}E_{i,i}\hfill \\ & & +\underset{ij,j}{}E_{i,i}\overline{}E_{j,j}\hfill \\ & & +(qq^1)\underset{i<j}{}\sigma _iE_{j,i}\overline{}E_{i,j}\hfill \\ & & (qq^1)\underset{i<j}{}\sigma _i\sigma _j\sigma _{i,j}((C^q)^1)_{j,j}C_{i,i}^qE_{j,i}\overline{}E_{j,i}\hfill \end{array}$$
$$\begin{array}{ccc}\hfill \widehat{R}& =& \underset{i}{}\sigma _iq^{\sigma _i}E_{i,i}\overline{}E_{i,i}+\underset{i}{}\sigma _iq^{\sigma _i}E_{i,i}\overline{}E_{i,i}\hfill \\ & & +\underset{ij,j}{}\sigma _iE_{j,i}\overline{}E_{i,j}\hfill \\ & & +(qq^1)\underset{i<j}{}E_{i,i}\overline{}E_{j,j}\hfill \\ & & (qq^1)\underset{i<j}{}\sigma _i\sigma _{i,j}((C^q)^1)_{j,j}C_{i,i}^qE_{j,i}\overline{}E_{j,i},\hfill \end{array}$$
and that the operator $`K`$ in $`VV`$ is equal to
$$K=\underset{i,j,k,\mathrm{}}{}\sigma _{k,j}\sigma _{k,\mathrm{}}((C^q)^1)_{i,j}C_{k,\mathrm{}}^qE_{i,k}\overline{}E_{j,\mathrm{}}.$$
Remark 5.1. In Ref. we have assumed that the base field $`𝕂`$ is equal to $``$ (since this assumption has been made by Yamane in Ref. ). Obviously, the $`R`$–matrix obtained in Ref. makes sense under the present more general assumptions, and the properties derived there are still true. Actually, also the definition of the quantum superalgebra $`U_q(𝔰𝔭𝔬(2n|2m))`$, its properties, and the discussion of the vector module and of its tensorial square all remain valid. (Indeed, according to my understanding Yamane’s entire paper is correct with $``$ replaced by $`𝕂`$.) For convenience, we still assume that $`q`$ is not a root of unity, although this condition could be relaxed.
The $`\sigma `$–bialgebras $`A(R)`$ don’t have an antipode (except for the degenerate case where $`I`$ is empty): This follows immediately by means of the $``$–gradation of these algebras. Hence some amendments are necessary if we want to obtain a $`\sigma `$–Hopf algebra. In view of the comments made in Section 2, it is easy to guess how this can be achieved under the present assumptions.
Let $`T^{\mathrm{st}}`$ be the super–transpose of the matrix $`T`$, defined by
$$(T^{\mathrm{st}})_{k,\mathrm{}}=\sigma (\eta _{\mathrm{}},\eta _{\mathrm{}}\eta _k)t_{\mathrm{},k}\text{for all }k,\mathrm{}I,$$
(see Eqn. (2.7)), and let $`F_r`$ be the $`I\times I`$–matrix over $`𝕂`$ , defined by
$$(F_r)_{i,j}=\sigma _{j,i}C_{i,j}^q=\sigma _iC_{i,j}^q=C_{i,j}^q\sigma _j\text{for all }i,jI$$
(see Eqn. (2.8)). Then the Eqns. (2.2), (2.5), and (2.6) suggest to define the $`I\times I`$–matrix (with elements in $`A(R)`$)
$$T^{}=F_r^1T^{\mathrm{st}}F_r,$$
explicitly:
$$t_{i,j}^{}=\sigma _i\sigma _{i,j}((C^q)^1)_{i,i}t_{j,i}C_{j,j}^q,$$
and to investigate, for all $`i,jI`$, the expressions
$$\begin{array}{ccc}\hfill Q_{i,j}^{}& =& \underset{kI}{}\sigma (\eta _k\eta _i,\eta _k\eta _j)t_{i,k}^{}t_{k,j}\hfill \\ \hfill Q_{i,j}^{\prime \prime }& =& \underset{kI}{}\sigma (\eta _k\eta _i,\eta _k\eta _j)t_{i,k}t_{k,j}^{}.\hfill \end{array}$$
This can be carried out as in Ref. . The basic observation is that the operator $`K`$ can be written as a polynomial in $`\widehat{R}`$ (see Eqn. (7.3) of Ref. ). This implies that the relations (4.10) are satisfied for $`S=K`$, and it is easy to evaluate these relations.
If $`ij`$ and $`k\mathrm{}`$, both sides of these relations are trivially equal to zero, and no information can be obtained.
If $`i=j`$ and $`k\mathrm{}`$, it follows that
$$\underset{a,bI}{}\sigma _{a,\mathrm{}}\sigma _{b,a}C_{a,b}^qt_{a,k}t_{b,\mathrm{}}=\mathrm{\hspace{0.17em}0}\text{for }k\mathrm{}.$$
If $`ij`$ and $`k=\mathrm{}`$, we find that
$$\underset{a,bI}{}\sigma _{i,b}((C^q)^1)_{a,b}t_{i,a}t_{j,b}=\mathrm{\hspace{0.17em}0}\text{for }ij.$$
If $`i=j`$ and $`k=\mathrm{}`$, we set
$$\begin{array}{ccc}\hfill Q^{}& =& (C_{k,k}^q)^1\underset{a,bI}{}\sigma _{a,k}\sigma _{b,a}C_{a,b}^qt_{a,k}t_{b,k}\hfill \\ \hfill Q^{\prime \prime }& =& (((C^q)^1)_{i,i})^1\sigma _i\underset{a,bI}{}\sigma _{i,b}((C^q)^1)_{a,b}t_{i,a}t_{i,b}.\hfill \end{array}$$
Then the relation under investigation says that
$$Q^{}=Q^{\prime \prime }.$$
Obviously, $`Q^{}`$ does not depend on $`i,j`$, and $`Q^{\prime \prime }`$ does not depend on $`k,\mathrm{}`$. This implies that $`Q^{}`$ and $`Q^{\prime \prime }`$ do not depend on $`i,j,k,\mathrm{}`$, and the equations above can be summarized as follows:
$`{\displaystyle \underset{a,bI}{}}\sigma _{\mathrm{},b}\sigma _{b,a}C_{a,b}^qt_{a,k}t_{b,\mathrm{}}`$ $`=`$ $`C_{k,\mathrm{}}^qQ^{}`$ (5.1)
$`{\displaystyle \underset{a,bI}{}}\sigma _{i,b}((C^q)^1)_{a,b}t_{i,a}t_{j,b}`$ $`=`$ $`\sigma _{i,j}((C^q)^1)_{i,j}Q^{\prime \prime }.`$ (5.2)
Remark 5.2. Using the generators $`\stackrel{~}{t}_{i,j}`$ as defined in Eqn. (4.4), i.e.,
$$\stackrel{~}{t}_{i,j}=\sigma _j\sigma _{j,i}t_{i,j}\text{for all }i,jI,$$
and defining the super–transpose $`\stackrel{~}{T}^{\mathrm{st}}`$ of the $`I\times I`$–matrix $`\stackrel{~}{T}=(\stackrel{~}{t}_{k,\mathrm{}})`$ by
$$(\stackrel{~}{T}^{\mathrm{st}})_{k,\mathrm{}}=\sigma _{\mathrm{}}\sigma _{k,\mathrm{}}\stackrel{~}{t}_{\mathrm{},k}\text{for all }k,\mathrm{}I$$
(see Eqn. (2.7)), it is easy to see that the relations (5.1) and (5.2) can be written in the form
$$\stackrel{~}{T}^{\mathrm{st}}C^q\stackrel{~}{T}=Q^{}C^q,\stackrel{~}{T}(C^q)^1\stackrel{~}{T}^{\mathrm{st}}=Q^{\prime \prime }(C^q)^1.$$
Now it is easy to see that
$$Q_{i,j}^{}=\delta _{i,j}Q^{},Q_{i,j}^{\prime \prime }=\delta _{i,j}Q^{\prime \prime }.$$
(5.3)
In the following, we set
$$Q^{}=Q^{\prime \prime }=Q.$$
Obviously, the element $`Q`$ is homogeneous of degree zero. Using the Eqns. (5.1) and (5.2) once again, we can also show that $`Q`$ is group–like, i.e., that
$$\mathrm{\Delta }(Q)=QQ,\epsilon (Q)=1.$$
Finally, we prove that $`Q`$ lies in the center of $`A(R)`$. Consider the following elements of the graded tensor product $`\text{Lgr}(V)\overline{}A(R)`$ :
$$T_0=\underset{i,jI}{}E_{i,j}t_{i,j},T_0^{}=\underset{i,jI}{}E_{i,j}t_{i,j}^{}.$$
Then the Eqns. (5.3) can be written in the following form:
$$T_0^{}T_0=T_0T_0^{}=𝕀Q$$
(5.4)
(where we have set $`𝕀=\text{id}_V`$). These equations imply that
$$T_0(𝕀Q)=T_0T_0^{}T_0=(𝕀Q)T_0.$$
Since $`Q`$ is homogeneous of degree zero, this is equivalent to
$$t_{i,j}Q=Qt_{i,j}\text{for all }i,jI,$$
which proves our claim.
After these preparations, we can define the quantum supergroup $`\text{SPO}_q(2n|2m)`$, as follows. Let $`J(Q)`$ be the ideal of $`A(R)`$ that is generated by $`Q1`$. Since $`Q`$ is homogeneous of degree zero and group–like, $`J(Q)`$ is a graded biideal of $`A(R)`$. Consequently, the quotient
$$\text{SPO}_q(2n|2m)=A(R)/J(Q)$$
(5.5)
inherits from $`A(R)`$ the structure of a bi–superalgebra (i.e., a $`\sigma `$–bialgebra with respect to the commutation factor $`\sigma `$ of supersymmetry). Let
$$\omega :A(R)\text{SPO}_q(2n|2m)$$
be the canonical map. Until further notice, we are going to use the notation
$$\omega (t_{i,j})=\overline{t}_{i,j}\text{for all }i,jI.$$
Similarly, we introduce the elements $`\overline{t}_{i,j}^{}`$ and the $`I\times I`$–matrices $`\overline{T}`$ and $`\overline{T}^{}`$.
Let us now show that $`\text{SPO}_q(2n|2m)`$ is a Hopf superalgebra. According to the preceding results we expect that the antipode $`S`$ of $`\text{SPO}_q(2n|2m)`$ is given by
$$S(\overline{T})=\overline{T}^{}$$
(5.6)
(to be interpreted element–wise). Consequently, we have to show that there exists a homomorphism of graded algebras
$$S:\text{SPO}_q(2n|2m)\text{SPO}_q(2n|2m)^{\mathrm{aop}}$$
(5.7)
such that Eqn. (5.6) is satisfied.
To prove this, we work in the graded tensor product
$$\text{Lgr}(V)\overline{}\text{Lgr}(V)\overline{}\text{SPO}_q(2n|2m)\text{Lgr}(VV)\overline{}\text{SPO}_q(2n|2m).$$
In that algebra, we have (with our usual notation)
$$(R1)\overline{T}_1\overline{T}_2=\overline{T}_2\overline{T}_1(R1),$$
and also
$$\overline{T}_1^1=\overline{T}_1^{},\overline{T}_2^1=\overline{T}_2^{}$$
(see Eqn. (5.4)). This implies that
$$(R1)\overline{T}_2^{}\overline{T}_1^{}=\overline{T}_1^{}\overline{T}_2^{}(R1).$$
A moment’s thought then shows that this is the RTT– relation (3.4) for $`\overline{T}^{}`$ in $`\text{SPO}_q(2n|2m)^{\mathrm{aop}}`$.
This result implies that there exists a homomorphism of graded algebras
$$S_0:A(R)\text{SPO}_q(2n|2m)^{\mathrm{aop}}$$
such that
$$S_0(T)=\overline{T}^{}.$$
It is easy to check that
$$S_0(Q^{})=\omega (Q^{\prime \prime })=1,S_0(Q^{\prime \prime })=\omega (Q^{})=1.$$
Consequently, $`S_0`$ induces the homomorphism (5.7) we are looking for. It is now easy to see that $`S`$ is an antipode: The identities to be proved hold (by construction) when evaluated on the generators $`\overline{t}_{i,j}`$ of $`\text{SPO}_q(2n|2m)`$, and this implies that they hold on all of $`\text{SPO}_q(2n|2m)`$.
Summarizing, we have shown that $`\text{SPO}_q(2n|2m)`$ is a Hopf superalgebra. In the sequel, we shall simplify the notation and write $`t_{i,j}`$ instead of $`\overline{t}_{i,j}`$. Then the antipode is uniquely fixed by the equation
$$S(T)=T^{},$$
where
$$T^{}=F_r^1T^{\mathrm{st}}F_r,$$
or more explicitly
$$t_{i,j}^{}=\sigma _i\sigma _{i,j}((C^q)^1)_{i,i}t_{j,i}C_{j,j}^q\text{for all }i,jI.$$
We note that
$$S^2(t_{i,j})=d_id_jt_{i,j}\text{for all }i,jI,$$
where
$$d_i=\sigma _i((C^q)^1)_{i,i}C_{i,i}^q\text{for all }iI.$$
Obviously, we have
$$d_i=d_i^1\text{for all }iI.$$
In $`\text{SPO}_q(2n|2m)`$, the relations (5.1) and (5.2) take the form
$`{\displaystyle \underset{a,bI}{}}\sigma _{\mathrm{},b}\sigma _{b,a}C_{a,b}^qt_{a,k}t_{b,\mathrm{}}`$ $`=`$ $`C_{k,\mathrm{}}^q`$ (5.8)
$`{\displaystyle \underset{a,bI}{}}\sigma _{i,b}((C^q)^1)_{a,b}t_{i,a}t_{j,b}`$ $`=`$ $`\sigma _{i,j}((C^q)^1)_{i,j}`$ (5.9)
for all $`i,j,k,\mathrm{}I`$. These relations have a simple interpretation. Eqn. (5.8) is equivalent to the fact that the linear form $`\stackrel{~}{b}^q:VV𝕂`$ associated to the bilinear form $`b^q`$ is a homomorphism of graded right $`\text{SPO}_q(2n|2m)`$–comodules, and also to the fact that the element
$$a^{}=\underset{i,jI}{}\sigma _{j,i}C_{i,j}^qe_i^{}e_j^{}$$
(5.10)
is an invariant of the graded left $`\text{SPO}_q(2n|2m)`$–comodule $`V^{\mathrm{gr}}V^{\mathrm{gr}}`$, in the well–known sense that
$$\delta _2^{}(a^{})=1a^{}.$$
Similarly, Eqn. (5.9) is equivalent to the fact that the linear form
$$\stackrel{~}{b}^q:V^{\mathrm{gr}}V^{\mathrm{gr}}𝕂$$
defined by
$$\stackrel{~}{b}^q(e_i^{}e_j^{})=\sigma _{i,j}((C^q)^1)_{i,j}\text{for all }i,jI,$$
is a homomorphism of graded left $`\text{SPO}_q(2n|2m)`$–comodules, and also to the fact that the element
$$a=\underset{i,jI}{}((C^q)^1)_{i,j}e_ie_j$$
(5.11)
is an invariant of the graded right $`\text{SPO}_q(2n|2m)`$–comodule $`VV`$, in the sense that
$$\delta _2(a)=a1.$$
Of course, this is exactly what we should guess.
According to Section 2, we expect that the Hopf superalgebra $`\text{SPO}_q(2n|2m)`$ is coquasitriangular. More precisely, we expect that the universal $`r`$–form $`\varrho `$ on $`A(R)`$ (as specified in Theorem 1) induces a universal $`r`$–form on $`\text{SPO}_q(2n|2m)`$.
In order to show this we have to prove that $`\varrho `$ vanishes on $`J(Q)\times A(R)`$ and on $`A(R)\times J(Q)`$ (where $`J(Q)`$ is the ideal of $`A(R)`$ used in the definition of $`\text{SPO}_q(2n|2m)`$, see Eqn. (5.5)). The basic ingredients of the proof are the equations
$$\varrho (C_{k,\mathrm{}}^q(Q^{}1),t_{i,j})=0,\varrho (t_{i,j},C_{k,\mathrm{}}^q(Q^{}1))=0,$$
which hold for all $`i,j,k,\mathrm{}I`$. These are equivalent to the equations
$$R^{\mathrm{st}_1}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(f_r𝕀)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}R=f_r𝕀,R^{\mathrm{st}_2}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(𝕀f_{\mathrm{}})\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}R=𝕀f_{\mathrm{}},$$
respectively, which have been proved in Ref. (see Eqn. (7.5) of that reference). We recall that $`f_{\mathrm{}}`$ and $`f_r`$ are the linear maps of $`V`$ into $`V^{\mathrm{gr}}`$ given by
$$f_{\mathrm{}}(e_j)=\underset{iI}{}C_{j,i}^qe_i^{},f_r(e_j)=\underset{iI}{}\sigma _{j,i}C_{i,j}^qe_i^{}$$
for all $`jI`$.
According to the preceding remarks, the bilinear form $`\varrho `$ induces a bilinear form $`\overline{\varrho }`$ on $`\text{SPO}_q(2n|2m)`$. Obviously, this form satisfies the Eqns. (2.17) – (2.19). In order to prove that $`\overline{\varrho }`$ is convolution invertible, one might think we have to show that the inverse $`\varrho ^{}`$ of $`\varrho `$ also induces a bilinear form on $`\text{SPO}_q(2n|2m)`$. However, it is not necessary to prove this independently: Since $`\text{SPO}_q(2n|2m)`$ is a Hopf superalgebra, it is known (and easy to see) that the bilinear form $`\overline{\varrho }^{}`$ on $`\text{SPO}_q(2n|2m)`$, defined by
$$\overline{\varrho }^{}(b,b^{})=\overline{\varrho }(S(b),b^{})\text{for all }b,b^{}\text{SPO}_q(2n|2m),$$
is an inverse of $`\overline{\varrho }`$.
Simplifying the notation, we have proved the following theorem.
Theorem 2. The Hopf superalgebra $`\text{SPO}_q(2n|2m)`$ is coquasitriangular. More precisely, there exists a unique universal $`r`$–form $`\varrho `$ on $`\text{SPO}_q(2n|2m)`$ such that
$$\varrho (t_{ik},t_j\mathrm{})=\sigma (\eta _i\eta _k,\eta _j\eta _{\mathrm{}})R_{ij,k\mathrm{}}\text{for all }i,j,k,\mathrm{}I.$$
(5.12)
(We are using the notation introduced in Eqn. (3.1). Since $`R`$ is homogeneous of degree zero, the factor $`\sigma (\eta _i\eta _k,\eta _j\eta _{\mathrm{}})`$ can be replaced by $`\sigma _i\sigma _k`$, and also by $`\sigma _j\sigma _{\mathrm{}}`$.)
Finally (and once again according to Section 2) we expect that the Hopf superalgebras $`U_q(𝔰𝔭𝔬(2n|2m))`$ and $`\text{SPO}_q(2n|2m)`$ are dual to each other in some sense. More precisely, we expect that there exists a Hopf superalgebra pairing of $`\text{SPO}_q(2n|2m)`$ and $`U_q(𝔰𝔭𝔬(2n|2m))`$. This can be proved, as follows.
Consider the vector module $`V`$ of $`U_q(𝔰𝔭𝔬(2n|2m))`$ (see Ref. ), and let $`\pi `$ be the graded representation of $`U_q(𝔰𝔭𝔬(2n|2m))`$ afforded by it. Define the linear forms $`\pi _{i,j}U_q(𝔰𝔭𝔬(2n|2m))^{}`$ by
$$\pi _{i,j}(h)=e_i^{},\pi (h)e_j\text{for all }i,jI\text{ and all }hU_q(𝔰𝔭𝔬(2n|2m)).$$
We know that $`\widehat{R}`$ is an endomorphism of the graded $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`VV`$. Defining the elements $`\pi _1`$ and $`\pi _2`$ of the superalgebra
$$\text{Lgr}(V)\overline{}\text{Lgr}(V)\overline{}U_q(𝔰𝔭𝔬(2n|2m))^{}\text{Lgr}(VV)\overline{}U_q(𝔰𝔭𝔬(2n|2m))^{}$$
(5.13)
as in Eqn. (2.13), it is easy to see that this is equivalent to the fact that the equation
$$(\widehat{R}\epsilon )\pi _1\pi _2=\pi _1\pi _2(\widehat{R}\epsilon )$$
holds in the algebra (5.13), where $`\epsilon `$ is the counit of $`U_q(𝔰𝔭𝔬(2n|2m))`$, i.e., the unit element of $`U_q(𝔰𝔭𝔬(2n|2m))^{}`$.
By definition of $`A(R)`$, this implies that there exists a unique algebra homomorphism
$$\chi _0:A(R)U_q(𝔰𝔭𝔬(2n|2m))^{}$$
such that
$$\chi _0(t_{i,j})=\pi _{i,j}\text{for all }i,jI.$$
Obviously, $`\chi _0`$ is a bi–superalgebra homomorphism, and we have
$$\chi _0(C_{k,\mathrm{}}^qQ^{})=\underset{a,bI}{}\sigma _{\mathrm{},b}\sigma _{b,a}C_{a,b}^q\pi _{a,k}\pi _{b,\mathrm{}}\text{for all }k,\mathrm{}I.$$
Since the bilinear form $`b^q`$ on $`V`$ is $`U_q(𝔰𝔭𝔬(2n|2m))`$–invariant (equivalently, since the associated linear form $`\stackrel{~}{b}^q`$ on $`VV`$ is a homomorphism of graded $`U_q(𝔰𝔭𝔬(2n|2m))`$–modules), the right hand side is equal to $`C_{k,\mathrm{}}^q\epsilon `$. Consequently, $`\chi _0`$ induces a bi–superalgebra homomorphism
$$\chi :\text{SPO}_q(2n|2m)U_q(𝔰𝔭𝔬(2n|2m))^{}$$
such that
$$\chi (t_{i,j})=\pi _{i,j}\text{for all }i,jI.$$
Since $`\text{SPO}_q(2n|2m)`$ and $`U_q(𝔰𝔭𝔬(2n|2m))^{}`$ are Hopf superalgebras, $`\chi `$ is known to be a Hopf superalgebra homomorphism. This implies (indeed, is equivalent to the fact) that the bilinear form
$$\phi :\text{SPO}_q(2n|2m)\times U_q(𝔰𝔭𝔬(2n|2m))𝕂$$
defined by
$$\phi (a,h)=\chi (a),h\text{for all }a\text{SPO}_q(2n|2m)\text{ and all }hU_q(𝔰𝔭𝔬(2n|2m))$$
(5.14)
is a Hopf superalgebra pairing.
Summarizing, we have proved the following theorem.
Theorem 3. There exists a unique Hopf superalgebra pairing
$$\phi :\text{SPO}_q(2n|2m)\times U_q(𝔰𝔭𝔬(2n|2m))𝕂$$
such that
$$\phi (t_{i,j},h)=\pi _{i,j}(h)\text{for all }i,jI\text{ and all }hU_q(𝔰𝔭𝔬(2n|2m)).$$
Let $`A`$ be the image of $`\text{SPO}_q(2n|2m)`$ in $`U_q(𝔰𝔭𝔬(2n|2m))^{}`$ under the homomorphism $`\chi `$. Obviously, $`A`$ is the subalgebra of $`U_q(𝔰𝔭𝔬(2n|2m))^{}`$ that is generated by the elements $`\pi _{i,j}`$; $`i,jI`$. Actually, $`A`$ is a sub–Hopf–superalgebra of $`U_q(𝔰𝔭𝔬(2n|2m))^{}`$ (see the remark below Eqn. (2.7)). By definition, $`A`$ is dense in $`U_q(𝔰𝔭𝔬(2n|2m))^{\mathrm{gr}}`$ if and only if it separates the elements of $`U_q(𝔰𝔭𝔬(2n|2m))`$. Using the definition (5.14), we see that the implication
$$\phi (a,h)=0\text{for all }a\text{SPO}_q(2n|2m)h=0$$
(5.15)
is true for all $`hU_q(𝔰𝔭𝔬(2n|2m))`$ if and only if $`A`$ is dense in $`U_q(𝔰𝔭𝔬(2n|2m))^{\mathrm{gr}}`$. On the other hand, the implication
$$\phi (a,h)=0\text{for all }hU_q(𝔰𝔭𝔬(2n|2m))a=0$$
(5.16)
is true for all $`a\text{SPO}_q(2n|2m)`$ if and only if $`\chi `$ is injective (and hence induces a Hopf superalgebra isomorphism of $`\text{SPO}_q(2n|2m)`$ onto $`A`$). Unfortunately, at present I am not able to prove or disprove one or even both of the implications above. Actually, there is a simple reason to expect that (5.16) is not correct (see Section 7).
## 6 Graded comodule algebras over $`\text{SPO}_q(2n|2m)`$
We keep the notation of the preceding sections, in particular, that of Section 5. In the following, we are going to apply the general results of Section 4 to the case of the Hopf superalgebra $`\text{SPO}_q(2n|2m)`$.
To begin with, we recall some well–known facts (actually, these are true in the general framework of $`\sigma `$–bialgebras). If $`W`$ is a graded right $`\text{SPO}_q(2n|2m)`$–comodule, the dual pairing of Theorem 3 can be used to convert $`W`$ into a graded left $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`W^{\mathrm{mod}}`$. The procedure is the same as that in Section 2. Let
$$\delta :WW\text{SPO}_q(2n|2m)$$
be the structure map of $`W`$, and let $`hU_q(𝔰𝔭𝔬(2n|2m))`$ and $`xW`$ be homogeneous elements, of degrees $`\eta `$ and $`\xi `$, respectively. If
$$\delta (x)=\underset{^s}{}x_s^0x_s^1,$$
with $`x_s^0W`$ and $`x_s^1\text{SPO}_q(2n|2m)`$, the action of $`h`$ on $`x`$ is defined by
$$hx=\sigma (\eta ,\xi )\underset{^s}{}x_s^0\phi (x_s^1,h).$$
If $`U`$ is a second graded right $`\text{SPO}_q(2n|2m)`$–comodule, and if $`f:UW`$ is a homomorphism of graded $`\text{SPO}_q(2n|2m)`$–comodules, then $`f:U^{\mathrm{mod}}W^{\mathrm{mod}}`$ is a homomorphism of graded $`U_q(𝔰𝔭𝔬(2n|2m))`$–modules. In particular, any graded $`\text{SPO}_q(2n|2m)`$–subcomodule of $`W`$ is a graded $`U_q(𝔰𝔭𝔬(2n|2m))`$–submodule of $`W^{\mathrm{mod}}`$. Also, if $`W`$ is the graded tensor product of the graded right $`\text{SPO}_q(2n|2m)`$–comodules $`W_1`$, …, $`W_p`$, then $`W^{\mathrm{mod}}`$ is the graded tensor product of the graded left $`U_q(𝔰𝔭𝔬(2n|2m))`$–modules $`W_1^{\mathrm{mod}}`$, …, $`W_p^{\mathrm{mod}}`$. In addition, as the reader will expect, if $`W`$ is the vector comodule $`V`$ of $`\text{SPO}_q(2n|2m)`$, then $`W^{\mathrm{mod}}`$ is the vector module $`V`$ of $`U_q(𝔰𝔭𝔬(2n|2m))`$. Finally, if $`B`$ is a graded right $`\text{SPO}_q(2n|2m)`$–comodule algebra, the construction above converts $`B`$ into a graded left $`U_q(𝔰𝔭𝔬(2n|2m))`$–module algebra (i.e., the structure maps of the algebra $`B`$ are homomorphisms of graded $`U_q(𝔰𝔭𝔬(2n|2m))`$–modules).
There is an analogous construction that converts graded left $`\text{SPO}_q(2n|2m)`$–comodules into graded right $`U_q(𝔰𝔭𝔬(2n|2m))`$–modules, and its properties are completely analogous to those mentioned above. We leave it to the reader to spell out the details.
After these preliminaries, we are ready to discuss the graded right $`\text{SPO}_q(2n|2m)`$–comodule algebras. Guided by Section 4, we should look for graded $`\text{SPO}_q(2n|2m)`$–subcomodules of $`T(V)`$. We restrict our attention to quadratic comodule algebras. Correspondingly, we shall first determine the graded $`\text{SPO}_q(2n|2m)`$–subcomodules of
$$T_2(V)=T_0(V)T_1(V)T_2(V).$$
In view of the introductory remarks, we can do this in two steps: First, we determine the graded $`U_q(𝔰𝔭𝔬(2n|2m))`$–submodules of $`T_2(V)`$, then we try to show that the $`U_q(𝔰𝔭𝔬(2n|2m))`$–submodules we have found are, in fact, graded $`\text{SPO}_q(2n|2m)`$–subcomodules.
In the following, the $`\text{SPO}_q(2n|2m)`$–comodules and $`U_q(𝔰𝔭𝔬(2n|2m))`$–modules will simply be called comodules and modules, respectively.
It is easy to see that every graded submodule of $`T_2(V)`$ is equal to $`U`$ or $`UT_1(V)`$, where $`U`$ is a graded submodule of $`T_0(V)T_2(V)`$. (Of course, we are not interested in the second case, since it leads to trivial comodule algebras.) Similarly, every graded submodule of $`T_0(V)T_2(V)`$ is equal to $`W`$ or $`W(VV)_s`$, where $`W`$ is a graded submodule of $`T_0(V)(VV)_a`$ (we are using the notation of Ref. ).
Finally, every graded submodule $`W`$ of $`T_0(V)(VV)_a`$ belongs to one of the following two classes.
1) The module $`W`$ is equal to $`W_a`$ or $`T_0(V)W_a`$, where $`W_a`$ is a graded submodule of $`(VV)_a`$.
2) We have
$$W=𝕂(gc)(VV)_a^0,$$
where $`g`$ is any element of $`(VV)_a`$ that does not belong to $`(VV)_a^0`$, and where $`c𝕂=T_0(V)`$ is different from zero.
Our next task consists in checking which of these submodules are, in fact, graded subcomodules of $`T(V)`$. This is done by means of the following observations. (Unfortunately, in the case $`n=m=1`$, I haven’t solved this problem completely, see below.)
Obviously, $`T_1(V)`$ is a graded subcomodule of $`T(V)`$.
Secondly, we know (see Ref. ) that the projector of $`VV`$ onto $`(VV)_s`$ with kernel $`(VV)_a`$ can be written as a polynomial in $`\widehat{R}`$. Since $`\widehat{R}`$ (and hence the projector) is a homomorphism of graded comodules, this implies that $`(VV)_s`$ and $`(VV)_a`$ are graded subcomodules.
We know that $`\stackrel{~}{b}^q:VV𝕂`$ is a homomorphism of graded comodules (see Section 5), and that $`(VV)_a^0`$ is the kernel of the restriction of $`\stackrel{~}{b}^q`$ onto $`(VV)_a`$ (see Ref. ). Hence $`(VV)_a^0`$ is a graded subcomodule as well.
Remark 6.1. We note that in the case $`n=m`$ there does not exist a projector of $`VV`$ onto $`(VV)_a^0`$ that would be a homomorphism of graded comodules. Hence in this case the projector method of Ref. cannot be applied.
In Section 5 we have also noted that the element $`a`$ (see Eqn. (5.11)) is an invariant of the comodule $`VV`$. This shows that $`𝕂a`$ is a subcomodule of $`VV`$.
In the case $`n=m=1`$, we also have to check whether $`V_4`$ and $`\overline{V}_4`$ are graded subcomodules. Unfortunately, I haven’t been able to solve this problem, but I think there are good reasons to conjecture that they are not (see Section 7).
Finally, each of the graded comodules of case 2) above is a graded subcomodule as well. In fact, it is equal to the kernel of the linear form
$$f:T_0(V)(VV)_a𝕂$$
defined by
$$f(\lambda +u)=\stackrel{~}{b}^q(g)\lambda +c\stackrel{~}{b}^q(u)\text{for all }\lambda T_0(V)\text{}u(VV)_a,$$
which is a homomorphism of graded comodules since $`\stackrel{~}{b}^q`$ is such (recall that $`\stackrel{~}{b}^q(g)`$ is not equal to zero).
Since, for each of the graded subcomodules of $`T_2(V)`$ found above, a basis has been given in Ref. , the defining relations of the corresponding graded right $`\text{SPO}_q(2n|2m)`$–comodule algebra $`B_r`$ can immediately be written down.
In the following, we shall concentrate on the case 2) above, but allow for $`c=0`$ (i.e., we include the case $`W=(VV)_a`$). Note that, for $`c0`$, the subcomodule $`W`$ only depends on the residue class $`c^1g+(VV)_a^0`$. We fix our conventions by choosing
$$W=𝕂(tq^{2d}c)(VV)_a^0.$$
Here, $`t(VV)_a`$ is the element specified in Eqn. (5.17) of Ref. , and $`c𝕂`$ is an arbitrary scalar. For convenience, we have included the factor $`q^{2d}`$; recall that $`d=nm`$. The corresponding graded right $`\text{SPO}_q(2n|2m)`$–comodule algebra $`B_r(W)`$ will be denoted by $`A_q^r(n|m;c)`$. For $`c=0`$ it might be called the superalgebra of functions on the symplecto–orthogonal quantum superspace, for $`c0`$ it is the quantum Weyl superalgebra we wanted to construct (see Remark 6.2).
Using the basis of $`(VV)_a^0`$ given in Ref. and the formula for $`t`$ given there, this comodule superalgebra can be described in terms of generators and relations, as follows.
By definition, $`A_q^r(n|m;c)`$ is the universal superalgebra, generated by elements $`x_i`$; $`iI`$, which are homogeneous of degree $`\eta _i`$ and satisfy the following relations:
The basis elements $`a_{i,i}`$ with $`\eta _i=\overline{1}`$ give
$$x_i^2=0\text{for }\eta _i=\overline{1},$$
the basis elements $`a_{i,j}`$ with $`i<j`$ but $`ij`$ give
$$x_ix_j=\sigma _{i,j}qx_jx_i\text{for }i<j,ij,$$
the tensors $`a_j`$, $`2jr`$, give
$$q^{\sigma _{j1}}x_{j+1}x_{j1}\sigma _{j1}qx_{j1}x_{j+1}=\sigma _{j1}\sigma _jx_jx_j\sigma _{j1}qq^{\sigma _j}x_jx_j,$$
and the tensor $`tq^{2d}c`$ gives
$$x_1x_1q^2x_1x_1=q^{2d}c(qq^1)\underset{i=2}{\overset{n+m}{}}((C^q)^1)_{i,i}x_ix_i.$$
The last $`(r1)+1`$ relations are equivalent to the system
$$\begin{array}{cccc}\hfill \text{(I)}& \hfill q^{\sigma _i}x_ix_i\sigma _iq^1x_ix_i& =& (qq^1)\sigma _iC_{i,i}^q\underset{j<i}{}((C^q)^1)_{j,j}x_jx_j+\sigma _iC_{i,i}^qq^{2d}c,\hfill \end{array}$$
which, in turn, can be shown to be equivalent to the system
$$\begin{array}{cccc}\hfill \text{(II)}& \hfill q^{\sigma _i}x_ix_i\sigma _iqx_ix_i& =& (qq^1)\sigma _iC_{i,i}^q\underset{j<i}{}((C^q)^1)_{j,j}x_jx_j+\sigma _iC_{i,i}^qc.\hfill \end{array}$$
Both systems hold for $`ri1`$ (and the summation is over all $`jI`$ such that $`j<i`$).
Finally, the structure map
$$\delta :A_q^r(n|m;c)A_q^r(n|m;c)\overline{}\text{SPO}_q(2n|2m)$$
is fixed by Eqn. (4.7).
It can be shown that the canonical image of the invariant $`aVV`$ is given by
$$\underset{jI}{}((C^q)^1)_{j,j}x_jx_j=(q^dq^d)(qq^1)^1q^dc.$$
(6.1)
For $`n=m`$, this element is equal to zero, as it should (since in this case, we have $`a(VV)_a^0`$).
Using the braid generator $`\widehat{R}`$, the defining relations of $`A_q^r(n|m;c)`$ can be written in the following form. Let $`M:A_q^rA_q^rA_q^r`$ denote the multiplication map of $`A_q^r(n|m;c)`$. Then the defining relations are equivalent to
$$M(\widehat{R}q𝕀𝕀)(x_ix_j)=C_{i,j}^qc\text{for all }i,jI.$$
(6.2)
In the purely bosonic case, this neat formula is due to Kulish . By multiplying this equation by $`((C^q)^1)_{i,j}`$, summing over $`i,j`$, and using some formulae given in Ref. , we can easily rederive Eqn. (6.1).
In view of the conciseness of Eqn. (6.2), we might want to use this formula to define the algebra $`A_q^r(n|m;c)`$. Then we would have to show that the Eqns. (4.7) define on $`A_q^r(n|m;c)`$ the structure of a graded right $`\text{SPO}_q(2n|2m)`$–comodule algebra. This will be the case if we can show that the relations (6.2) stem from a graded subcomodule of $`T(V)`$, in the way described in Section 4. In order to prove this, we consider the linear map
$$F:T_2(V)T_0(V)T_2(V),$$
defined by
$$F(u)=\widehat{R}(u)quc\stackrel{~}{b}^q(u)\text{for all }uT_2(V).$$
Since $`\widehat{R}`$ and $`\stackrel{~}{b}^q`$ are homomorphisms of graded comodules, so is $`F`$, and its image is the graded subcomodule of $`T(V)`$ we are looking for.
We close this section by a series of remarks.
Remark 6.2. Suppose that $`𝕂=`$ (or, more generally, that $`𝕂`$ contains the square roots of all of its elements). Then, for fixed numbers $`m`$ and $`n`$, the graded right $`\text{SPO}_q(2n|2m)`$–comodule algebras $`A_q^r(n|m;c)`$ with $`c0`$ are all isomorphic. (This follows by simply rescaling the generators with a suitable overall factor.) Thus, with a slight abuse of language, we then may call this superalgebra the $`\text{SPO}_q(2n|2m)`$–covariant quantum Weyl superalgebra and denote it by $`W_q(n|m)`$. Needless to say, there are really innumerable papers on quantum oscillator (super)algebras of all types, and well–known methods to relate the various versions (see Ref. for an overview). Thus the main contribution of the present work to this field is to present an $`\text{SPO}_q(2n|2m)`$–covariant quantum oscillator superalgebra and to provide the mathematical basis of this concept (including the construction of the quantum supergroup $`\text{SPO}_q(2n|2m)`$ itself).
Remark 6.3. The defining relations of the algebra $`A_q^r(n|m;c)`$ can be simplified in various ways. For example, choose elements $`r_i𝕂`$; $`iI`$, such that
$$r_ir_i=((C^q)^1)_{i,i}\text{for }ri1,$$
and define
$$x_i^{}=r_ix_i\text{for all }iI.$$
Multiplying the $`i`$th relation of the system (I) by $`((C^q)^1)_{i,i}`$, we obtain the following equivalent system for the generators $`x_i^{}`$:
$$q^1x_i^{}x_i^{}\sigma _iq^{\sigma _i}x_i^{}x_i^{}=(qq^1)\underset{j<i}{}x_j^{}x_j^{}q^{2d}c\text{for }ri1.$$
Obviously, for $`ij`$ the generators $`x_i^{}`$ and $`x_j^{}`$ satisfy the same relation as $`x_i`$ and $`x_j`$. On the other hand, the system (II) is more complicated for the generators $`x_i^{}`$. Of course, we could also simplify the system (II), at the price that now the system (I) becomes more complicated.
Remark 6.4. In the case $`n=m=1`$, a $`U_q(𝔰𝔭𝔬(2|2))`$–covariant quantum Weyl superalgebra of the type considered in this section has been constructed in Ref. . However, contrary to the claim of the authors, they consider the other version of $`U_q(𝔰𝔭𝔬(2|2))`$, namely, the one associated to a basis of the root system consisting of two odd roots. For this version of $`U_q(𝔰𝔭𝔬(2|2))`$, I have gone through all the steps of Ref. and of the present work, and finally have found (after adequate adjustments) the quantum Weyl superalgebra of Ref. .
Remark 6.5. It should be clear that results similar to those obtained above can also be derived in the case of graded left $`\text{SPO}_q(2n|2m)`$–comodule algebras (and I have done that almost completely). Actually, most of the necessary technical tools have been mentioned in the present paper. Thus the interested reader should have no difficulties to carry out the details. He/she might then proceed to study the duality between the two classes of comodule superalgebras thus obtained.
## 7 Discussion
In the following, it will be useful to include the undeformed case $`q=1`$ in our discussion. The corresponding Hopf superalgebra $`\text{SPO}_1(2n|2m)`$ is easily described. For $`q=1`$, we have $`R=\text{id}_{VV}`$, hence $`A(R)`$ is the universal supercommutative algebra, generated by elements $`t_{ij}`$; $`i,jI`$, which are homogeneous of degree $`\eta _j\eta _i`$ . The Hopf superalgebra $`\text{SPO}_1(2n|2m)`$ is defined by requiring in addition the relations (5.8) and (5.9) (with $`q=1`$), or their matrix equivalents (see Remark 5.2).
A look at the supergroup $`\text{SPO}_1(2n|2m)`$ helps to understand how the quantum supergroup $`\text{SPO}_q(2n|2m)`$ should be interpreted. The Lie group “contained” in $`\text{SPO}_1(2n|2m)`$ is equal to $`\mathrm{SP}(2n)\times \mathrm{O}(2m)`$. Correspondingly, the algebra of functions $`\text{SPO}_1(2n|2m)`$ contains an element $`B=B_1`$, the Berezinian = superdeterminant, such that $`B^2=1`$, but $`B1`$. It is important to recall that the Berezinian is multiplicative. By setting $`B`$ equal to one, we obtain the algebra of functions on a supergroup, whose associated Lie group is equal to $`\mathrm{SP}(2n)\times \mathrm{SO}(2m)`$, and which, accordingly, might be called $`\text{SPSO}_1(2n|2m)`$.
It is tempting to conjecture that something analogous should be true in the deformed case. More precisely, there should exist an element $`B_q\text{SPO}_q(2n|2m)`$, which is group–like, homogeneous of degree zero, central, and satisfies $`B_q^2=1`$ but $`B_q1`$. (We might even hope to find a preimage of $`B_q`$ in $`A(R)`$ with analogous properties.) The quotient of $`\text{SPO}_q(2n|2m)`$ modulo the graded Hopf ideal generated by $`B_q1`$ would then be a Hopf superalgebra $`\text{SPSO}_q(2n|2m)`$. Moreover, we expect that the Hopf superalgebra pairing $`\phi `$ given by Theorem 3 satisfies
$$\phi (B_q1,h)=0\text{for all }hU_q(𝔰𝔭𝔬(2n|2m)),$$
which implies that $`\phi `$ induces a Hopf superalgebra pairing $`\phi _0`$ of $`\text{SPSO}_q(2n|2m)`$ and $`U_q(𝔰𝔭𝔬(2n|2m))`$. In particular, the pairing of Theorem 3 would not satisfy (5.16), and the best we could hope for is that $`\phi _0`$ would satisfy the implication analogous to (5.16). Similarly, we also expect that the universal $`r`$–form $`\varrho `$ on $`\text{SPO}_q(2n|2m)`$ described in Theorem 2 induces a universal $`r`$–form on $`\text{SPSO}_q(2n|2m)`$.
The foregoing is closely related to the unsolved problem of whether, in the case $`n=m=1`$, the $`U_q(𝔰𝔭𝔬(2|2))`$–submodules $`V_4`$ and $`\overline{V}_4`$ are $`\text{SPO}_q(2|2)`$–subcomodules of $`VV`$. It is easy to see that $`V_4`$ (resp. $`\overline{V}_4`$) is an $`\text{SPO}_q(2|2)`$–subcomodule of $`VV`$ if and only if $`t_{2,2}^2=0`$ (resp. $`t_{2,2}^2=0`$). (The fact that $`V_4`$ and $`\overline{V}_4`$ are $`U_q(𝔰𝔭𝔬(2|2))`$–submodules of $`VV`$ implies that the matrix elements $`\pi _{2,2}=\chi (t_{2,2})`$ and $`\pi _{2,2}=\chi (t_{2,2})`$ satisfy the equation $`\pi _{2,2}^2=\pi _{2,2}^2=0`$ in $`U_q(𝔰𝔭𝔬(2|2))^{}`$.) It can be shown that in $`\text{SPO}_q(2|2)`$ (even in $`A(R)`$) we have
$$t_{2,2}^2t_{2,2}^2=t_{2,2}^2t_{2,2}^2=t_{2,2}^2t_{2,2}^2=t_{2,2}^2t_{2,2}^2=\mathrm{\hspace{0.17em}0}.$$
In the case $`q=1`$, it is the relation $`B_1=1`$ which implies that $`t_{2,2}`$ and $`t_{2,2}`$ are invertible and hence that $`t_{2,2}^2=t_{2,2}^2=0`$. Accordingly, we conjecture that $`V_4`$ and $`\overline{V}_4`$ are $`\text{SPSO}_q(2|2)`$–subcomodules but not $`\text{SPO}_q(2|2)`$–subcomodules of $`VV`$.
The construction of the quantum Berezinian $`B_q`$ (if it exists) is a difficult problem. One may hope to solve it by use of a suitable Koszul complex. In the undeformed case, the method has been described in Ref. , the deformed case has been discussed in Ref. . Note, however, that in the latter reference, the case of Iwahori, Hecke type $`R`$–matrices is considered, whereas here we are in the Birman, Wenzl, Murakami setting.
We close this discussion by some general remarks. In the present work, we have studied the symplecto–orthogonal quantum supergroups $`\text{SPO}_q(2n|2m)`$ and their graded comodule algebras. Actually, the results of the Sections 3 and 4 have been derived under much more general assumptions. Moreover, also the discussion of the Sections 5 and 6 should be applicable to more general cases. A look into Section 5 shows that only some general properties of the $`R`$–matrix have been used. Thus, once the $`R`$–matrices of (some version of) the orthosymplectic quantum superalgebras in the vector representation have been calculated, it should be easy to extend the results of the Sections 5 and 6 to these cases.
The reader may have noticed that I have been careful not to use the duality of $`\text{SPO}_q(2n|2m)`$ and $`U_q(𝔰𝔭𝔬(2n|2m))`$ beyond the extent to which it is really established (see Theorem 3). This should help to investigate this duality more carefully.
We have not given a basis of the quantum Weyl superalgebra $`A_q^r(n|m;c)`$. Using the diamond lemma (see Ref. ), it should not be difficult to prove that (with respect to some suitable ordering of the indices) the ordered monomials in the generators $`x_i`$ form a basis of this algebra; of course, the exponent of $`x_i`$ should be restricted to $`\{0,1\}`$ if $`x_i`$ is odd.
Acknowledgements
Throughout the investigations which led to the present work, I have had various highly fruitful discussions with Ruibin Zhang on most of the topics considered here. I am very grateful to him for sharing with me his insights into the theory of quantum supergroups. Part of this research has been carried out during a visit of the author to the Department of Mathematics at the University of Queensland. The kind invitation by Mark Gould and the hospitality extended to the author by the members of the department are gratefully acknowledged.
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# 1 Introduction.
## 1 Introduction.
The search for kinks in N-component $`(1+1)`$-dimensional field theory is a very difficult endeavour, see Reference pp. 23-24. Mathematically, this problem is equivalent to a N-dimensional mechanical system via dimensional reduction to (0+1)-dimensions of space-time. Therefore, only if the mechanical system is completely integrable is there a hope of finding all the kinks of the field theory that are in one-to-one correspondence with the separatrix trajectories (developing finite action in infinite time) of the mechanical system. The authors have investigated this line of research on one-dimensional topological defects in several earlier papers, see .
After the seminal paper of Olive and Witten on extended supersymmetry algebras and solitons, , the paradigmatic N=1 kinks have been understood as BPS states of $`𝒩=1`$ supersymmetric quantum field theory. In $`(1+1)`$-dimensions, the N=1 theory is based in an extended supersymmetric algebra because there exist Majorana-Weyl fermions that could be used to generalize the purely bosonic theory supporting kinks to a Bose-Fermi system enjoying supersymmetry. Very recently, domain walls have been studied in the Wess-Zumino model, , and $`𝒩=1`$ SUSY QCD, . These theories arise in the low energy limit of string theory and the domain walls can be seen as membranes or other extended structures. On the other hand, both the WZ model and SQCD contain more than one real scalar field. Thus, their dimensional reduction is a $`(1+1)`$-dimensional either $`𝒩=1`$ or $`𝒩=2`$ supersymmetric field theory. Because the kinks are the actual domain walls if they are looked at from a $`(3+1)`$D point of view, in order to identify the BPS kinks (hence the domain walls, hence the d-branes) it is convenient to study $`𝒩=2`$ or $`𝒩=4`$ supersymmetric classical mechanics, starting from the finite action trajectories of these mechanical systems.
Accordingly, in this work we shall study the $`𝒩=2`$ supersymmetric extension of the N=2 Liouville systems. These supersymmetric extensions of completely integrable mechanical systems arise in the dimensional reduction of $`(1+1)`$-dimensional bosonic theories that have a very rich manifold of kinks. Our aim is to elucidate wheter these supersymmetric models in classical mechanics have a second invariant; if the answer is affirmative, then in priciple all the finite action trajectories that give the kinks and their supersymmetric partners can be found.
## 2 Supersymmetric Classical Mechanics
Let us start with a Grassman algebra $`B_L`$ with generators $`\theta _a^iB_L`$, $`a=1,2`$, $`i=1,\mathrm{},N`$, that satisfy the anti-commutation rules:
$$\theta _a^i\theta _b^j+\theta _b^j\theta _a^i=0$$
Any element of $`B_L`$ is a combination of the $`L=2N`$ generators $`\theta _a^i`$ that can be written in the form:
$$b=b_0\mathrm{𝟏}+b_{i_1\mathrm{}i_m}^{a_1\mathrm{}a_m}\theta _{a_1}^{i_1}\mathrm{}\theta _{a_m}^{i_m},$$
where the coefficients $`b_{i_1\mathrm{}i_m}^{a_1\mathrm{}a_m}`$ are real numbers. We shall distinguish between odd and even elements of $`B_L`$ according to the parity of the number of Grassman generators.
### 2.1 Lagrangian Formalism
The configuration space of the system is $`𝒞=^N\times (B_{2N})`$. If we choose $`(x^i,\theta _1^i,\theta _2^i)`$ as local coordinates in $`𝒞`$, the Lagrangian of our system has the “supernatural” form:
$$L=\frac{1}{2}\dot{x}^j\dot{x}^jU(\stackrel{}{x})+\frac{i}{2}\theta _a^j\dot{\theta }_a^j+iR_{jk}(\stackrel{}{x})\theta _1^j\theta _2^k$$
Observe that the Lagrangian is defined on even elements of $`𝒞`$ of two types: there is a bosonic contribution collecting the kinetic and potential energies of the $`x^i`$ coordinates. There are also two terms that contain fermionic variables: besides the Grassman kinetic energy, there is a Yukawa coupling $`R_{jk}`$ that we assume to be symmetric in the $`j,k`$ indices, see . The Euler-Lagrange equations
$$\ddot{x}^k+\frac{U}{x^k}i\frac{R_{jl}}{x^k}\theta _1^j\theta _2^l=0\dot{\theta }_1^i=R_{ij}\theta _2^j\dot{\theta }_2^i=R_{ji}\theta _1^j$$
determine the dynamics of our mechanical system. Via Noether’s theorem, one finds the energy as the invariant associated to invariance under time translations:
$$I=H=\frac{1}{2}\dot{x}^j\dot{x}^j+U(\stackrel{}{x})iR_{jk}(\stackrel{}{x})\theta _1^j\theta _2^k$$
(1)
### 2.2 Hamiltonian Formalism
We shall now briefly discuss the Hamiltonian formalism . The usual definition of generalized momentum is extended to the Grassman variables:
$$\pi _{\theta _a^j}=L\frac{\stackrel{}{}}{\dot{\theta }_a^j}=\frac{i}{2}\theta _a^j$$
(2)
We note the dependence of the fermionic generalized momenta on their Grassman variables. In the $`6N`$-dimensional phase space $`T^{}𝒞`$ with coordinates $`(x^i,\theta _1^i,\theta _2^i,p_i,\pi _{\theta _1^i},\pi _{\theta _2^i})`$, the definition of Grassman generalized momenta (2) provides $`2N`$ second class constraints. Instead of working in the reduced $`4N`$-dimensional phase space (the Grassman sub-space of the phase space coincides with the Grassman sub-space of the configuration space because the Lagrangian is first-order in the velocities) we implement the constraints in the Hamiltonian by means of Grassman-Lagrange multipliers :
$$H_T=\frac{1}{2}\dot{x}^j\dot{x}^j+U(\stackrel{}{x})iR_{jk}(\stackrel{}{x})\theta _1^j\theta _2^k(\pi _{\theta _a^j}\frac{i}{2}\theta _a^j)\lambda _a^j$$
The Hamilton equations are:
$$\dot{x}^j=\frac{H_T}{p_j},\dot{p}_j=\frac{H_T}{x^j},\dot{\theta }_a^j=\frac{H_T}{\pi _{\theta _a^j}},\dot{\pi }_{\theta _a^j}=\frac{H_T}{\theta _a^j}$$
Note the difference in sign between the bosonic and fermionic canonical equations. Solving for the Lagrange multipliers, we find
$$H_T=\frac{1}{2}\dot{x}^j\dot{x}^j+U(\stackrel{}{x})+R_{jk}(\stackrel{}{x})(\pi _{\theta _2^k}\theta _1^j\pi _{\theta _1^j}\theta _2^k),$$
which provides the right Hamiltonian flow in the phase space.
Defining the Poisson brackets for the generic functions $`F`$ and $`G`$ in phase space in the usual way:
$`\{F,G\}_P`$ $`=`$ $`{\displaystyle \frac{F}{p_j}}{\displaystyle \frac{G}{q^j}}{\displaystyle \frac{F}{q^j}}{\displaystyle \frac{G}{p_j}}+iF{\displaystyle \frac{\stackrel{}{}}{\theta _a^j}}{\displaystyle \frac{\stackrel{}{}}{\theta _a^j}}G{\displaystyle \frac{1}{2}}F{\displaystyle \frac{\stackrel{}{}}{\pi _{\theta _a^j}}}{\displaystyle \frac{\stackrel{}{}}{\theta _a^j}}G`$
$`{\displaystyle \frac{1}{2}}F{\displaystyle \frac{\stackrel{}{}}{\theta _a^j}}{\displaystyle \frac{\stackrel{}{}}{\pi _{\theta _a^j}}}G{\displaystyle \frac{i}{4}}F{\displaystyle \frac{\stackrel{}{}}{\pi _{\theta _a^j}}}{\displaystyle \frac{\stackrel{}{}}{\pi _{\theta _a^j}}}G,`$
the canonical equations read:
$$\begin{array}{ccc}\frac{dx^j}{dt}=\{H_T,x^j\}_P\hfill & & \frac{d\theta _a^j}{dt}=\{H_T,\theta _a^j\}_P\hfill \\ \frac{dp_j}{dt}=\{H_T,p_j\}_P\hfill & & \frac{d\pi _{\theta _a^j}}{dt}=\{H_T,\pi _{\theta _a^j}\}_P\hfill \end{array}$$
In general, the time-dependence of any observable $`F`$ is ruled by:
$$\frac{dF}{dt}=\{H_T,F\}_P$$
Therefore, physical observables are constants of motion or invariants if and only if:
$$\{H_T,I\}_P=0$$
In practical terms, it is better to work on the reduced phase space and define the reduced Poisson brackets as:
$$\{F,G\}_P=\frac{F}{p_j}\frac{G}{q^j}\frac{F}{q^j}\frac{G}{p_j}+iF\frac{\stackrel{}{}}{\theta _a^j}\frac{\stackrel{}{}}{\theta _a^j}G$$
Thus, we receive the following Poisson structure:
$$\{p_j,x^k\}_P=\delta _j^k\{x^j,x^k\}_P=\{p_j,p_k\}_P=0\{\theta _a^j,\theta _b^k\}_P=i\delta ^{jk}\delta _{ab}$$
and the canonical equations and the invariant observables are referred to the reduced Hamiltonian $`H`$ given in (1).
### 2.3 Supersymmetry
The question arises: are there transformations in the configuration or phase space such that they mix the bosonic and fermionic variables and leave invariant the Lagrangian or the Hamiltonian?. If the answer is yes, then the mechanical system can be said to enjoy supersymmetry . Instead of using the elegant superfield/superspace formalism, we take a direct approach. Consider the following infinitesimal variations in the configuration space defined in terms of a Grassman parameter $`ϵ`$:
$$\begin{array}{ccc}\begin{array}{c}\text{Variation 1:}\\ \{\begin{array}{c}\delta _1x^j=\epsilon \theta _1^j\hfill \\ \delta _1\theta _1^j=i\epsilon \dot{x}^j\hfill \\ \delta _1\theta _2^j=i\epsilon f^j(\stackrel{}{x})\hfill \end{array}\end{array}\hfill & & \begin{array}{c}\text{Variation 2:}\\ \{\begin{array}{c}\delta _2x^j=\epsilon \theta _2^j\hfill \\ \delta _2\theta _1^j=i\epsilon g^j(\stackrel{}{x})\hfill \\ \delta _2\theta _2^j=i\epsilon \dot{x}^j\hfill \end{array}\end{array}\hfill \end{array}$$
The variations induced on the Lagrangian are:
$`\delta _1L`$ $`=`$ $`{\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{1}{2}}\dot{x}^j\epsilon \theta _1^j\right)+{\displaystyle \frac{1}{2}}f^j\epsilon \dot{\theta }_2^j\left(R_{kj}\dot{x}^k+{\displaystyle \frac{1}{2}}\dot{f}^j\right)\epsilon \theta _2^j`$
$``$ $`\left({\displaystyle \frac{V}{x^j}}+R_{jk}f^k\right)ϵ\theta _1^j+i{\displaystyle \frac{R_{jk}}{x^l}}\epsilon \theta _1^l\theta _1^j\theta _2^k`$
$`\delta _2L`$ $`=`$ $`{\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{1}{2}}\dot{x}^j\epsilon \theta _2^j\right){\displaystyle \frac{1}{2}}g^j\epsilon \dot{\theta }_1^j+\left(R_{kj}\dot{x}^k+{\displaystyle \frac{1}{2}}\dot{g}^j\right)\epsilon \theta _1^j`$
$``$ $`\left({\displaystyle \frac{V}{x^j}}+R_{jk}g^k\right)ϵ\theta _2^j+i{\displaystyle \frac{R_{jk}}{x^l}}\epsilon \theta _2^l\theta _1^j\theta _2^k`$
There is symmetry with respect to the variations $`\delta _1`$ and $`\delta _2`$ if and only if $`\delta _1L`$ and $`\delta _2L`$ are pure divergences. This happens if
$$U(\stackrel{}{x})=\frac{1}{2}\frac{W}{x^j}\frac{W}{x^j}R_{jk}=\frac{^2W}{x^jx^k}f^j=g^j=\frac{W}{x^j},$$
where $`W(\stackrel{}{x})`$ is a function defined in the bosonic piece of the configuration space called the superpotential. A supersymmetric Lagrangian has the form,
$$L=\frac{1}{2}\dot{x}^j\dot{x}^j+\frac{i}{2}\theta _a^j\dot{\theta }_a^j\frac{1}{2}\frac{W}{x^j}\frac{W}{x^j}i\frac{^2W}{x^jx^k}\theta _1^j\theta _2^k$$
The Noether charges associated with the symmetries with respect to $`\delta _1`$ and $`\delta _2`$ are respectively
$$Q_1=\dot{x}^j\theta _1^j\frac{W}{x^j}\theta _2^jQ_2=\dot{x}^j\theta _2^j+\frac{W}{x^j}\theta _1^j$$
Going to Hamiltonian formalism, one easily checks that these fermionic supercharges close the $`N=2`$ superalgebra
$$\{Q_1,Q_1\}_P=2H\{Q_1,Q_2\}_P=0\{Q_2,Q_2\}_P=2H$$
We immediately find a bosonic invariant, the Hamiltonian $`H`$ itself, and two fermionic constants of motion- the supercharges $`Q_1`$ and $`Q_2`$; their Poisson brackets with $`H`$ are zero. If the number of bosonic degrees of freedom is $`N`$, it is easy to check that there is other bosonic invariant, $`I_3=_{i=1}^N\theta _1^i\theta _2^i`$ .
If the bosonic piece of the configuration space is a general Riemannian manifold $`M^N`$ equipped with a metric tensor $`g_{ij}`$, the Grassman variables $`\vartheta _a^i`$ are the components of a contravariant vector. The definition of the supercharges is generalized in the form:
$$Q_1=g_{jk}\dot{x}^j\vartheta _1^k\frac{W}{x^j}\vartheta _2^jQ_2=g_{jk}\dot{x}^j\vartheta _2^k+\frac{W}{x^j}\vartheta _1^j$$
The supersymmetric algebra dictates the form of the Hamiltonian
$$H=\frac{1}{2}g_{jk}\dot{x}^j\dot{x}^k+\frac{1}{2}g^{jk}\frac{W}{x^j}\frac{W}{x^k}+iW_{j;k}\vartheta _1^j\vartheta _2^k,$$
if
$$W_{j;k}=\frac{^2W}{x^jx^k}\mathrm{\Gamma }_{jk}^l\frac{W}{x^l},$$
and the inverse Legendre transformation leads to the Lagrangian
$$L=\frac{1}{2}g_{jk}\dot{x}^j\dot{x}^k+\frac{i}{2}g_{jk}\vartheta _a^jD_t\vartheta _a^k+\frac{1}{4}R_{jkln}\vartheta _1^j\vartheta _2^l\vartheta _1^k\vartheta _2^n\frac{1}{2}g^{jk}\frac{W}{x^j}\frac{W}{x^k}iW_{j;k}\vartheta _1^j\vartheta _2^k,$$
(3)
where the covariant derivative is defined as $`D_t\vartheta _a^j=\dot{\vartheta }_a^j+\mathrm{\Gamma }_{lk}^j\dot{x}^l\theta _a^k`$.
## 3 From Liouville to SuperLiouville Models
We shall focus on mechanical systems of two bosonic degrees of freedom; $`N=2`$. In particular, we shall analyze the supersymmetric extensions of the classical Liouville models. These models are completely integrable and their classical invariants are well known . Our goal is to study what the invariants are in their supersymmetric extensions.
In fact, the classical Liouville models are not only completely integrable but Hamilton-Jacobi separable. The Hamilton-Jacobi principle thus provides all the solutions of the dynamics in these models. This is achieved by using appropriate coordinate systems. There are four possibilities:
$``$ Liouville Models of Type I: Let us consider the map $`\xi ^{}:D^2`$, where $`D`$ is an open sub-set of $`^2`$ with coordinates $`(u,v)`$. These variables are the elliptic coordinates of the bosonic system if the map is defined as: $`\xi ^{}(x^1)=\frac{1}{\mathrm{\Omega }}uv`$, $`\xi ^{}(x^2)=\frac{1}{\mathrm{\Omega }}\sqrt{(u^2\mathrm{\Omega }^2)(\mathrm{\Omega }^2v^2)}`$ and $`u[\mathrm{\Omega },\mathrm{})`$, $`v[\mathrm{\Omega },\mathrm{\Omega }]`$. In the new variables, the Lagrangian of a Liouville model of Type I reads:
$$L=\frac{1}{2}\frac{u^2v^2}{u^2\mathrm{\Omega }^2}\dot{u}\dot{u}+\frac{1}{2}\frac{u^2v^2}{\mathrm{\Omega }^2v^2}\dot{v}\dot{v}\frac{u^2\mathrm{\Omega }^2}{u^2v^2}f(u)\frac{\mathrm{\Omega }^2v^2}{u^2v^2}g(v)$$
(4)
Observe that apart from a common factor the contribution to the Lagrangian of the $`u`$ and $`v`$ variables splits completely. The common factor can be interpreted as a metric (of zero curvature): $`g_{ij}=(u^2v^2)\delta _{ij}`$.
$``$ Liouville Models of Type II: In polar coordinates $`\zeta ^{}(x^1)=\rho \mathrm{cos}\chi `$, $`\zeta ^{}(x^2)=\rho \mathrm{sin}\chi `$, $`\rho [0,\mathrm{})`$, $`\chi [0,2\pi )`$ the Lagrangian of the Liouville models of Type II reads:
$$L=\frac{1}{2}\dot{\rho }\dot{\rho }+\frac{1}{2}\rho ^2\dot{\chi }\dot{\chi }f(\rho )\frac{1}{\rho ^2}g(\chi )$$
(5)
Again, besides the metric factor $`g_{11}=1`$, $`g_{22}=\rho ^2`$, the contributions of $`\rho `$ and $`\chi `$ appear completely separated in the Lagrangian.
$``$ Liouville Models of Type III: Parabolic coordinates $`u(\mathrm{},\mathrm{})`$, $`v[0,\mathrm{})`$ are defined through the map $`\gamma ^{}:D^2`$ such that $`\gamma ^{}(x^1)=\frac{1}{2}(u^2v^2)`$, $`\gamma ^{}(x^2)=uv`$. A Liouville model of Type III obeys a Lagrangian of the form:
$$L=\frac{1}{2}(u^2+v^2)\left(\dot{u}\dot{u}+\dot{v}\dot{v}\right)\frac{1}{u^2+v^2}\left(f(u)+g(v)\right)$$
(6)
There is a metric factor $`g_{ij}=(u^2+v^2)\delta _{ij}`$ and separate contributions of $`u`$ and $`v`$ to $`L`$.
$``$ Liouville Models of Type IV: In these models the Lagrangian is directly separated in Cartesian coordinates
$$L=\frac{1}{2}\dot{x}^1\dot{x}^1+\frac{1}{2}\dot{x}^2\dot{x}^2f(x^1)g(x^2)$$
(7)
and a Euclidean metric $`g_{ij}=\delta _{ij}`$ can be understood.
The definition of SuperLiouville models is a two step process:
$`(i)`$ Define a supersymmetric N=2 Lagrangian system on a Riemannian manifold $`M^2`$, which is $`^2`$ equipped with the metric induced by the maps $`\chi ^{}`$, $`\zeta ^{}`$ and $`\gamma ^{}`$ for Types I, II, and III, and the Euclidean metric for Type IV. Consider also Grassman variables that transform as $`\vartheta _a^i=\frac{x_{}^{}{}_{}{}^{i}}{x^j}\theta _a^j`$ under these changes of coordinates.
$`(ii)`$ A model defined in this way is a SuperLiouville model if the superpotential splits in such a manner that the bosonic part of the Lagrangian coincides with the Lagrangian of a Liouville model.
$``$ SuperLiouville Models of Type I: The Lagrangian of the bosonic sector of this Type of model contains two contributions:
$$L_B=\frac{1}{2}\frac{u^2v^2}{u^2\mathrm{\Omega }^2}\dot{u}\dot{u}+\frac{1}{2}\frac{u^2v^2}{\mathrm{\Omega }^2v^2}\dot{v}\dot{v}\frac{1}{2}\frac{u^2\mathrm{\Omega }^2}{u^2v^2}\left(\frac{W}{u}\right)^2\frac{1}{2}\frac{\mathrm{\Omega }^2v^2}{u^2v^2}\left(\frac{W}{v}\right)^2,$$
where the superpotential $`W`$ provides the potential $`U_B`$ through the identity,
$$U_B=\frac{1}{2}\frac{u^2\mathrm{\Omega }^2}{u^2v^2}\left(\frac{W}{u}\right)^2+\frac{1}{2}\frac{\mathrm{\Omega }^2v^2}{u^2v^2}\left(\frac{W}{v}\right)^2$$
In the fermionic sector the Lagrangian takes the form $`L_F=T_F+L_{BF}^I`$, with
$$T_F=\frac{i}{2}\frac{u^2v^2}{u^2\mathrm{\Omega }^2}\vartheta _a^uD_t\vartheta _a^u+\frac{i}{2}\frac{u^2v^2}{\mathrm{\Omega }^2v^2}\vartheta _a^vD_t\vartheta _a^v$$
The Bose-Fermi interaction adds to the Lagrangian the Yukawa terms
$`L_{FB}^I`$ $`=`$ $`i\left[{\displaystyle \frac{^2W}{uu}}+{\displaystyle \frac{\mathrm{\Omega }^2v^2}{(u^2v^2)(u^2\mathrm{\Omega }^2)}}\left(u{\displaystyle \frac{W}{u}}v{\displaystyle \frac{W}{v}}\right)\right]\vartheta _1^u\vartheta _2^u`$
$`i\left[{\displaystyle \frac{^2W}{uv}}+{\displaystyle \frac{1}{u^2v^2}}\left(v{\displaystyle \frac{W}{u}}u{\displaystyle \frac{W}{v}}\right)\right](\vartheta _1^u\vartheta _2^v+\vartheta _1^v\vartheta _2^u)`$
$`i\left[{\displaystyle \frac{^2W}{vv}}+{\displaystyle \frac{(u^2\mathrm{\Omega }^2)}{(u^2v^2)(\mathrm{\Omega }^2v^2)}}\left(u{\displaystyle \frac{W}{u}}v{\displaystyle \frac{W}{v}}\right)\right]\vartheta _1^v\vartheta _2^v`$
Definition: A system in supersymmetric classical mechanics is a SuperLiouville model of Type I if the map given by the change from Cartesian to elliptic coordinates acting on the Cartesian superpotential is such that:
$$\xi ^{}W_C=W_1(u)\pm W_2(v)$$
$``$ SuperLiouville Models of Type II: The bosonic Lagrangian is:
$$L_B=\frac{1}{2}\dot{\rho }\dot{\rho }+\frac{1}{2}\rho ^2\dot{\chi }\dot{\chi }\frac{1}{2}\left(\frac{W}{\rho }\right)^2\frac{1}{2\rho ^2}\left(\frac{W}{\chi }\right)^2$$
and the superpotential $`W`$ is related to the potential through
$$U_B=\frac{1}{2}\left(\frac{W}{\rho }\right)^2+\frac{1}{2\rho ^2}\left(\frac{W}{\chi }\right)^2$$
In the fermionic sector, the Lagrangian takes the form:
$$T_F=\frac{i}{2}\vartheta _a^\rho D_t\vartheta _a^\rho +\frac{i}{2}\rho ^2\vartheta _a^\chi D_t\vartheta _a^\chi ,$$
and the Yukawa Bose-Fermi couplings are:
$`L_{FB}^I`$ $`=`$ $`i{\displaystyle \frac{^2W}{\rho \rho }}\vartheta _1^\rho \vartheta _2^\rho i\left({\displaystyle \frac{^2W}{\chi \chi }}+\rho {\displaystyle \frac{W}{\rho }}\right)\vartheta _1^\chi \vartheta _2^\chi `$
$`i\left({\displaystyle \frac{^2W}{\rho \chi }}{\displaystyle \frac{1}{\rho }}{\displaystyle \frac{W}{\chi }}\right)(\vartheta _1^\rho \vartheta _2^\chi +\vartheta _1^\chi \vartheta _2^\rho )`$
Definition: a system in supersymmetric classical mechanics is a SuperLiouville Model of Type II if the map given by the change from Cartesian to polar coordinates acting on the Cartesian superpotential produces:
$$\varsigma ^{}W_C=W_1(\rho )\pm W_2(\chi )$$
$``$ SuperLiouville Models of Type III: The bosonic Lagrangian is:
$$L_B=\frac{1}{2}(u^2+v^2)\left(\dot{u}\dot{u}+\dot{v}\dot{v}\right)\frac{1}{2(u^2+v^2)}\left[\left(\frac{W}{u}\right)^2+\left(\frac{W}{v}\right)^2\right]$$
The potential is determined from the superpotential $`W`$ as:
$$U_B=\frac{1}{2(u^2+v^2)}\left[\left(\frac{W}{u}\right)^2+\left(\frac{W}{v}\right)^2\right]$$
The purely fermionic contribution to the Lagrangian is:
$$T_F=\frac{i}{2}(u^2+v^2)\vartheta _a^uD_t\vartheta _a^u+\frac{i}{2}(u^2+v^2)\vartheta _a^vD_t\vartheta _a^v$$
and the Yukawa couplings are:
$`L_{BF}^I`$ $`=`$ $`i\left[{\displaystyle \frac{^2W}{uu}}{\displaystyle \frac{1}{u^2+v^2}}\left(u{\displaystyle \frac{W}{u}}v{\displaystyle \frac{W}{v}}\right)\right]\vartheta _1^u\vartheta _2^u`$
$`i\left[{\displaystyle \frac{^2W}{uv}}{\displaystyle \frac{1}{u^2+v^2}}\left(u{\displaystyle \frac{W}{v}}+v{\displaystyle \frac{W}{u}}\right)\right](\vartheta _1^u\vartheta _2^v+\vartheta _1^v\vartheta _2^u)`$
$`i\left[{\displaystyle \frac{^2W}{vv}}+{\displaystyle \frac{1}{u^2+v^2}}\left(u{\displaystyle \frac{W}{u}}v{\displaystyle \frac{W}{v}}\right)\right]\vartheta _1^v\vartheta _2^v`$
Definition: A system in supersymmetric classical mechanics is a SuperLiouville Model of Type III if the map given by the change from Cartesian to parabolic coordinates acting on the Cartesian superpotential is such that:
$$\rho ^{}W_C=W_1(u)\pm W_2(v)$$
$``$ SuperLiouville Models of Type IV: Finally, the definition of SuperLiouville Model of Type IV is straightforward.
Definition: A system in supersymmetric classical mechanics belongs to Type IV SuperLiouville models if the superpotential $`W(x^1,x^2)`$ is of the form:
$$W(x^1,x^2)=W_1(x^1)\pm W_2(x^2)$$
The Lagrangian is:
$$L=\frac{1}{2}\dot{x}^j\dot{x}^j+\frac{i}{2}\theta _a^j\dot{\theta }_a^j\frac{1}{2}\frac{W}{x^j}\frac{W}{x^j}i\frac{^2W}{x^1x^1}\theta _1^1\theta _2^1i\frac{^2W}{x^2x^2}\theta _1^2\theta _2^2,$$
and the system can be be understood as a $`𝒩=2𝒩=2`$ SUSY in $`(0+1)`$ dimensions.
As a common feature, observe that the potential is insensitive to the relative signs of the separated parts of the superpotential. Therefore, all the Liouville models are supersymmetrizable by means of two different superpotentials.
## 4 On the Bosonic Invariants
It is well known that Liouville models have a second invariant in involution with the energy -the first invariant- that guarantees complete integrability in the sense of the Liouville theorem. We shall now show that SuperLiouville models also have a second invariant of bosonic nature. Our strategy in the search of such an invariant, $`[I,H]=0`$, follows the general pattern shown in the literature: see . The ansatz for invariants at the highest quadratic level in the momenta is:
$$\begin{array}{ccc}I\hfill & =& \hfill \frac{1}{2}H^{ij}p_ip_j+V(x_1,x_2)+F_{ij}\theta _1^i\theta _2^j+G_{ij}\theta _1^i\theta _1^j+J_{ij}\theta _2^i\theta _2^j+\\ & & \hfill +L_{jk}^ip_i\theta _1^j\theta _1^k+M_{jk}^ip_i\theta _2^j\theta _2^k+N_{jk}^ip_i\theta _1^j\theta _2^k+S_{ijkl}\theta _1^i\theta _2^k\theta _1^j\theta _2^l\end{array}$$
Here, we assume that:
\- $`i)`$ $`H^{ij}`$ is a symmetric tensor depending on $`x^i`$. There are three independent functions to determine.
\- $`ii)`$ $`L_{jk}^i`$ and $`M_{jk}^i`$ also depend only on $`x^i`$ and are antisymmetric in the indices $`j`$ and $`k`$: $`L_{jk}^i=L_{kj}^i`$, $`M_{jk}^i=M_{kj}^i`$. They include four independent functions.
\- $`iii)`$ $`G_{ij}`$ and $`J_{ij}`$ are antisymmetric functions of $`x^i`$ in the indices: $`G_{ij}=G_{ji}`$ and $`J_{ij}=J_{ji}`$. $`F_{ij}(x^i)`$, however, is neither symmetric nor antisymmetric; it contains four independent functions.
\- $`iv)`$ Finally, $`S_{ijkl}(x^i)`$ is antisymmetric in the exchange of the indices $`i,j`$ and $`k,l`$ and symmetric in the exchange of the pairs $`ij,kl`$. There is only one independent function to determine in this tensor.
The commutator with the Hamiltonian is:
$`[I,H]={\displaystyle \frac{1}{2}}{\displaystyle \frac{H^{jk}}{x^l}}p_lp_jp_k+\left(H^{lj}{\displaystyle \frac{^2W}{x^jx^k}}{\displaystyle \frac{W}{x_k}}+{\displaystyle \frac{V}{x_l}}\right)p_l+`$
$`+\left(iH^{nj}{\displaystyle \frac{^3W}{x^kx^lx^j}}+{\displaystyle \frac{F_{kl}}{x_i}}+2L_{km}^i{\displaystyle \frac{^2W}{x_mx_l}}+2M_{ln}^i{\displaystyle \frac{^2W}{x_nx^k}}\right)p_n\theta _1^k\theta _2^l+`$
$`+\left({\displaystyle \frac{^2W}{x_lx^k}}F_{lj}M_{kj}^n{\displaystyle \frac{^2W}{x_nx_l}}{\displaystyle \frac{W}{x^l}}\right)\theta _2^k\theta _2^j+\left({\displaystyle \frac{J_{lj}}{x_k}}+N_{mj}^k{\displaystyle \frac{^2W}{x_mx^l}}\right)p_k\theta _2^l\theta _2^j`$
$`\left({\displaystyle \frac{^2W}{x^jx_k}}F_{nk}+L_{nj}^l{\displaystyle \frac{^2W}{x^lx^k}}{\displaystyle \frac{W}{x_k}}\right)\theta _1^n\theta _1^j+\left({\displaystyle \frac{G_{nj}}{x_l}}N_{nk}^l{\displaystyle \frac{^2W}{x_jx_k}}\right)p_l\theta _1^n\theta _1^j+`$
$`+\mathrm{\hspace{0.17em}2}\left(G_{nj}{\displaystyle \frac{^2W}{x_jx^k}}+J_{kl}{\displaystyle \frac{^2W}{x^nx_l}}{\displaystyle \frac{1}{2}}N_{nk}^j{\displaystyle \frac{^2W}{x^jx^l}}{\displaystyle \frac{W}{x_l}}\right)\theta _1^n\theta _2^k+`$
$`+{\displaystyle \frac{L_{jk}^n}{x_l}}p_lp_n\theta _1^j\theta _1^k+{\displaystyle \frac{M_{jk}^n}{x_l}}p_lp_n\theta _2^j\theta _2^k+{\displaystyle \frac{N_{jk}^n}{x_l}}p_lp_n\theta _1^j\theta _2^k+`$
$`iN_{jk}^n{\displaystyle \frac{^3W}{x^nx^lx^m}}\theta _1^j\theta _2^k\theta _1^l\theta _2^mi{\displaystyle \frac{S_{njkl}}{x_m}}p_m\theta _1^n\theta _2^k\theta _1^j\theta _2^l`$
Therefore, $`[I,H]=0`$, and $`I`$ is a second invariant if and only if the following equations are satisfied:
| BOX 1 | a) $`{\displaystyle \frac{H^{ij}}{x^k}}+{\displaystyle \frac{H^{kj}}{x^i}}=0`$ |
| --- | --- |
| BOX 2 | a) $`H^{ij}{\displaystyle \frac{^2W}{x^jx^k}}{\displaystyle \frac{W}{x_k}}={\displaystyle \frac{V}{x_i}}`$ |
| BOX 3 | a) $`ϵ^{jk}{\displaystyle \frac{L_{jk}^i}{x_l}}+ϵ^{jk}{\displaystyle \frac{L_{jk}^l}{x_i}}=0`$ |
| | b) $`ϵ^{jk}{\displaystyle \frac{M_{jk}^i}{x_l}}+ϵ^{jk}{\displaystyle \frac{M_{jk}^l}{x_i}}=0`$ |
| BOX 4 | a) $`H^{nj}{\displaystyle \frac{^3W}{x^kx^lx^j}}+i{\displaystyle \frac{F_{kl}}{x_n}}+2iL_{km}^n{\displaystyle \frac{^2W}{x_mx_l}}+2iM_{lj}^n{\displaystyle \frac{^2W}{x_jx^k}}=0`$ |
| | b) $`ϵ^{ij}\left({\displaystyle \frac{^2W}{x_ix^k}}F_{kj}M_{ij}^k{\displaystyle \frac{^2W}{x_kx_l}}{\displaystyle \frac{W}{x^l}}\right)=0`$ |
| | c) $`ϵ^{ij}\left({\displaystyle \frac{^2W}{x^jx_k}}F_{ik}+L_{ij}^l{\displaystyle \frac{^2W}{x^kx^l}}{\displaystyle \frac{W}{x_k}}\right)=0`$ |
| BOX 5 | a) $`ϵ^{ij}\left({\displaystyle \frac{G_{ij}}{x_l}}N_{jk}^l{\displaystyle \frac{^2W}{x_ix_k}}\right)=0`$ |
| | b) $`ϵ^{ij}\left({\displaystyle \frac{J_{ij}}{x_k}}+N_{mj}^k{\displaystyle \frac{^2W}{x_mx^i}}\right)=0`$ |
| | c) $`G_{ij}{\displaystyle \frac{^2W}{x_jx^k}}+J_{kl}{\displaystyle \frac{^2W}{x^ix_l}}{\displaystyle \frac{1}{2}}N_{ik}^j{\displaystyle \frac{^2W}{x^jx^l}}{\displaystyle \frac{W}{x_l}}=0`$ |
| | d) $`{\displaystyle \frac{N_{jk}^i}{x_l}}+{\displaystyle \frac{N_{jk}^l}{x_i}}=0`$ |
| | e) $`ϵ^{ij}ϵ^{lk}N_{jk}^m{\displaystyle \frac{^3W}{x^ix^lx^m}}=0`$ |
| BOX 6 | a) $`ϵ^{ij}ϵ^{kl}{\displaystyle \frac{S_{ijkl}}{x_m}}=0`$ |
### 4.1 General properties of the solution
We deal with a overdetermined system of partial differential equations: there are 31 PDE relating 15 unknown functions. Moreover, some sub-systems can be solved for some sub-set of functions. We proceed in a recurrent way:
\- $`i)`$ The equations in BOXES 1 and 2 are sufficient to find $`H^{ij}`$ and $`V`$. We recover the information about the second invariant of the purely bosonic sector: the Liouville model.
\- $`ii)`$ The equations in BOX 3 are solved if the independent components of $`L_{jk}^i`$ and $`M_{jk}^i`$ have the form,
$$L_{12}^i=Cϵ^{ij}x_j+A_iM_{12}^i=Dϵ^{ij}x_j+B_i,$$
where $`A_i`$, $`B_i`$, $`C`$ y $`D`$ are constants.
\- $`iii)`$ The equations in BOX 4, together with the previous information, leads to the computation of $`F_{ij}`$. The identity $`\frac{^2F_{kl}}{x_1x_2}=\frac{^2F_{kl}}{x_2x_1}`$ and the equation 4a) requires that
$$ϵ^{mn}\frac{}{x_m}\left[L_{jk}^n\frac{^2W}{x_jx^l}+M_{jl}^n\frac{^2W}{x^k_j}+\frac{i}{2}H^{nj}\frac{^3W}{x^jx^kx^l}\right]=0$$
Moreover, if we restrict $`F_{ij}`$ to be symmetric under the exchange of indices and then identify $`L_{jk}^i=M_{jk}^i`$, equation 4b) becomes equal to 4c).
\- $`iv)`$ The equations of BOX 5 are satisfied if :
$$G_{ij}=J_{ij}=N_{ijk}=0$$
\- $`v)`$ Equation 6a) by itself, BOX 6, sets the only independent component of $`S_{ijkl}`$ to be constant; $`S_{1212}=\text{cte}`$. Then:
$$I_3=\theta _1^1\theta _1^2\theta _2^1\theta _2^2$$
is a constant of motion, an invariant.
### 4.2 Invariants in SuperLiouville models
We now apply the previous results to the computation of the supersymmetric extensions of the second invariant of Liouville models. In general they have the form:
$$I_2=I_2^{(B)}+I_2^{(F)},$$
where $`I_2^{(B)}`$ is the “body”, already present in the Liouville model, and $`I_2^{(F)}`$ is the “soul”-containing Grassman variables- of the second invariant in the SuperLiouville models. We find:
#### 4.2.1 SuperLiouville Models of Type I:
$`I_2^{(B)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\left(x^2\dot{x}^1x^1\dot{x}^2\right)^2\mathrm{\Omega }^2\dot{x}^2\dot{x}^2+\left(x^2{\displaystyle \frac{W}{x^1}}x^1{\displaystyle \frac{W}{x^2}}\right)^2\mathrm{\Omega }^2{\displaystyle \frac{W}{x^2}}{\displaystyle \frac{W}{x^2}}\right]`$
$`I_2^{(F)}`$ $`=`$ $`i(x^2\dot{x}^1x^1\dot{x}^2)\theta _a^1\theta _a^2+i\left(2x^1{\displaystyle \frac{^2W}{x^2x^2}}{\displaystyle \frac{W}{x^1}}x^2{\displaystyle \frac{^2W}{x^1x^2}}\right)\theta _1^2\theta _2^2+`$
$`+`$ $`i\left(x^1x^2{\displaystyle \frac{^2W}{x^2x^2}}+x^2{\displaystyle \frac{W}{x^1}}+x^2x^2{\displaystyle \frac{^2W}{x^1x^2}}\right)(\theta _1^1\theta _2^2+\theta _1^2\theta _2^1)+`$
$`+`$ $`i\left(x^2{\displaystyle \frac{W}{x^2}}x^1x^2{\displaystyle \frac{^2W}{x^1x^2}}+x^2x^2{\displaystyle \frac{^2W}{x^1x^1}}\right)\theta _1^1\theta _2^1`$
#### 4.2.2 SuperLiouville Models of Type II:
$`I_2^{(B)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(x^2\dot{x}^1x^1\dot{x}^2\right)^2+{\displaystyle \frac{1}{2}}\left(x^2{\displaystyle \frac{W}{x^1}}x^1{\displaystyle \frac{W}{x^2}}\right)^2`$
$`I_2^{(F)}`$ $`=`$ $`i\left(x^2\dot{x}^1x^1\dot{x}^2\right)\theta _a^1\theta _a^2+ix^2\left(x^2{\displaystyle \frac{W}{x^1x^1}}x^1{\displaystyle \frac{^2W}{x^1x^2}}{\displaystyle \frac{W}{x^2}}\right)\theta _1^1\theta _2^1+`$
$`+`$ $`ix^2\left(x^2{\displaystyle \frac{^2W}{x^1x^1}}+{\displaystyle \frac{W}{x^1}}x^1{\displaystyle \frac{^2W}{x^2x^2}}\right)(\theta _1^1\theta _2^2+\theta _1^2\theta _2^1)+`$
$`+`$ $`ix^1\left(x^1{\displaystyle \frac{^2W}{x^2x^2}}{\displaystyle \frac{W}{x^1}}x^2{\displaystyle \frac{^2W}{x^1x^2}}\right)\theta _1^2\theta _2^2`$
#### 4.2.3 SuperLiouville Models of Type III:
$`I_2^{(B)}`$ $`=`$ $`\left(x^1\dot{x}^2x^2\dot{x}^1\right)\dot{x}^2+\left(x^1{\displaystyle \frac{W}{x^2}}x^2{\displaystyle \frac{W}{x^1}}\right){\displaystyle \frac{W}{x^2}}`$
$`I_2^{(F)}`$ $`=`$ $`i\dot{x}^2\theta _a^1\theta _a^2ix^2{\displaystyle \frac{^2W}{x^1x^2}}\theta _1^1\theta _2^1ix^2{\displaystyle \frac{^2W}{x^2x^2}}(\theta _1^1\theta _2^2+\theta _1^2\theta _2^1)+`$
$`+i\left(2x^1{\displaystyle \frac{^2W}{x^2x^2}}{\displaystyle \frac{W}{x^1}}x^2{\displaystyle \frac{^2W}{x^1x^2}}\right)\theta _1^2\theta _2^2`$
#### 4.2.4 SuperLiouville Models of Type IV:
$$I_2^{(B)}=\frac{1}{2}\dot{x}^1\dot{x}^1+\frac{1}{2}\frac{W}{x^1}\frac{W}{x^1}I_2^{(F)}=i\frac{^2W}{x^1x^1}\theta _1^1\theta _2^1$$
Finally, we briefly comment on the geometrical and physical meaning of the second invariant. Usually, it is related to transformation that is termed as a hidden symmetry. We see that by introducing the generalized momenta $`\mathrm{\Pi }_j=\dot{x}_j+i\frac{W}{x^j}`$, the second invariant of the Type I model is:
$$I_2^{(B)}=\frac{1}{2}\left[\left|x^2\mathrm{\Pi }_1x^1\mathrm{\Pi }_2\right|^2\mathrm{\Omega }^2\left|\mathrm{\Pi }_2\right|^2\right],$$
which is no more than the modulus of the generalized angular momentum to the square added to $`\mathrm{\Omega }^2`$ times the square of $`|\mathrm{\Pi }_2|`$. Similar considerations are easily applied to the second invariant of the other Types. A generalized momentum such as $`\mathrm{\Pi }_j`$ can be obtained if one adds the complex topological piece:
$$L_T^{(B)}=i\dot{x}^j\frac{W}{x^j}$$
to the bosonic Lagrangian.
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# Linear Collider Signal of a Wino LSP in Anomaly Mediated Scenarios
## Abstract:
Selectron (smuon) pair-production in a next generation Linear Collider, yielding a fast electron (muon) trigger, a visible heavily ionizing track and/or a resolved soft pion impact parameter and overall $`E_T/`$, is shown to provide a smoking gun signature for Anomaly Mediated Supersymmetry Breaking models with a neutral Wino as the Lightest Supersymmetric Particle, nearly mass-degenerate with the lighter chargino.
Supersymmetry Breaking, Beyond Standard Model preprint: TIFR/TH/00-18
hep-ph/0004127
Understanding how supersymmetry breaks in the real world from a deeper, more fundamental, standpoint is a challenge in high energy physics today. An interesting recent idea in this direction has been that of Anomaly Mediated Supersymmetry Breaking (AMSB) , based on which a whole class of supersymmetric models \[ReferencesReferences\] have emerged. A crucial signal in a high energy Linear Collider, namely $`e^+e^{}e^\pm (\mu ^\pm )+soft\pi ^{}+E_T/`$, to test such a scenario, is proposed in this Letter.
AMSB models are strongly motivated by String Theory which is defined in a higher dimensional spacetime and is valid at a very high energy scale. It is quite natural from that point of view to expect a low energy description of the physical world in four dimensions to inherit some of the features of the higher dimensional theory. This is indeed the case with AMSB scenarios. AMSB occurs when in such a higher dimension, one has a supergravity theory defined on two separated parallel 3-branes ($`3+1`$ dimensional subspaces) in a way that the Standard Model (SM) particles are localized on one of these while the supersymmetry breaking sector is localized on the other. There are no tree-level couplings between these two branes and thus the supersymmetry breaking sector is truly hidden. Gravity propagates in the bulk and the breakdown of supersymmetry is communicated from the hidden to the visible sector through the loop-induced super-Weyl anomaly. In the absence of tree-level interactions between the two 3-branes, this is the dominant contribution to the soft supersymmetry breaking parameters determining the masses of various superparticles. In the more commonly used Gravity Mediated Supersymmetry Breaking scenario , supergravity interactions directly communicate supersymmery breaking between the hidden and observable sectors at the tree level, so that loop-induced contributions from the super-Weyl anomaly, though present, are subdominant. A characteristic feature of AMSB models is that the stable LSP or Lightest Supersymmetric Particle ($`\stackrel{~}{\chi }_1^0`$) is almost exclusively a neutral Wino which is nearly mass-degenerate with the lighter chargino ($`\stackrel{~}{\chi }_1^\pm `$), also predominantly a Wino. Models with only AMSB have the problem of tachyonic sleptons; however, modified versions exist in which the sleptons have physical masses.
Though the quest for supersymmetry has been a major preoccupation of collider experimenters and phenomenologists alike, most of the searches and simulation studies so far have been based on the assumption of the LSP being predominantly a Bino; i.e. the superpartner of the $`U(1)_Y`$ gauge boson. Various produced superparticles are expected to decay into the LSP accompanied by other particles of the Standard Model (SM). The LSP escapes detection carrying off missing transverse energy or $`E_T/`$ which becomes the classic signature. There have, however, been a few papers \[ReferencesReferences\] of late which have considered detection possibilities for scenarios in which a largely Wino LSP occurs. Our investigation belongs to this genre. We consider the pair production of left-selectrons (smuons) in $`e^+e^{}`$ interactions, followed by their decays<sup>1</sup><sup>1</sup>1These decay channels need not make all of the selectron (smuon) width. There are regions of parameter space where $`\stackrel{~}{\chi }_2^{\pm ,0}`$ can be lighter than the selectron (smuon), but decays like $`\stackrel{~}{e}(\stackrel{~}{\mu })e(\mu )\stackrel{~}{\chi }_2^0`$, $`\nu \stackrel{~}{\chi }_2^\pm `$, which open up there, are only a few percent of the branching ratio. Even so, we take these into account in our calculations. $`\stackrel{~}{e}(\stackrel{~}{\mu })e(\mu )+\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{e}(\stackrel{~}{\mu })\nu +\stackrel{~}{\chi }_1^\pm `$; $`\stackrel{~}{\chi }_1^\pm `$ will further decay into $`\stackrel{~}{\chi }_1^0+\pi ^\pm `$. Finally, there will be a fast $`e^\pm (\mu ^\pm )`$ trigger, a displaced vertex which can be inferred from the impact parameter of a visible soft $`\pi ^\pm `$ and/or a heavily ionizing track with high momentum (i.e. nearly straight in the magnetic field) and large $`E_T/`$. Similar considerations just with selectrons can also be made for $`e^{}e^{}`$ collision.
The chargino and the neutralino masses, in any version of the Minimal Supersymmetric Standard Model (MSSM) , are controlled by the following supersymmetry parameters at the weak scale: the Bino mass $`M_1`$, the Wino mass $`M_2`$, the Higgsino mass parameter $`\mu `$ and the ratio $`\mathrm{tan}\beta `$ of the two Higgs VEVs. The situation, with the LSP ($`\stackrel{~}{\chi }_1^0`$) being largely the neutral Wino, obtains when one has
$$|M_2|<|M_1||\mu |.$$
(1)
One should emphasize that, within a MSSM framework, the mass-hierarchy (1) is very characteristic of AMSB models. For instance, in such models, after taking into account next-to-leading order corrections to the gaugino mass parameters, one finds that $`M_1:M_22.8:1`$ as contrasted with $`M_1:M_21:2`$ in gauge or usual supergravity mediated supersymmetry breaking models with gaugino masses unified at the grand unifying scale. The next-to-lightest superparticle in the AMSB case is the lighter chargino ($`\stackrel{~}{\chi }_1^\pm `$) which is almost exclusively a Wino. Then the masses of $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_1^\pm `$ are verly close and the small mass-splitting $`\mathrm{\Delta }M`$ has the form :
$`\mathrm{\Delta }M={\displaystyle \frac{M_W^4\mathrm{tan}^2\theta _W}{(M_1M_2)\mu ^2}}\mathrm{sin}^22\beta \left[1+𝒪({\displaystyle \frac{M_2}{\mu }},{\displaystyle \frac{M_W^2}{\mu M_1}})\right]`$
$`+{\displaystyle \frac{\alpha M_2}{\pi \mathrm{sin}^2\theta _W}}\left[f\left({\displaystyle \frac{M_W^2}{M_2^2}}\right)\mathrm{cos}^2\theta _Wf\left({\displaystyle \frac{M_Z^2}{M_2^2}}\right)\right],`$ (2)
with
$`f(x)`$ $``$ $`{\displaystyle \frac{x}{4}}+{\displaystyle \frac{x^2}{8}}\mathrm{ln}(x)+{\displaystyle \frac{1}{2}}(1+{\displaystyle \frac{x}{2}})\sqrt{4xx^2}[\mathrm{tan}^1\left({\displaystyle \frac{2x}{\sqrt{4xx^2}}}\right)`$
$`\mathrm{tan}^1\left({\displaystyle \frac{x}{\sqrt{4xx^2}}}\right)].`$
The second term in the RHS of Eq. 2 is the one-loop contribution which is dominated by gauge boson loops.
The mass-splitting $`\mathrm{\Delta }M`$ of Eq. 2 has been investigated numerically in various region of the parameter space consistent with Eq. 1. The general conclusion is that
$$165\mathrm{MeV}\stackrel{<}{}\mathrm{\Delta }M\stackrel{<}{}1\mathrm{GeV},$$
(3)
with $`lim_\mu \mathrm{}`$ $`\mathrm{\Delta }M`$ being $`165`$ MeV. On the other hand, if a radiative electroweak (EW) symmetry breakdown is sought to be implemented in the AMSB scenario, the ratio $`|\mu /M_2|`$ has to be approximately between 3 and 6. Given the LEP lower limit of 56 GeV on the mass of the lighter chargino for nearly degenerate $`\stackrel{~}{\chi }_1^\pm ,\stackrel{~}{\chi }_1^0`$ (in the anomaly-mediated Wino LSP scenario), it then follows that the upper limit on $`\mathrm{\Delta }M`$ cannot be much in excess of 800 $`\mathrm{MeV}`$. In that case we can take
$$165\mathrm{MeV}\stackrel{<}{}\mathrm{\Delta }M\stackrel{<}{}800\mathrm{MeV}.$$
(4)
Eq. 4 means that the decay $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0+\pi ^\pm `$ is kinematically allowed. The corresponding branching ratio is found to vary in the range $`93\%96\%`$, the balance being largely due to the decay modes $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0+e+\nu _e`$, $`\stackrel{~}{\chi }_1^0+\mu +\nu _\mu `$. The resulting soft pion with a sub-GeV energy may be detectable, in which case its impact parameter may allow one to infer a displaced vertex. On the other hand, the $`\stackrel{~}{\chi }_1^\pm `$ may have a long enough decay length to show a high momentum heavily ionizing track which stops in some of the layers in the vertex detector. The experimental issues concerning methods of observing this decay have been discussed in the third paper of Ref. and in Refs. . If the decay length<sup>2</sup><sup>2</sup>2Here $`c\tau =c\mathrm{}p_{\stackrel{~}{\chi }^\pm }(M_{\stackrel{~}{\chi }^\pm }\mathrm{\Gamma }_{\stackrel{~}{\chi }^\pm })^1`$ with $`p_{\stackrel{~}{\chi }^\pm }`$, $`M_{\stackrel{~}{\chi }^\pm }`$ and $`\mathrm{\Gamma }_{\stackrel{~}{\chi }^\pm }`$ respectively being the momentum, mass and width (all in GeV) of the chargino. $`c\tau `$ of $`\stackrel{~}{\chi }_1^\pm `$ is greater than $`3`$ cm., it could be observable<sup>3</sup><sup>3</sup>3A CCD or APS vertex detector of radius 2.5 cm and a beam pipe of radius 2 cm, have been proposed for TESLA. The chargino track should be identifiable if it covers several layers and also ends in the vertex detector. though the $`\pi ^\pm `$ may be too soft to be detected. Contrariwise, if $`c\tau <3`$ cm, the track may not be observable but the soft charged pion is likely to be visible with its impact parameter $`b`$ resolved. Thus the event, proposed by us, can be triggered by the fast charged lepton emanating from the decay of one of the sleptons while it can be identified uniquely in terms of the displaced vertex determined by the heavily ionizing charged track, which should be nearly straight in the magnetic field because of the high momentum, and/or the impact parameter $`b`$ of the soft pion coming from the two-step decay of the other slepton.
In the anomaly mediated case, gaugino masses are proportional to the coefficients of the one-loop beta functions of the corresponding gauge couplings (generically denoted as $`g`$), while scalar masses are determined in terms of anomalous dimensions and beta functions of both gauge and Yukawa couplings (generically denoted as $`y`$). The expressions for the anomaly induced contributions to the soft masses are
$$M_\lambda =\frac{\beta _g}{g}m_{3/2},$$
(5)
$$m_{\stackrel{~}{f}}^2=\frac{1}{4}\left(\frac{\gamma }{g}\beta _g+\frac{\gamma }{y}\beta _y\right)m_{3/2}^2,$$
(6)
$$A_y=\frac{\beta _y}{y}m_{3/2},$$
(7)
where gaugino masses are denoted by $`M_\lambda `$, scalar masses are given the generic symbol $`m_{\stackrel{~}{f}}`$, $`m_{3/2}`$ is the mass of the gravitino which here is quite heavy ($``$ tens of TeV) and $`A_y`$ are the trilinear soft parameters defined with the convention of the third paper of Ref. . The renormalization group beta and gamma functions are defined as $`\gamma (g,y)d\mathrm{ln}Z/dt`$, $`\beta _g(g,y)dg/dt`$ and $`\beta _y(g,y)dy/dt`$, $`t`$ being the logarithmic scale variable.
The most striking feature of this AMSB scenario is the invariance of the expressions for soft SUSY breaking mass parameters Eqs. (57) under renormalization group (RG) evolution. Thus, these parameters can be evaluated at any scale with the appropriate values of the gauge couplings at that particular scale. However, the mass squares of the sleptons, calculated in this way, turn out to be negative. These tachyonic sleptons constitute a major problem of this scenario. The most simple and economical way by which these slepton mass-squares can be made positive is to add a common $`m_0^2`$ to all scalars and this is what we consider. However, our signal is also present for models where this positive term is nonuniversal and arises from the D-term of a broken $`U(1)`$ gauge symmetry. Of course, the addition of any such term destroys the RG invariance of Eq. 6. Then, in order to get the correct values of the mass-squares of the scalars at the EW scale, one must take into account the RG evolution of these soft masses from a very high scale. In our calculations we have taken this to be the unification scale ($`1.5to2.0\times 10^{16}\mathrm{GeV}`$) where all the three gauge couplings meet and the evolution of these couplings reproduces the measured values at the EW scale with $`\alpha _s0.118`$. The evolution of gauge and Yukawa couplings has been determined by two-loop RG equations. The detailed expressions for scalar and gaugino masses as well as the trilinear A-parameters are given in Refs. References and References. The Higgsino mass parameter $`\mu `$ has been computed using complete one-loop correction terms of the effective potential at the scale $`Q`$ in such a way that it reproduces the correct pattern of EW symmetry breaking with $`Q`$ chosen to be the geometric mean of the $`t`$-squark masses $`\sqrt{m_{\stackrel{~}{t}_1}m_{\stackrel{~}{t}_2}}`$. The supersymmetric correction to the mass of the bottom quark (sizable for large $`\mathrm{tan}\beta `$) has also been computed to one-loop. We have, moreover, accounted for the constraints coming from charge and color conservation as well as from the experimental lower limits on various sparticle masses including $`m_{\stackrel{~}{\chi }_1^\pm }>56`$ GeV and also from the requirement of the stau not being the LSP.
We have determined the slepton and chargino/neutralino sector of the MSSM mass spectrum completely in terms of $`m_{3/2}`$, $`m_0`$, $`\mathrm{tan}\beta `$ (the ratio of the two Higgs vacuum expectation values) and the sign of $`\mu `$. We have checked that our results agree with those of previous authors for $`\mathrm{tan}\beta =3`$ with $`\mu <0`$ and $`\mu >0`$ as well as for $`\mathrm{tan}\beta =30`$ with $`\mu <0`$ and $`\mu >0`$. The LSP $`\stackrel{~}{\chi }_1^0`$ and the lighter chargino $`\stackrel{~}{\chi }_1^\pm `$ are found to be very nearly degenerate, as suggested by Eq. 2. Indeed, we find $`\mathrm{\Delta }M`$ not only to obey the inequality (3); but also to be a decreasing function of $`m_{3/2}`$, asymptotically reaching the lower bound of (3) when the latter gets very high. This function is quite insensitive to the value of $`m_0`$. The left and right selectron masses are also found to be almost degenerate. The tiny mass-difference between the latter comes mainly from one-loop corrections at the electroweak scale since the anomaly induced as well as $`D`$-term contributions are negligible in comparison. This, again, is a distinguishing feature of the AMSB scenario which is based on the assumption of a universal contribution to the mass-squared for scalars added to make the sleptons non-tachyonic. An important point is that the region of parameter space where the masses of the selectrons (smuons) do not lie between those of $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\chi }_2^{\pm ,0}`$, the latter being the higher chargino/neutralino, is not small, though this does not affect<sup>4</sup><sup>4</sup>4See, footnote 1. our analysis. The other important aspect of the superparticle spectrum in such an AMSB scenario is that the squarks are significantly heavier (typically by at least a factor of four) than the sleptons, the squark masses being pushed up by the QCD coupling. This means that sleptons should be easier to discover in such models. This is why we have chosen to study slepton pair-production in a linear collider. Of course, one could also directly study the pair-production of charginos $`\stackrel{~}{\chi }_1^\pm `$, each decaying into $`\stackrel{~}{\chi }_1^0`$ and a soft pion. However, one would then need to have an additional hard initial-state-radiated (ISR) photon to act as a trigger. The event rate there would be significantly less than that of slepton pair-production on account of the former process being radiative.
We have calculated the left-selectron (these would be mass eigenstates because of negligible left-right mixing) pair production cross section at an $`e^+e^{}`$ CM energy of $`1`$ TeV for two values of $`\mathrm{tan}\beta `$, namely, $`3`$ and $`30`$, for $`\mu <0`$ and $`\mu >0`$. We have then folded into it the branching fractions for the decays mentioned in the first paragraph. The selection cuts that have been used on the decay products are as follows : (1) the transverse momentum of the lepton $`p_T^{\mathrm{}}>5`$ GeV, (2) the pseudorapidities of the lepton and the pion $`|\eta |<2.5`$, (3) the electron-pion isolation variable $`\mathrm{\Delta }R=\sqrt{(\mathrm{\Delta }\eta )^2+(\mathrm{\Delta }\varphi )^2}>0.4`$, (4) the missing transverse energy $`E_T/>20`$ GeV and $`(5)`$ $`p_T^\pi >200`$ MeV for a detectable soft pion (N.B. the total momentum of the pion is in the range of hundreds of MeV). Contour plots in the $`m_0m_{3/2}`$ plane for various values of cross-sections (in fb) are shown in Fig. 1. The shaded regions are excluded by the constraints mentioned earlier; in addition, the selectron mass has been required not to exceed $`500`$ GeV which is the kinematic limit for observability in a $`1`$ TeV Linear Collider. The allowed region is somewhat smaller for large $`\mathrm{tan}\beta `$ because of stronger left-right mixing in the stau sector. We see that quite interesting regions in the $`m_0m_{3/2}`$ plane are covered for cross sections ranging from $`10fb`$ to $`125fb`$. Our signal should thus generate $`𝒪(10^4)`$ events for an integrated luminosity of $`500(fb)^1`$. These calculations have been done with projected TESLA parameters in mind . For a scaled down linear collider, e.g. with a CM energy of $`500`$ GeV and an integrated luminosity of $`50(fb)^1`$, we would expect $`𝒪(10^3)`$ events.
We have also plotted the decay length $`c\tau `$ distribution of the chargino track in Fig. 2 with the same selection cuts as used in Fig. 1; in addition, we have chosen characteristic sample values of $`m_0=230`$ GeV and $`m_{3/2}=43`$ TeV, $`\mathrm{tan}\beta =3`$ and $`\mu <0`$ corresponding to $`\mathrm{\Delta }M=182.8`$ MeV. We observe a plateau in the $`c\tau `$ distribution in the range $`8.5`$ to $`9.9`$ cm which can cover several layers in the vertex detector. Thus there is a reasonable chance of a direct observation of the chargino track. The transverse momentum ($`p_T^\pi `$) and the impact parameter $`b`$ distributions of the soft pion are plotted in Figs. 3a and 3b respectively with the same input parameters and selection cuts as for Fig. 2. The $`b`$-distribution extends till about $`9.9`$ cm and peaks at around $`b=8.5`$ cm. It is clear from the $`p_T^\pi `$ distribution that the minimum $`p_T^\pi `$ cut of $`200`$ MeV still leaves a substantial part of the allowed phase space for study. For such values of $`p_T^\pi `$, the $`3\sigma `$ impact parameter resolutions are typically $`𝒪(10^1)`$ cm. Of course, we have chosen a particularly favorable region of the allowed MSSM parameter space. The numbers are not always so good in other regions. We have nonetheless checked that $`b`$ is always significantly above the impact parameter resolution value. Hence the prospects of resolving the displaced vertex by measuring the soft pion impact parameter here are quite high. Let us comment finally that, if selectrons are replaced by smuons (with a fast muon used as a trigger), event rates are reduced typically by a factor of five on account of s-channel suppression.
An alternative MSSM scenario of nearly degenerate $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_1^\pm `$ (and $`\stackrel{~}{\chi }_2^0`$ as well) can arise when $`|\mu ||M_{1,2}|`$. In such a case a mass-difference $`\mathrm{\Delta }M(\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0)\stackrel{<}{}\mathrm{\hspace{0.25em}1}`$ GeV can be obtained with $`m_{\stackrel{~}{\chi }_1^\pm }>51`$ GeV by setting $`|M_{1,2}|\stackrel{>}{}\mathrm{\hspace{0.25em}5}`$ TeV and $`|\mu |\stackrel{>}{}M_Z/2`$. Though this is a rather unnatural scenario and quite difficult to obtain in a phenomenologically viable model, we can ask whether our signal can be mimicked here. The answer is no. The two-body decays of selectrons, relevant for us, are highly suppressed in this other scenario on account of the factor $`m_e/M_W`$ in the concerned couplings. The latter arises because $`\stackrel{~}{\chi }_1^\pm `$, $`\stackrel{~}{\chi }_{1,2}^0`$ are all almost exclusively higgsinos here. So selectrons primarily have three-body decays $`\stackrel{~}{e}\nu _eW\stackrel{~}{\chi }_{1,2}^0`$, $`eZ\stackrel{~}{\chi }_{1,2}^0`$ mediated by virtual heavier charginos/neutralinos $`(\stackrel{~}{\chi }_2^\pm /\stackrel{~}{\chi }_2^0)`$, which are gauginos, with finals states dominated by jets. One can easily estimate the ratio of the partial widths of left selectron decays into two-body and three-body channels to be $`𝒪(10^4)`$ in this scenario demonstrating that the desired two-body decays would be unobservable. Therefore, unlike the soft pion plus hard ISR photon signal studied in Ref. , our final state of a fast electron (muon) and a soft pion distinguishes AMSB models from the light higgsino scenario. We would like to highlight this new result which has emerged from the present work.
Let us also discuss the question of background to our signal. The signal can be classified into two categories. There is one in which we see a heavily ionizing nearly straight charged track ending with a soft pion with large impact parameter and $`E_T/`$, the signal being triggered with a fast electron or a muon. In the other case, while the other aspects remain the same, one may not see the heavily ionizing charged track but the impact parameter of the soft pion can be resolved and measured to be large. In the first case the heavily ionizing charged track is due to the passage of a massive chargino with a very large momentum. Due to this reason the charged tarck will be nearly straight in the presence of the magnetic field. One cannot imagine a similar situation in the SM with such a nearly straight heavily ionizing charged track due to a very massive particle. An ionized charged track can possibly arise from the flight of a low energy charged pion, kaon or proton but it will curl significantly in the magnetic field. Another distinguishing feature of the charged track in our signal is that it will be terminated after a few layers in the vertex detector and there will be a soft pion at the end. In the second case, where the ionizing track is unseen, possible SM backgrounds can come from the following processes: $`e^++e^{}\tau ^++\tau ^{}`$ and $`e^++e^{}W^++W^{}`$. In the case of $`e^++e^{}\tau ^++\tau ^{}`$, one $`\tau `$ can have the three body decay $`\tau e\nu _e\nu _\tau `$ or $`\mu \nu _\mu \nu _\tau `$ and the other $`\tau `$ can go via the two body channel $`\tau \pi +\nu _\tau `$. Thus we can have a final state of the type $`e(\mu )+\pi +E_T/`$. Since we are considering an $`(e^+e^{})`$ CM energy of 1 TeV, and the pion comes from a sequence of two-body production and decay, it will have a fixed high momentum much in excess of 1 GeV. This will clearly separate this type of background from our signal since in our case the resulting pion is very soft with a momentum in the range of hundreds of MeV. In the case of $`e^++e^{}W^++W^{}`$ a similar argument follows. Here one $`W`$ can go to $`e(\mu )+\nu _e(\nu _\mu )`$ and the other one can go to $`\tau +\nu _\tau `$. The $`\tau `$ can subsequently go to one $`\pi `$ and a $`\nu _\tau `$, thereby producing the final state $`e(\mu )+\pi +E_T/`$. As we have discussed just now, the resulting pion will have a very large momentum and again one can clearly separate the background from the signal.
In conclusion, we claim to have pinpointed a fast electron (muon) trigger, overall $`E_T/>20`$ GeV and a displaced vertex emitting a soft pion in the final state configuration as a distinct and unique linear collider signal of the AMSB scenario with a Wino LSP. A more detailed discussion of this as well as other linear collider signals of AMSB models will be given elsewhere.
This work came out of a study-project on Anomaly Mediated Supersymmetry Breaking at the Workshop on High Energy Physics Phenomenology WHEPP-6 (Chennai, January, 2000). We thank the organizers as well as the other members of the project namely D. Choudhury, S. King, A. Kundu, B. Mukhopadhyaya, S. Raychaudhuri and K. Sridhar for many fruitful discussions. We also thank U. Chattopadhyay for the use of his codes and M. Maity and N. K. Mondal for discussions of experimental issues.
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# SUE: A Special Purpose Computer for Spin Glass Models
## 1 Introduction
Two approaches have become popular in the field of computer design for scientific calculations: special or general purpose computers. Lattice Monte Carlo in Quantum Field Theory and Statistical Mechanics requires large computational power in relatively general purpose computers and the processing can often be parallelized. Various groups have developed their own parallel machines for those simulations . Those general purpose computers require continuous technological upgrading and investment to obtain competitive results. On the other hand, special purpose computers can approach very specific problems, achieving better performance than general computers.
The emergence in the market of Complex Programmable Logic Devices (CPLD) makes it possible to design dedicated machines with low cost and high performance. In this paper we describe a CPLD-based machine, dedicated to three-dimensional spin glass models with variables belonging to $`Z_2`$ and couplings to first neighbours, and report on the reliability tests which have been carried out.
Our machine is called SUE, for Spin Updating Engine, because its task is to generate sets of updated spin configurations in the Monte Carlo simulation. In a previous work , we presented the prototype for the two-dimensional model, and introduced the first ideas about the final version. After checking that the 2d version worked properly, we have designed, constructed and tested the 3d version, which differs in some aspects from the 2d version as we will see below. The performance of the 3d machine is improved due to the fact that it can run more than a single model: the lattice size or the action of the physical model can be easily changed using the on-board reprogrammability of the CPLDs. A device devoted to generate a 32-bit random number has been developed and included in every SUE board. This device (described below) enables SUE to operate with both canonical and microcanonical algorithms.
At present, spin glass models are a progressing area of Statistical Mechanics. They are related to neural networks, spin models, some High $`T_c`$ superconductivity models, etc. There is large activity in the 3d models because of the uncertainty in the vacuum structure at low temperature. Monte Carlo simulations of spin glass systems have been used to study the phase transition, the ultrametric structure and the dynamics out of equilibrium . Only sizes up to $`L=16`$ have been simulated , due to the slow dynamics of the systems and the strong slowing-down as the size grows. Yet, those simulations requiring very simple calculations, they are easily implementable in a dedicated machine. That way the computational power needed to obtain results in larger lattices is obtained.
A standard way of studying spin glasses is the use of independent lattices with the same quenched couplings, called replicas. The overlap between two replicas acts as the order parameter in that model. A great improvement on the usual Monte Carlo scheme is the parallel tempering method . The basic idea is to move in temperature space: the system changes its temperature, goes up to the paramagnetic phase and eventually goes back to lower temperature. With high probability in its motion through temperature the system will visit new local minima. That scheme has been implemented in SUE: Replicas at different temperatures are simulated, and systems running at adjacent levels can be swapped according to an appropriate probability distribution.
An essential tool for the analysis of results is Finite Size Scaling , which requires the use of different volumes. In that sense, SUE is capable of working with different lattice sizes by reprogramming its CPLD devices.
The main differences with respect to the 2d prototype presented in are:
* Larger and faster devices.
* Multi-Layer instead of Double-Layer Printed Circuits.
* On-board reprogrammability.
* Dedicated device for 32-bit random number generation.
* Demon and Heat Bath algorithm support.
* Parallel Tempering implementation.
* Driver and Software development for easy (transparent) use.
At present we have built 12 boards and tested them using the Demon and Heat Bath algorithms in different lattice sizes. Each board simulates 8 lattices, updating 8 spins every $`20.8ns`$ cycle. The update speed of a single board is therefore $`2.6ns/spin`$. The cost of each board is 2400 Euros (500 for PCB and mounting, and 1900 for components.)
The summary of this paper is as follows: We start introducing the physical model in the next section. In section 3 We describe the electronic architecture of SUE, design considerations and software support. The development process is outlined in section 4. Last section is devoted to discuss the performance.
## 2 The Physical Model
We want to simulate the 3D Edwards-Anderson model with first neighbour couplings (see for a detailed description of the model). The action of this model is given by
$$E=\underset{i,j}{}\sigma _i\sigma _jJ_{ij},$$
(1)
where the value of the Ising spins $`\sigma `$ can be $`1`$ or $`1`$, and the couplings $`J_{ij}`$ are random variables taking the values $`\pm 1`$ with equal probability. For a fixed set of couplings $`\{J_{ij}\}`$ the partition function is
$$Z(\beta ,\{J_{ij}\})=\underset{\{\sigma \}}{}\mathrm{exp}\beta E(\{J_{ij}\},\{\sigma \}).$$
(2)
We study the existence of phase transitions using as order parameter the overlap between two independent systems (replicas) with the same set of couplings $`J_{ij}`$. We should finally average over different realizations of the disorder ($`J_{ij}`$) to obtain physical results about the system.
To calculate (2) we must sum over $`2^V`$ possible configurations, where $`V`$ is the volume of the lattice, which is a very large number for any computer. The standard way to compute the partition function is to run an algorithm that selects only a representative set of configurations. There are different appropiate algorithms, see for instance the chapter by Sokal in . For pure spin systems (all $`J_{ij}`$ equal to $`1`$) some cluster algorithms are very efficient, but for a general spin glass model only local algorithms achieve good efficiency.
Typically, we must run the algorithm and generate millions of different representative configurations in order to obtain accurate results. The autocorrelation time $`\tau `$ is a measure of the correlation between configurations: a run of length $`n`$ provides only $`n/\tau `$ effectively independent samples. Near the critical point, $`\tau `$ diverges as $`\tau L^z`$, where $`L`$ is the size of the system and $`z`$, the dynamical critical exponent, has been found to be around $`6`$ (while in the pure Ising model it is close to $`2`$). This strong slowing-down, due to the existence of many pure and metastable states, and the absence of non local algorithms, makes this problem really hard from a computational point of view.
Two different updating algorithms have been implemented in our design, one microcanonical (Demon) and one canonical (Heat Bath).
The Demon algorithm keeps the sum of the lattice energy and a demon energy constant. In order to generate the representative set of samples, we start from a spin configuration with an action $`S`$ and a demon energy equal to zero. Now we use the algorithm to change the spins to generate new configurations (one for every $`V`$ updates). The update of a spin is as follows: if the flip lowers the spin energy, the demon takes that energy and the flip is accepted. On the other hand, if the flip increases the spin energy, the change is only made if the demon energy is sufficient to transfer that energy to the system. The conservation of the total energy (lattice plus demon), has been useful in the programming/test stage, allowing fast tests of proper function.
In the Heat Bath algorithm , the new spin value for each site $`\sigma _i`$ is independent from the old one, and its probability distribution is that of a single Ising spin $`\sigma _i`$ in the effective magnetic field produced by the fixed neighbouring spins $`\sigma _j`$:
$$P(\sigma _i\{\sigma _j\}_{ji})=\frac{\mathrm{exp}(\beta \sigma _i_jJ_{ij}\sigma _j)}{\mathrm{exp}(\beta _jJ_{ij}\sigma _j)+\mathrm{exp}(\beta _jJ_{ij}\sigma _j)}.$$
(3)
The drawback of this canonical algorithm is the necessity of a random number to decide the acceptance of the new spin.
The algorithms being local, to update a spin only the nearest neighbours are needed. Because of simplicity in the electronic design, we use helicoidal boundary conditions. Let us consider a lattice of side $`L`$ and volume $`V=L^3`$ with sites labelled in the standard way: the site $`[x,y,z]`$ (with $`0x,y,zL1`$) gets the index $`n=x+y\times L+z\times L^2`$. We will call $`x_+`$ ($`y_+`$, $`z_+`$) the neighbour in the positive direction along the $`x`$ ($`y`$, $`z`$) axis. With our helicoidal boundary conditions the neighbours of the site $`n`$ are simply:
$$\begin{array}{ccc}x_+& =& (n+1)modV\\ y_+& =& (n+L)modV\\ z_+& =& (n+L^2)modV\end{array}$$
(4)
We define in an analogous manner the remaining neighbours $`x_{}`$, $`y_{}`$ and $`z_{}`$.
## 3 Operation and General Structure of SUE
The SUE machine is connected to a Host Computer (HC) running under Linux. SUE performs the update of the configurations, but the measurements and analysis are made by the HC. SUE is set up with initial spin configurations, couplings and several simulation parameters. Then SUE is started and simulation begins. After a certain number of iterations, SUE is stopped to download the configuration to the HC and SUE keeps the updating process. In this sense, SUE and the HC work in parallel: while SUE is updating the system the HC processes the previously read configurations.
Fig. 1 shows a simple diagram of the whole machine which consists of the HC and $`n`$ SUE boards (the figure is for $`n=8`$, but the final system consists of 12 processing modules). They are connected to the HC through a PCI Data Acquisition Card. Every processing module contains the hardware to store and update eight lattices in parallel. Note that there are two degrees of parallelism: inside the processing module and between the modules.
Every clock cycle, the random number generator device included in each board provides a pseudo-random number which is shared for the update of the eight lattices, so the replicas (systems with the same couplings $`J_{ij}`$) must be simulated in different boards. We can then think of each pair of boards as a unit, allowing us to simulate eight pairs of replicas (corresponding to eight realizations of disorder $`J_{ij}`$). Periodically, the configurations are read and the relevant measurements carried out and stored.
Parallel tempering requires the simulation of pairs of replicas with the same couplings at different values of $`\beta `$. With 12 boards we could then simulate replicas corresponding to eight sets of couplings at 6 values of $`\beta `$ at once. Parallel tempering requires more temperature levels, so each time the configuration is read the $`\beta `$ is changed (the corresponding probability table is loaded), and the configuration to be updated is loaded onto the board (in the meantime it was stored in the HC). Different temperatures are then sequentially simulated.
The HC controls this mechanism, and is responsible for deciding whether the configurations being simulated at adjacent temperatures are interchanged. Given the configurations $`X`$ at temperature $`\beta `$ and $`X^{}`$ at temperature $`\beta `$’, we compute
$$\mathrm{\Delta }=(\beta ^{}\beta )(E(X)E(X^{}))$$
(5)
and use a Metropolis like test: if $`\mathrm{\Delta }<0`$ we accept the change, otherwise we swap the configurations with probability $`\mathrm{exp}(\mathrm{\Delta })`$.
By processing 8 spins in parallel on 12 modules (96 spins in total) within one clock cycle (clock period of 48 MHz), we obtain an update speed of 217 ps/spin. The time spent reading, writing and processing (Meassurement and Paralell Tempering) the configurations is around 4% of the computation time in the smallest simulable lattice ($`L=20`$), and decreases steeply with size (it is less than 1% of the total time for $`L=30`$).
Let us describe the main characteristics of a SUE board. Devices used are listed in table 1, apart from passive components (resistors, diodes, capacitors, leds, etc.). The main electronic devices are the Altera 10K CPLDs .
The photograph of one of the boards can be seen in fig. 2. It contains four devices FLEX 10K30 responsible for the core of the Monte Carlo simulation (UPDATE area on the figure). That four devices have the same electronic logic inside, which is prepared to update two lattices in parallel.
On the right of these chips are the static memory devices (SRAM) which store the couplings of the lattices (J-MEMORY). On the left, each UPDATE device has two SRAM devices which store the spin variables (SPIN MEMORY). Latch devices are used as tristate devices to manage the polarity of the data buses at high frequency.
The RNG device is a FLEX 10K50 where a random number generator is programmed, allowing the use of canonic simulations. Addressing of the memories and sincronization between the devices are the main tasks of the fifth FLEX 10K30 (ADDRESS). The coupling memories are addressed through latch devices to avoid fan-out problems.
External communication is provided by three EPM7032 chips placed near the 68-pin connector. One of them controls the board when the on-board programmable devices are not yet programmed. It responds to basic commands sent from the HC, allowing to select and program the board. The programmed logic establishes 4 control lines in each direction allowing communication between the HC an the ADDRESS device, and a 32-bit data bus common to the HC and the UPDATE and RNG devices. The two lower bits in this bus reach ADDRESS too, and act as extra control lines when needed.
The clock signal is distributed to all the synchronous devices in the board through Cypress 2308-4 PLL devices. In the upper right corner, a set of leds permits us to visualize the state of the machine. The connection to the HC is made through a 68 pin (SCSI-2 type) connector. The SUE boards can share the same bus for an easy management from the HC.
Once the general architecture of a SUE board has been outlined, the internal details are explained more deeply in the next subsections.
### 3.1 Updating Logic
Four Altera 10K30 devices are responsible for the update. To each one of those devices lines are assigned to access the spin and coupling memories and the 32-bit bus through which the random number is provided. That bus is used also to write and read the memories, the demon energy or the probability table from the HC.
In order to obtain an updated spin every clock cycle we have designed a pipeline structure that performs the algorithm step by step: A state machine runs over a 10 states cycle during the simulation, one spin being put into the updating pipeline at each step.
We have already mentioned that both algorithms are local. Indeed, the devices that actually perform the updating ignore where in the lattice is the site being updated, or the size of the system. They just process their input data and output the updated spins. There is another device (ADDRESS) which takes care of the geometry and addresses the memories accordingly. That component is also responsible for stopping the simulation when the desired number of configurations has been calculated.
### 3.2 Memory Scheme
Two different memory banks, for spins and couplings, are available to each UPDATE device (SPIN MEMORY and J-MEMORY areas in fig. 2).
Couplings are not dynamic variables (their values remain constant during the simulation), so the coupling memories are always in reading mode during update. When a site is to be updated, the couplings with its neigbours $`x_+`$, $`x_{}`$, $`y_+`$, $`y_{}`$, $`z_+`$, $`z_{}`$ must be supplied to the updating engine. We organize therefore the 18 SRAM devices as a single bank of width $`3\times 16`$ bits and depth $`6\times 64K`$. Each lattice takes 6 of the 48 bits to store the six needed couplings. The maximum simulable volume is then limited by the depth to L=73.
The spins changing during the simulation, an appropiate mechanism is needed in order to read and write the configurations simultaneously: The spin memory is duplicated (P and Q banks), and while one memory bank is read from the other is written on. Because of that, two memory devices are connected to every UPDATE component, each one capable to store $`64K\times `$ 18 bits.
In order to understand how the spin memory is managed, let us consider each column along the $`x`$ axis of the lattice divided in blocks of fixed length $`l`$.
To update one of the blocks, the block itself and its four $`y`$ and $`z`$ neighbours have to be supplied. So, five blocks must be read to update one, implying that, if an updated spin every clock cycle is wanted, the block length has to be at least five (see subsection 3.3 below).
Each spin memory device of 18 bit words stores two lattices, so 9 bits are available for each lattice. Each block can contain from 5 to 9 spins, and the maximum lattice size $`L(l\times 64K)^{1/3}`$ that can be stored is $`68`$, $`73`$ , $`77`$, $`80`$ or $`84`$, depending on the selected block size. This limitation, together with the one we found from the coupling memory and the fact that the number of blocks must be even, yields the range of simulable sizes shown in table 2.
The spin memories are arranged in the following way: Each 9-bit word contains $`l`$ consecutive spins ($`5l9`$), being consecutive (along x axis) lattice blocks stored in consecutive memory addresses. The $`V/l`$ words are not read consecutively, but following a pattern that makes the block to be updated and its neighbouring blocks available to the UPDATE component, as explained in the next subsection.
On the other hand, the coupling memories store in the $`n^{th}`$ 6-bit word the couplings of the $`n^{th}`$ spin with its six neighbours, and is read sequentially as the $`V`$ spins are updated.
### 3.3 Pipelined Updating
In this subsection we describe the logic programmed in every UPDATE device. We consider the case in which the demon algorithm is used with a block length $`l=5`$ (see fig. 3). In this case, the algorithm runs over a state machine with ten states. In each of those states a block is read from bank Q, which is in reading mode, and stored in one of the internal registers A…J. Let us suppose we have already been in states $`0\mathrm{}4`$, and some registers are already loaded: A ($`z_{}`$ neighbouring block), B ($`y_{}`$ neighbouring block), C (block to be updated), D ($`y_+`$ neighbouring block) and E ($`z_+`$ neighbouring block).
In state $`5`$, we send for update the first spin in the block stored in C. The $`x_+`$ neighbour is in the same block, the $`x_{}`$ neighbour is the previous updated spin, which is still in the updating process (and will not be needed until the last step), and the rest of the neighbours are stored in blocks A,B,D,E. We read simultaneously the $`z_{}`$ neighbouring block of the next block to be updated and store it in register F.
In states $`6`$ to $`8`$, we continue sending for update the spins second to fourth in block C, and loading registers G (next block $`y_{}`$ neighbour), H (next block to be updated) and I (next block $`y_+`$ neighbour). In state $`9`$, the last spin in block C is sent into the update pipeline. It is no longer true that the $`x_+`$ neighbour is in the same block, but it is in the block we have already stored in register H. Register J (next block $`z_+`$ neighbour) is loaded. When we return to state $`0`$, the updating of the block registered in H starts.
After some cycles, the updated value of the block that was stored in C has been calculated and is written on the appropiate memory position in bank P, in writing mode. Writing follows the same scheme as reading, not only the updated values are written but also the unchanged neighbours.
When a whole column (containing $`L/l`$ blocks of $`l`$ spins each) has been updated, we change the role of the memories: we will now write on bank Q, the memory bank we were previously reading from, and read from bank P, the bank we were writing on. Bank P stores now the correct new configuration. Bank Q stores the old configuration, which shall be overwritten with the result of updating the column read from bank P.
We have seen that the writing and reading sequences are equal, although writing is obviously delayed with respect to reading several cycles. Due to this delay, to avoid problems in the change of role of the spin memory banks the first block of the updated column, which should be read while its bank is still in writing mode, is stored in a cache memory inside the UPDATE devices. This mechanism requires at least four blocks, making the minimum simulable size to be $`L=20`$.
### 3.4 Addressing Logic
The ADDRESS device in fig. 2 controls and addresses the memories, establishes the functioning mode of the UPDATE and RNG devices and takes charge of the communications with the HC. As we said above, The board is accessed through a communication port with 32 bidirectional lines devoted to data transfer and 8 control lines (4 in each direction). Data lines are connected through tri-state circuits to the bus connecting the UPDATE and RNG devices. Control lines are connected to the ADDRESS chip, which controls the board according to the commands sent from the HC.
The implemented instruction set allows us to:
* program the devices.
* read/write the spin and coupling configurations.
* read/write the demon energies.
* load the number generator initialization table.
* load the probability tables used in the Heat Bath algorithm.
* set the number of iterations to run.
* start the simulation.
The ADDRESS device controls the UPDATE and RNG devices to carry out that operations. A 3-bit wide bus is used to encode the instructions for the UPDATE devices.
### 3.5 Random Number Generator
The Altera 10K50 device (RNG chip in fig. 2) is a 32-bit pseudo-random number generator of the R250 kind. Those generators are known to suffer some problems in Monte Carlo simulations, but only with non-local algorithms .
In the C implementation, a vector is initialized with a conventional pseudo-random number generator. Using the macro instruction RANDOM we run over the wheel, getting a new number and changing one of the values in the wheel:
```
#define RANDOM ( (irr[ip++]=irr[ip1++]+irr[ip2++])^irr[ip3++] )
```
The variables involved need to be properly initialized before using the defined macro:
```
/* random number generator initialization */
unsigned int irr[256];
unsigned char ip, ip1, ip2, ip3;
ip=128;
ip1=ip-24;
ip2=ip-55;
ip3=ip-61;
for(i=0; i<256; i++)
irr[i]=(unsigned int) rand();
```
In the RNG device (see Fig. 4), the `irr[i]` wheel becomes a 32-bit wide shift register, reproducing that way the effect of incrementing ip,ip1,ip2 and ip3. An adder sums the words WordA and WordB and stores the result in the first position IN\_SR. This result also serves as input to a XOR function, together with the value in the last register WordC. The result of this function provides us with the pseudo-random number, every clock cycle.
The seeds loading process is controlled by the ADDRESS component, which also enables the random number generation during the simulation.
### 3.6 Software
The boards are connected to the HC through a data acquisition card PCI-DIO32HS from National Instruments. To access the DAQ, a Linux driver has been programmed, and also a user library allowing to operate with the boards in an easy way.
The functions available to the user are the following:
* dioinit : initializes the DAQ boards to be used by the HC.
* boardsel : selects one board among those connected to the HC.
* ws : writes the spin configuration corresponding to one of the UPDATE devices in the selected board.
* rs : reads the spin configuration.
* wj : writes the couplings of the lattices in the selected board.
* rj : reads the couplings in the selected board.
* rd : reads the Demon energy.
* wd : writes the Demon energy.
* wmesfr : sets the number of iterations in each run.
* wrng : writes the initial random number table.
* wprob : writes the probability table 2 on UPDATE.
* startsue : starts the simulation in the selected board.
* waitsue : waits for one of the boards to finish.
The functions to access the memories get the arguments as arrays of bytes, where each element corresponds to one site in the lattice and the eight bits in the element to each one of the lattices in the board, as is usual in multi-spin code. The user needs not worry about SUE internal details.
### 3.7 Design considerations
#### 3.7.1 Programming method
CPLDs are electronic devices which can be programmed as many times as needed. They lose their program code every time the board is switched off, so they have to be reprogrammed after switch on. To manage the programming task, these devices are connected sequentially creating a JTAG chain which is controlled by the HC through the communication port described above. No extra cables are needed, providing easy on-board reprogrammability controlled from the HC. This feature was extremelly important during the debug process of the boards
#### 3.7.2 Printed Circuit Board
The printed circuit board surface is $`24.5\times 30.5cm^2`$, and it is $`2mm`$ thick. Manufactured in FR4 fiber, it consists of eight layers (four dedicated to signal transimision and the others to powering).
The board satisfies the ATX standard. In the final version, the 12 boards are mounted in a rack and are fed by a 800 Wat source at 5 V. Current, voltage and temperature are monitored. Full operation values are 90 Amp at 5 V.
#### 3.7.3 Frequency
The proper working of the circuit requires perfect synchronization between the active devices. The working frequency is 48 MHz, and the clock signal should reach the 32 components spread over a $`747cm^2`$ surface.
The clock distribution is made through CY2308-4 devices (3.3V Zero Delay Buffers), provided with a PLL mechanism (Phase Locked Loop) that allows to double the input frequency and supply eight outputs.
A 12 MHz oscillator is connected to a PLL device that doubles its frequency. Five outputs are driven into PLL components that double the frequency again and feed the neighbouring components. In this way, the clock is distributed across the circuit at low frequency, and the frequency doubled near the final components.
#### 3.7.4 Transmission lines
As a consequence of the large size of the circuit, there exist connections with a large total trace length. The rise times of the signals determine whether the transmission line behaves like a distributed circuit or not.
The effective length associated with the rise time of a signal is
$$l=\frac{T_r}{D}$$
(6)
where $`T_r`$ is the rise time and $`D`$ the propagation delay, characteristic of the material. We must consider a distributed circuit if the length of the transmission line is greater than a quarter of the effective length.
In our board, diode barriers protect the traces addressing the coupling memories, the data bus connecting UPDATE and RNG devices, and the connector for external communication. The rest of the signals, generated by memory devices or Altera 10K components (the latter allow the user to set the rise time), have rise times short enough for the system to behave in a lumped fashion.
## 4 Development Process
Initially, only one board was manufactured. In a first stage, we tested its general performance. After being able to communicate with the machine, the programming mechanism was implemented. Different test programs were written and compiled to program the CPLDs, using Altera’s MAXPlus+II development enviroment. Once we checked that all the components worked properly (fixing some electrical bugs in the way), the Demon algorithm was progressively implemented. We chose the demon algorithm because it is microcanonical, so the random number generator is not needed and the conservation of the total energy provides a fast test mechanism. Additional functionalities were added and the algortithm scheme was fine-tuned, until the program was complete.
When the rest of the boards were available, they were tested with this Demon program, and the Heat Bath algorithm was then implemented. The structure of the algorithm remained almost the same, although some details in the algorithm had to be changed and some new functions were added to the user library. The main novelty was the random number generator usage, which worked finally with an appropriate pipelining scheme both in RNG and UPDATE chips.
In the early debugging stage, development had been carried at 24 MHz, so we switched to high frequency. Some fine-tuning in the programs was needed and the CPLD logic layout was carefully studied in order to reach the design-goal of 48 MHz.
To make sure of the proper working of the machine beyond any doubt, an emulator was developed to run in a PC, so the machine configurations can be compared with those obtained with the PC emulation. This test proved that both the updating algorithm and the random number generator worked as intended.
## 5 Performance
Table 3 compares the update speed achieved by SUE with that of some simulations run by our group in different computers, and with the performance obtained running highly optimized multi-spin code in a Cray T3E supercomputer as reported in . We can see that the whole machine matches the computational power of one hundred processors of a CrayT3E.
## 6 Preliminary Physical Results
In this section we present some preliminary results obtained with SUE. We have run an $`L=20`$ lattice in $`4`$ boards at $`48`$ Mhz. We have simulated $`1600`$ sets of $`\{J_{i,j}\}`$, with $`2`$ replicas at $`12`$ different values of $`\beta `$, which previously have been controlled to have a correct transfer probability between them with the parallel tempering method. We measure every $`16384`$ sweeps, collecting $`800`$ measurements. These results have been obtained in $`60`$ days.
In Fig. 5 we plot the value of the average squared overlap.
The errors are plotted only at the simulations points. We are working around the critical region, as we reach high values for $`q^2`$. The extrapolated lines connect properly and the different values evolve smoothly for different $`T`$ values, as corresponds to a good thermalization and a high transition probability from parallel tempering.
At the moment of writing, we are running $`L=20`$ in 12 boards and almost finished the runs. Afterwards we will start the simulation in the $`L=30`$ system, and estimate that the time needed to obtain good results is around one year.
Acknowledgements
We wish to thank H.G. Ballesteros, J.M. Carmona, L.A. Fernández, D. Iñiguez, and J.J. Ruiz-Lorenzo for useful discussions. Partially supported by DGA (P46/97) and CICyT (AEN97-1768 and AEN99-0990).
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# Relevance of memory in Minority Games
## I Introduction
The Minority Game has been designed as the most drastic possible simplification of Arthur’s El Farol’s bar problem . It is believed to capture some essential and general features of competition between adaptative agents, which is found for instance in financial markets. In this model, agents have to take each time step one of two decisions; they share a common piece of information $`\mu \{0,\mathrm{},P1\}`$ that encodes the state of the world, use it to make their choice, and those who happen to be in minority are rewarded. In its original formulation, the piece of information is the binary encoding of the $`M`$ last winning choices, hence $`P=2^M`$. Hence, the dynamics of $`\mu `$ is coupled to the dynamics of agents.
Cavagna claimed that all quantities of the system “are completely independent from the memory of the agents”. This means that replacing the dynamics of $`\mu `$ induced by agents by a random history $`\mu `$ drawn at random at each time step, one finds the same results. While this statement has turned out to be wrong for many extensions of the MG , it has been helpful as a first approximation for the analytical understanding of the standard MG: an exact solution for random histories has been found in the “thermodynamic” limit . Interestingly, this solution shows that all quantities depend on the frequencies $`\{\rho ^\mu \}`$ of visit of histories. The random history case is recovered if $`\rho ^\mu =1/P`$, but in the real dynamics of the MG the distribution $`\rho ^\mu `$ is determined by the behavior of agents (indeed modifying the behavior of agents may have strong effects on $`\rho ^\mu `$ as shown in ref. ). It turns out, that the frequencies $`\rho ^\mu `$ are not uniform for all parameters of the MG.
In this paper we study quantitatively this problem. The first step is to characterize the properties of the dynamics of real histories, which amounts to study randomly biased diffusion on De Bruijn graphs. Depending on the asymmetry of the bias, we quantify the deviation $`\delta \rho ^\mu =\rho ^\mu 1/P`$ from the uniform distribution. Then we move to the MG and quantify the bias which agents induce on the dynamics of $`\mu `$ in the asymmetric phase. Using a simple parameterization of $`\rho ^\mu `$ which is inferred from numerical data, we generalize the calculations of refs. . This leads to a self-consistent equation between the asymmetry of the game and the diffusion bias, which we can solve. The results are in excellent agreement with numerical simulations and show a systematic deviation from the random history MG. Hence, our conclusion is that, even though the random history MG is qualitatively similar to the original MG, memory is actually not irrelevant, and one can quantify the difference between the two cases.
## II De Bruijn graphs
Let us begin with the definition of some elementary concepts. A binary sequence $`\mu (t)`$ of length $`m`$ consists of $`m`$ ordered elements $`\{b(tm),\mathrm{},b(t1)\}`$ where $`b`$ is a letter belonging to the alphabet $`\{0,1\}`$. $`\mu (t+1)`$ is obtained by adding $`b(t)`$ to the right of $`\mu (t)`$ and erasing $`b(tm)`$. Thus, for a given $`\mu (t)`$, there are two possible $`\mu (t+1)`$, which we call “next neighbours”. This updating rule defines the De Bruijn graph of order $`m`$ (see Fig. 1 for an example).
Let $`G`$ be the $`P\times P`$ adjacency matrix of the De Bruijn graph of order $`m`$. if we adopt the convention that its elements are indiced by the decimal value of the binary strings, that is, $`\mu =0,\mathrm{},P1`$,
$$G_{\mu ,\nu }=\delta _{[2\mu \%P],\nu }+\delta _{[2\mu \%P]+1,\nu }$$
(1)
where $`A\%B`$ stands for the remainder of the division of $`B`$ by $`A`$ and $`\delta _{i,j}`$ is the Kronecker symbol. The adjacency matrix for $`m=3`$ is:
$$G=\left(\begin{array}{cccccccc}1& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 1\\ 1& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 1\end{array}\right)$$
(2)
## III Unbiased diffusion
The unbiased diffusion is defined as follows : a particle moves on the directed De Bruijn graph $`G`$ and at each time step $`t`$, it jumps with equal probability to one of the next neighbours of the vertex it stands on at this time. Thus the transition probabilities matrix is $`W_0=G/2`$. In the long run, the fraction of time spent on vertex $`\nu `$ is given by $`[(W_0)^{\mathrm{}}]_{0,\nu }`$. It can be seen (see appendix) that
$$[(W_0)^k]_{\mu ,\nu }=\frac{1}{2^k}\underset{n=0}{\overset{2^k1}{}}\delta _{[2^k\mu \%P]+n,\nu }.$$
(3)
In particular, $`(W_0^{M+k})_{\mu ,\nu }=\frac{1}{P}`$ for all $`k0`$, that is, all strings $`\mu `$ are visited with the same frequency $`\rho ^\mu =1/P`$.
In order to have a intuitive feeling of those graphs, we write them for $`M=3`$ :
$$W_0^2=\frac{1}{4}\left(\begin{array}{cccccccc}1& 1& 1& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 1& 1& 1\\ 1& 1& 1& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 1& 1& 1\\ 1& 1& 1& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 1& 1& 1\\ 1& 1& 1& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 1& 1& 1\end{array}\right)W_0^3=W_0^4=W_0^5=\mathrm{}=\frac{1}{8}\left(\begin{array}{cccccccc}1& 1& 1& 1& 1& 1& 1& \\ 1& 1& 1& 1& 1& 1& 1& \\ 1& 1& 1& 1& 1& 1& 1& \\ 1& 1& 1& 1& 1& 1& 1& \\ 1& 1& 1& 1& 1& 1& 1& \\ 1& 1& 1& 1& 1& 1& 1& \\ 1& 1& 1& 1& 1& 1& 1& \\ 1& 1& 1& 1& 1& 1& 1& \end{array}\right)$$
(4)
## IV Randomly biased diffusion
The perturbations are introduced by adding a term to the transition probabilities matrix $`W_ϵ=W_0+ϵW_1`$ where $`ϵ`$ quantifies the asymmetry and $`W_1`$ contains the disorder $`\xi `$
$$(W_1)_{\mu ,\nu }=(1)^\nu \xi _\mu (W_0)_{\mu ,\nu }$$
(5)
where the $`\xi `$ are iid from the pdf $`P(\xi )=1/2\delta (\xi 1)+1/2\delta (\xi +1)`$ and the $`(1)^\nu `$ comes from the normalization of the perturbed probabilities. We are looking for the stationary transition probabilities, i.e., $`W_ϵ^{\mathrm{}}`$ such that $`W_ϵ^{\mathrm{}}=lim_k\mathrm{}\left(W_ϵ\right)^k`$. It exists since $`W_ϵ`$ is a bounded operator. Its formal series expansion in $`ϵ`$ is noted by $`W_ϵ^{\mathrm{}}=_{k0}ϵ^kW_k^{\mathrm{}}`$ where $`W_0^{\mathrm{}}`$ is a matrix whose all coefficients are equal to $`1/P`$ (see above). The relationship $`W_ϵ^{\mathrm{}}=W_ϵ^{\mathrm{}}W_ϵ`$ provides the recurrence
$$W_k^{\mathrm{}}=W_k^{\mathrm{}}W_0+W_{k1}^{\mathrm{}}W_1$$
(6)
Since $`W_k^MW_0^{\mathrm{}}=0`$, we iterate $`m1`$ times this equation by replacing $`W_k^{\mathrm{}}`$ with $`W_k^{\mathrm{}}W_0+W_{k1}^{\mathrm{}}W_1`$ in the r.h.s., yielding to
$$W_k^{\mathrm{}}=W_{k1}^{\mathrm{}}W_1V=W_0^{\mathrm{}}\left[W_1V\right]^k,$$
(7)
where $`V=_{c=0}^{M1}\left(W_0\right)^c`$. At this point, it is useful to remark that multiplying a matrix on the left by $`W_0^{\mathrm{}}`$ is equivalent to averaging its columns :
$$(W_0^{\mathrm{}}A)_{\mu ,\nu }=\underset{a=0}{\overset{P1}{}}(W_0^{\mathrm{}})_{\mu ,a}A_{a,\nu }=\frac{1}{P}\underset{a=0}{\overset{P1}{}}A_{a,\nu }=\text{average of the }\text{ν}\text{-th column of }\text{A}$$
(8)
thus the matrices $`W_k^{\mathrm{}}`$ consist of averages of columns of $`(W_1V)^k`$. Therefore, $`(W_k^{\mathrm{}})_{\mu ,\nu }`$ is the $`k`$-th order correction to the frequency of vertex $`\nu `$, that will be called $`\rho _{(k)}^\nu `$ in the following. Note that $`\rho _{(k)}^\nu _\xi =0`$ for all $`k1`$. The square root of the second moment of $`\rho _{(k)}^\nu `$ averaged over the disorder gives an indication of the typical value of $`\rho _{(k)}^\nu `$. In appendix B we obtain the approximation
$$\rho _{(k)}^2_\xi \frac{(11/P)^k}{P}.$$
(9)
which is exact for the first order perturbation. Therefore $`\rho _{(k)}^\nu `$ is of the same order as the unperturbed $`\rho _{(0)}^\nu `$, thus it cannot be neglected. Fig 2 shows that the behavior predicted by Eq (9) is indeed correct for large $`P`$.
Finally, one can estimate the second moment of $`\rho ^\nu `$. If one supposes that the perturbations at different orders are independent, one obtains
$$\mathrm{\Delta }\rho ^2=\frac{1}{P}\underset{\nu =0}{\overset{P1}{}}\left[[\rho ^\nu ]^2_\xi \rho ^\nu _\xi ^2\right]\frac{1}{P^2}\left[\frac{1}{1(11/P)ϵ^2}1\right]\frac{1}{P^2}\frac{ϵ^2}{1ϵ^2}$$
(10)
## V Application to MG
Let us first define the game<sup>*</sup><sup>*</sup>*See refs for more details: MG consists of $`N`$ agents trying to be at each time step in minority. Each agent has $`S`$ strategies, or lookup tables $`a_{i,s}`$, $`s=1,\mathrm{},S`$ , and dynamically assigns a score to each of them . At each time step $`t`$, the system’s history $`\mu (t)`$ is made available to all agents; the latter use their best strategythe one with the highest score $`s_i(t)`$ take the decision $`a_{i,s_i(t)}^\mu (t)=+1`$ or $`1`$ and a market maker sums up all decisions into the aggregate quantity $`A(t)=_{i=1}^Na_{i,s_i(t)}^\mu (t)`$ .
Macroscopic quantities of interest include the temporal averages of $`A(t)`$ conditional to $`\mu (t)=\mu `$, for all $`\mu `$, noted by $`A^\mu `$. The MG undergoes a second order phase transition with symmetry breaking as the control parameter $`\alpha =P/N`$ is varied : the system is in the symmetric phase ($`A^\mu =0`$ for all $`\mu `$) if $`\alpha <\alpha _c`$ and it is in the asymmetric phase for $`\alpha >\alpha _c`$. One convenient order parameter is$`\overline{R}=_\mu \rho ^\mu R^\mu `$ is the notation for the weighted average over the histories. $`H=\overline{A^2}`$: it is equal to zero in the symmetric phase, and grows monotonically with $`\alpha `$ in the asymmetric phase (see Fig. 4). One other relevant macroscopic quantity is the fluctuations $`\sigma ^2=\overline{A^2}`$ which quantifies the performance of the agents.
Before doing any analytic calculations, it is worth looking at Figures 3 and 4 which clearly show that Cavagna’s assertion is right as long as the system is in the symmetric phase. Indeed, if $`A^\mu =0`$, the transition probabilities from $`\mu `$ to its next neighbours are unbiased, that is $`ϵ^\mu =0`$; therefore in the symmetric phase, where $`A^\mu =0`$ for all $`\mu `$, the frequencies of visit are uniform $`\rho ^\mu =1/P`$. Accordingly, numerical simulations show that these quantities collapse on the same curve.
As $`\alpha `$ increases, the critical point is crossed, and $`A^\mu 0`$ for some $`\mu `$. The dynamics of the history is biased on all such histories and consequently all macroscopic quantities are significantly different: both $`\sigma ^2/N`$ and $`H/N`$ are lower for real histories than for uniformly sampled histories. This can be understood by the facts that $`\sigma ^2/N`$ and $`H`$ are increasing functions of $`\alpha `$ and that the biases on the De Bruijn graph of histories reduce the effective number of histories, that can be defined as $`2^{\overline{\mathrm{log}_2\rho }}`$: in other words, effective $`\alpha `$ of MG with real histories is smaller than that of MG with uniform histories. This explanation is indeed confirmed by Fig 5; this shows the fraction of frozen agents<sup>§</sup><sup>§</sup>§See : they are agents that stop to be adaptative. $`\varphi `$ which is a decreasing function of $`\alpha `$ in the asymmetric phase. As expected from the above argument, $`\varphi `$ of MG with real histories is larger than that of MG with uniformly sampled histories.
The bias $`ϵ^\mu `$ on a particular history can be estimated for large $`N`$: in this limit $`A^\mu `$ is a Gaussian variable with average $`A^\mu `$ and variance $`(A^\mu )^2A^\mu ^2`$, leading to
$$ϵ^\mu =\text{sgn}(A^\mu )ϵ_{th}^\mu =\text{erf}\left(\sqrt{\frac{A^\mu ^2}{2[(A^\mu )^2A^\mu ^2]}}\right).$$
(11)
Fig 6 confirms the validity of Eq. (11). The figure also shows that $`ϵ^\mu `$ are unevenly distributed and they are not equal even if the system is deep in the asymmetric phase ($`\alpha 8.5`$ in this figure). Indeed, as a function of $`\mu `$, $`A^\mu `$ is a random variable with average 0 and variance $`H`$, which is an increasing function of $`\alpha `$. Since we studied diffusion of perturbed graphs with only one parameter $`ϵ`$, we have to map all $`ϵ^\mu `$ onto a scalar quantity, so that we define $`ϵ`$ as the non weighted averageThis is clearly an important assumption, but the diffusion on De Bruijn graphs with one $`ϵ^\mu `$ per site leads to a much greater complexity. As it appears on Fig 3, 4 and 9, this assumption is not unrealistic. of $`ϵ^\mu `$ over the histories. For large $`P`$, $`ϵ`$ can be approximated by
$$\overline{ϵ}_{th}=2_0^{\mathrm{}}\text{d}A\frac{e^{\frac{A^2}{2H}}}{\sqrt{2\pi H}}\text{erf}\left(\frac{A}{\sqrt{2(\sigma ^2H)}}\right).$$
(12)
Here both $`H`$ and $`\sigma ^2`$ can be computed analytically with the method of refs. (see the appendix). However the solution depends on the distribution $`\rho ^\mu `$. In order to make Eq. (12) a self-consistent equation for $`\overline{ϵ}_{th}`$, we need to parameterize the distribution of $`\rho ^\mu `$ by $`\overline{ϵ}_{th}`$ itself.
We could not find ab initio the analytic form of the pdf of $`\{\rho ^\mu \}`$, but Fig 7 shows that
$$P(\tau )\frac{(\lambda +1)^{\lambda +1}}{\mathrm{\Gamma }(\lambda +1)}\tau ^\lambda e^{(\lambda +1)\tau }$$
(13)
is a very good approximation for the pdf of $`\rho =\tau /P`$. The parameter $`\lambda `$ is easily connected with $`\overline{ϵ}_{th}`$:
$$\tau ^2\tau ^2=\frac{1}{1+\lambda }=P^2\mathrm{\Delta }\rho ^2\frac{\overline{ϵ}_{th}^2}{1\overline{ϵ}_{th}^2}$$
(14)
where we used Eq. (10). This gives $`\lambda (12\overline{ϵ}_{th}^2)/\overline{ϵ}_{th}^2`$. Note that this approximation requires $`\overline{ϵ}_{th}<1/\sqrt{2}`$.
This turns Eq. (12) into an equation for $`\overline{ϵ}_{th}`$, and the theory is self consistent. Figure 8 reports measured $`ϵ`$ and its approximation $`\overline{ϵ}_{th}`$. What clearly appears from this figure is that $`ϵ`$ is far from being negligible, and that $`\overline{ϵ}_{th}`$ is a quite good approximation to $`ϵ`$.
We can also check the validity of Eq. (10) against the self-consistent theory. Fig 9 shows that Eq. (10) is in good agreement with numerical simulations as long as all histories are visited. Moreover the approximation $`\overline{ϵ}_{th}`$ for $`ϵ`$ leads to qualitatively similar results, but underestimates $`\mathrm{\Delta }\rho ^2`$ because $`\overline{ϵ}_{th}<ϵ`$ (see figure 8).
The self-consistent replica calculation for the Minority Game of refs. with the ansatz $`\rho =\tau /P`$ and $`\tau `$ given by the pdf (13) is discussed in the appendix. Fig 3 and 4 indicate that analytic predictions are well supported by numerical simulations.
In the asymmetric phase, which is arguably the most relevant and interesting in the MG , all quantities of MG change significantly if one replaces real histories with random uniform histories. A dependence on the frequencies $`\rho ^\mu `$ does not necessarily imply the relevance of the detailed dynamics of the histories. If the histories $`\mu `$ where drawn randomly from the “correct” distribution $`\rho ^\mu `$, the results would be the same (actually it suffices to know the pdf of $`\rho ^\mu `$). The problem is that the distribution $`\rho ^\mu `$ depends on the asymmetry $`A^\mu `$, which in turn depend on the microscopic constitution of all agents . In other words, $`\rho ^\mu `$ is a self-consistently determined quantity and hence it is only known a posteriori.
## VI Conclusion
We have shown that the dynamics of histories cannot be considered as irrelevant. Indeed, even for the canonical MG, it is relevant and cannot be replaced by randomly drawn histories. In addition, for many extensions and variations of the MG, the dynamics of histories is not only relevant, but crucial.
We acknowledge fruitful discussions with Philippe Flajolet and Paolo De Los Rios. This work has been partially supported by the Swiss National Science Foundation under Grant Nr 20-46918.98.
## A
Let us prove by induction that
$$(W_0^k)_{\mu ,\nu }=\frac{1}{2^k}\underset{n=0}{\overset{2^k1}{}}\delta _{[2^k\mu \%P]+n,\nu }$$
(A1)
It is sufficient to calculate explicitly $`(W_0^k)_{\mu ,\nu }`$ from $`(W_0^{k1})_{\mu ,\nu }`$
$`(W_0^k)_{\mu ,\nu }`$ $`=`$ $`{\displaystyle \underset{\tau =0}{\overset{P1}{}}}(W_0)_{\mu ,\tau }(W_0^{k1})_{\tau ,\nu }`$ (A2)
$`=`$ $`{\displaystyle \underset{n=0}{\overset{2^{k1}1}{}}}\left\{\delta _{[2^{k1}([2\mu \%P])\%P]+n,\nu }+\delta _{[2^{k1}([2\mu \%P]+1)\%P]+n,\nu }\right\}`$ (A3)
$`=`$ $`{\displaystyle \frac{1}{2^k}}{\displaystyle \underset{n=0}{\overset{2^k1}{}}}\delta _{[2^k\mu \%P]+n,\nu }`$ (A4)
since $`A(B\%P)\%P=AB\%P`$ and $`(2^k\mu +2^{k1})\%P=[2^k\mu \%P]+2^{k1}`$ if $`P=2^M`$ and $`km1`$.
## B
In order to simplify the notations, we define
$$(X^c)_{\mu ,\nu }=\underset{n=0}{\overset{2^c1}{}}\delta _{[2^{c+1}\mu \%P]+n,\nu }\delta _{[2^{c+1}\mu \%P]+n+2^c,\nu }$$
(B1)
This matrix is such that
$$(X^c)_{\mu ,\nu }=\{\begin{array}{ccc}\hfill 1& \text{if}& 2^{c+1}\mu \%P\nu <[2^{c+1}\mu \%P]+2^c\}\hfill \\ \hfill 1& \text{if}& [2^{c+1}\mu \%P]+2^c\nu <2^{c+1}(\mu +1)\%P\hfill \\ \hfill 0& \text{else}& \end{array}$$
(B2)
With this formalism, one can write $`W_1V`$ as
$$(W_1V)_{\mu ,\nu }=\frac{\xi _\mu }{2}\underset{c=0}{\overset{M1}{}}\frac{1}{2^c}(X^c)_{\mu ,\nu }$$
(B3)
Let us calculate the perturbation at order 1: one has to compute $`\rho _{(1)}^2`$ in order to have an estimation of the typical value of a generic $`\rho _{(1)}^\nu `$: since the $`\xi `$ are uncorrelated and $`_{\nu =0}^{P1}(X^c)_{\mu ,\nu }(X^d)_{\mu ,\nu }=2^{c+1}\delta _{c,d}`$ ,
$$\rho _{(1)}^2_\xi =\frac{1}{4P^2}\underset{\mu ,\nu =0}{\overset{P1}{}}\underset{c=0}{\overset{M1}{}}\frac{[(X^c)_{\mu ,\nu }]^2}{2^{2c}}=\frac{(11/P)}{P}$$
(B4)
The next orders of perturbation are much harder to handle. However, for large $`P`$, one can approximate them by supposing that
$$\rho _{(k)}^2_\xi (11/P)\rho _{(k1)}^2_\xi =\frac{(11/P)^k}{P}.$$
(B5)
Consequently, $`\rho _{(k)}^\nu (11/P)^{k/2}\frac{1}{P}\frac{1}{P}`$ at leading order
## C
Since agents actually minimize $`H/N`$, one can consider this quantity as a Hamiltonian and find its ground state. This is possible by methods of statistical physics such as replica trick . The generalization of the calculus of refs to $`\rho ^\mu =\tau ^\mu /P`$ drawn from the pdf given by Eq (13) and (12) is straightforward; the free energy reads in the thermodynamic limit
$`F(\beta ,Q,q,R,r)`$ $`=`$ $`<{\displaystyle \frac{\alpha }{2\beta }}\mathrm{log}[1+\chi \tau ]>_\tau +{\displaystyle \frac{1+q}{2}}<{\displaystyle \frac{1}{\frac{1}{\tau }+\chi }}>_\tau `$ (C1)
$`+`$ $`{\displaystyle \frac{\alpha \beta }{2}}(RQrq){\displaystyle \frac{1}{\beta }}\mathrm{log}{\displaystyle _1^1}𝑑se^{\beta (\zeta s^2\sqrt{\alpha r}zs)}_z`$ (C2)
where $`\chi =\beta (Qq)/\alpha `$ and $`\zeta =\sqrt{\alpha /r}\beta (Rr)`$. Next, the $`\beta \mathrm{}`$ limit is taken while keeping finite $`\chi `$ and $`\zeta `$. One obtains
$$H=\frac{1+Q}{2}\left[<\frac{1}{\frac{1}{\tau }+\chi }>_\tau \chi <\frac{1}{[\frac{1}{\tau }+\chi ]^2}>_\tau \right]$$
(C3)
and
$$\sigma ^2=H+\frac{1Q}{2}$$
(C4)
where $`Q`$ and $`\chi `$ take their saddle point values, given by the solution of
$`Q(\zeta )`$ $`=`$ $`1\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{e^{\zeta ^2/2}}{\zeta }}\left(1{\displaystyle \frac{1}{\zeta ^2}}\right)\text{erf}\left({\displaystyle \frac{\zeta }{\sqrt{2}}}\right)`$ (C5)
$`Q(\zeta )`$ $`=`$ $`{\displaystyle \frac{1}{\alpha }}\left[{\displaystyle \frac{\text{erf}(\zeta /\sqrt{2})}{\chi \zeta }}\right]^2{\displaystyle \frac{1}{\frac{1}{[1/\tau +\chi ]^2}_\tau }}1`$ (C6)
$`\chi <{\displaystyle \frac{1}{\frac{1}{\tau }+\chi }}>_\tau `$ $`=`$ $`{\displaystyle \frac{\text{erf}(\zeta /\sqrt{2})}{\alpha }}`$ (C7)
Eqs (C6) and (C7), together with Eq (12), form a closed set of equations that has to be solved numerically. Note that as in the random histories case, $`\chi `$ becomes infinite at the critical point, where $`\alpha _c=\text{erf}(\zeta /\sqrt{2})`$.
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# Cooling neutron stars with localized protons
<sup>1</sup><sup>1</sup>institutetext: A.F. Ioffe Physical Technical Institute, Politekhnicheskaya 26, 194021, St.Petersburg, Russia
e-mail: baiko@astro.ioffe.rssi.ru <sup>2</sup><sup>2</sup>institutetext: N. Copernicus Astronomical Center, Bartycka 18, 00-716, Warsaw, Poland
e-mail: haensel@camk.edu.pl (Accepted 20th January 2000 )
## Abstract
We analyze cooling of neutron stars, assuming the presence of localized protons in the densest region of their cores. Choosing a single threshold density for proton localization and adjusting neutron star mass, we reproduce the observational data on effective surface temperatures of Vela and PSR 0656+14, with or without an accreted hydrogen envelope. However, the presence of a tiny hydrogen envelope is mandatory, in this model, for reproducing the Geminga data.
###### Key Words.:
neutron stars–cooling–localized protons offprints: P. Haensel
The model of dense neutron star matter with localized protons was proposed by Kutschera & Wójcik (KW93 (1993)). According to these authors, in high-density matter with low proton fraction the “zero-point” (Fermi) kinetic energy of protons, estimated from the uncertainty principle, becomes small, compared to the energy of proton interactions with density waves of the neutron background. These density waves resemble hydrodynamic sound or long wavelength acoustic phonons in a solid. The coupling of protons to the neutron density waves results in a significant increase of proton effective mass in close analogy with polaron behaviour of a slow electron in a polar solid which also acquires a large effective mass because of the interactions with lattice phonons.
Moreover, considering one proton in neutron matter, Kutschera & Wójcik have shown that above some critical density $`\rho _{\mathrm{lp}}4\rho _0`$ ($`\rho _0=2.8\times 10^{14}`$ g cm<sup>-3</sup>) the energy of a quantum state, in which a proton wave function is localized around some spatial point associated with the neutron density minimum, is lower than the energy of a state, where proton wave function is nonlocalized and the distribution of neutrons is uniform. Such a state also has an analog in solid state physics, the so-called small polaron state. It occurs when an electron interacts with a lattice so strongly that the latter deforms and this deformation traps the electron in a self-consistent manner.
This result by Kutschera & Wójcik may be applied to neutron star cores, where one has finite admixture of protons (rather than one proton in neutron matter), if the proton fraction $`x_p`$ is so small that wave functions of neighbouring localized protons do not overlap. Small proton fractions in high density neutron-star matter were obtained by Wiringa et al. (WFF88 (1988)), who performed extensive many-body calculations of the ground state, based on variational method and assuming realistic nuclear hamiltonian. In particular, their UV14+TNI model, which will be used further, predicts a proton fraction below 5% for densities $`\rho >3\rho _0`$, and, eventually, complete disappearance of protons at about $`7\rho _0`$. This result is in contradiction with some other calculations of equation of state (EOS) of neutron star matter \[based on the relativistic mean-field approach, see Glendenning (G96 (1996)), or performed within the Brueckner-Bethe-Goldstone many-body theory, e.g., Baldo et al. (BBB97 (1997))\] which yield a proton fraction growing monotonically with baryon density $`n_b`$. The high density behaviour of $`x_p(n_b)`$ is, therefore, subject to a significant theoretical uncertainty. Arguments in favor of $`x_p(n_b)`$, which decreases and eventually vanishes at high $`n_b`$, were presented by Kutschera (K94 (1994)).
The localization of protons (if confirmed) leads to dramatic changes of basic kinetic properties of matter in the neutron star core. The reason is that the localized protons are very efficient scatterers of the main transport agents, electrons and neutrons. In effect, in the low-temperature regime, the thermal conductivity $`\kappa `$, for instance, behaves as $`T`$, instead of the conventional $`T^1`$ dependence when scattering is provided by degenerate particles. As shown by Baiko & Haensel (BH99 (1999)) (hereafter Paper I), this extends the time scale of thermal equilibration in the neutron star core by about two orders of magnitude independently of the state of neutrons (normal or superfluid).
The neutrino emissivity of matter with localized protons is also very non-standard. First of all, the conventional beta-processes which involve charged currents (direct or modified Urca processes) are forbidden since at given baryon density the minimum of energy corresponds to a specific proton fraction, while any process which changes this fraction would require rearrangement of the whole state and, consequently, much larger energy than the thermal energy available, say, at $`T10^9`$ K. Thus, the only possible processes are those which involve neutral currents. These are bremsstrahlung of neutrino–antineutrino pairs in $`nn`$, $`np`$, and $`ep`$ collisions.
The kinetic properties of matter with localized protons were considered in Paper I. In particular, simple analytical formulae were derived for the electron and neutron thermal conductivities and shear viscosities, for the electron electrical conductivity as well as for the neutrino emissivity from $`np`$ and $`ep`$ bremsstrahlung. All these quantities were calculated under two simplifying assumptions. Firstly, it was assumed that the localized protons did not exibit any correlations (i.e., behaved like impurities). Secondly, the in-vacuum matrix elements of strong interactions were used. Both assumptions lead to an increase of the reaction rates by a factor of a few, since it is generally expected that the medium effects and correlation of scatterers reduce the matrix elements of elementary scattering processes, at least for small momentum transfers. Thus the results of Paper I give the upper bounds to the neutrino emissivities and the lower bounds to the transport coefficients. In other words, they represent the maximum effect that one can expect from the proton localization.
The aim of the present paper is to apply the results of Paper I to modelling neutron star cooling in order to check if the hypothetical proton localization is compatible with observational data on the neutron star effective surface temperatures. To do that we have to specify the neutron star model. In the core (at densities $`\rho >\rho _{\mathrm{cr}}=2.5\times 10^{14}`$ g cm<sup>-3</sup>), we take the UV14+TNI EOS by Wiringa et al. (WFF88 (1988)), while in the crust (below $`\rho _{\mathrm{cr}}`$) we take the EOS by Negele & Vautherin (NW73 (1973)) above neutron-drip density, and the EOS by Baym et al. (BPS71 (1971)) below neutron-drip density. The model is then determined by the only parameter, the central mass density $`\rho _\mathrm{c}`$, and is constructed by integrating the Tolman-Oppenheimer-Volkoff equation outward the center.
We use the cooling code described, e.g., in Levenfish et al. (LSY99 (1999)). The code includes explicitly effects of General Relativity (GR) and traces a steady thermal evolution of a spherically symmetric neutron star. Matter at densities $`\rho >\rho _\mathrm{b}=10^{10}`$ g cm<sup>-3</sup> is assumed to be isothermal with constant temperature $`T_\mathrm{i}^{\mathrm{}}`$ (as measured at infinity). The local temperature, $`T_\mathrm{i}`$, appears to be dependent on the radial coordinate $`r`$ due to the GR effect. The cooling is governed by the thermal balance equation with energy losses due to neutrino emission from the core and photon emission from the stellar surface. The local effective surface temperature $`T_\mathrm{s}`$ is obtained from the local internal temperature, $`T_\mathrm{b}`$, at $`\rho =\rho _\mathrm{b}`$ using the formula by Potekhin et al. (PCY97 (1997)). The latter formula fits the results of numerical simulations of heat transport through the crust. The formula depends also on the mass $`\mathrm{\Delta }M`$ of the outer envelope made of light elements (light envelope). Such an envelope may, in principle, reside on the neutron star surface and due to its higher thermal conductivity the envelope would decrease the thermal insulation of the inner stellar regions and modify the $`T_\mathrm{s}(T_\mathrm{b})`$ relationship. The surface temperature $`T_\mathrm{s}^{\mathrm{}}`$, as measured by a distant observer, is related to $`T_\mathrm{s}`$ as $`T_\mathrm{s}^{\mathrm{}}=T_\mathrm{s}\sqrt{12GM/Rc^2}`$, where $`M`$ and $`R`$ are, respectively, the mass and radius of the neutron star.
For our purposes we have introduced a number of modifications into the cooling code. All the changes are related to matter at densities above $`\rho _{\mathrm{lp}}`$. First of all, we have removed the proton contribution to the heat capacity. Then we have revised contributions to neutrino emissivity. Above $`\rho _{\mathrm{lp}}`$, if neutrons are nonsuperfluid (which is the case in study), the main contribution to the neutrino emissivity is due to $`nn`$ and $`np`$ neutrino-pair bremsstrahlung. The rate of the $`nn`$ process is not changed by the proton localization, while the emissivity of the $`np`$ process reads (Paper I):
$`Q_{\mathrm{Brem}}^{np}`$ $`=`$ $`1.3\times 10^{21}\left({\displaystyle \frac{n_b}{4n_0}}\times {\displaystyle \frac{x_p}{0.01}}\right)\left({\displaystyle \frac{n_n}{4n_0}}\right)^{4/3}\left({\displaystyle \frac{m_n^{}}{m_n}}\right)^2`$ (1)
$`\times `$ $`S_{np}^\mathrm{t}T_9^6\mathrm{ergs}\mathrm{s}^1\mathrm{cm}^3,`$
where $`x_p`$ is the fraction of protons by number, $`n_n`$ is the neutron number density, $`m_n`$ and $`m_n^{}`$ are the bare and effective neutron masses, respectively; $`n_0=0.16`$ fm<sup>-3</sup>, and $`T_9`$ is (local) temperature in units of $`10^9`$ K. Finally, $`S_{np}^\mathrm{t}`$ is a function of $`n_n`$ related to the total $`np`$ cross-section in vacuum,
$$S_{np}^\mathrm{t}=26.09+3.444\frac{4n_0}{n_n}.$$
(2)
Typical cooling curves obtained within this model are presented in Fig. 1. The effect of the proton localization can be easily seen by comparing the dotted curves, which represent the cooling without the proton localization, with curves of other types, which show a faster cooling due to the localization effect. The rate of cooling without the proton localization is virtually independent of the mass: the dotted curves for 1.44 $`M_{}`$ and 1.48 $`M_{}`$ stars coincide. If the effect of the proton localization is taken into account, the more massive star cools more rapidly because at fixed threshold density, $`\rho _{\mathrm{lp}}`$, the more efficient neutrino pair bremsstrahlung in neutron-localized proton collisions is operative in a larger region of the neutron star core leading to a larger overall neutrino luminosity.
The results of our cooling simulations may be used for interpretation of thermal emission from the surface of cooling neutron stars of ages $`t10^4`$ yr, which is the typical thermal equilibration time for neutron stars with localized protons (Paper I).
At present there are 4 sources, identified as isolated neutron stars, which show thermal soft X-ray emission and have characteristic ages $`t`$ in excess of $`10^4`$ yrs. These are Vela, PSR 0656+14, Geminga and PSR 1055-52. The data that we use are summarized in Table 1. The effective surface temperatures $`T_\mathrm{s}^{\mathrm{}}`$ are determined from fitting the observed spectra of the sources in two different ways, either by black-body spectra ($`T_{\mathrm{bb}}^{\mathrm{}}`$ in Table 1) or by spectra obtained in realistic models of magnetized neutron star atmospheres ($`T_H^{\mathrm{}}`$ in Table 1 for magnetized hydrogen atmosphere). The inferred values of $`T_\mathrm{s}^{\mathrm{}}`$ appear to be quite different. Radiation emerging from an atmosphere composed of light elements (hydrogen and helium) has harder spectrum than the black-body radiation with the same $`T_\mathrm{s}^{\mathrm{}}`$ due to the decrease of opacity at higher energies (Romani R87 (1987)). Consequently, the effective temperatures predicted by the atmosphere models are 2-3 times lower than the black-body ones which require lower $`d/R`$ (distance to radius) ratio. With the stellar radius fixed around the standard value of 10 km this translates into about 10 times smaller distances to the objects. The atmospheres composed of iron, on the contrary, produce radiation with a softer spectrum which is more similar on average to that of a black body.
The magnetic fields inferred from the spin-down rates are above $`10^{12}`$ G for all 4 stars. Such fields must have significant effect on the atmospheric opacities and therefore will modify the emerging spectrum (Shibanov et al. SZPV92 (1992)). In general, the magnetic field makes the spectrum softer and somewhat elevates the effective temperature compared to the non-magnetized atmosphere case. However, no accurate magnetized atmosphere model have been developed so far for the relevant temperature range ($`T_\mathrm{s}<10^6`$ K) where one should take into account effects of atom motion on the opacity (e.g., Pavlov & Zavlin, PZ98 (1998)). These effects were neglected while obtaining the values reproduced in Table 1.
The surface temperatures $`T_\mathrm{s}^{\mathrm{}}`$ are plotted in Fig. 1. The filled and open circles correspond to black-body models and simplified magnetized hydrogen atmosphere models, respectively. As one concludes from Table 1 and the figure, the overall theoretical uncertainty in surface temperatures appears to be quite large; for instance $`T_\mathrm{s}^{\mathrm{}}(26)\times 10^5`$ K for Geminga. This does not allow one to draw any definite conclusion about the cooling scenario of the neutron star except that it probably rejects the “rapid” cooling (via direct Urca process unsuppressed by superfluidity in dense matter). We note that the surface temperatures obtained from the hydrogen atmosphere models are rather low and cannot be explained within the “standard” cooling scenario (the neutrino losses via modified Urca process from non-superfluid matter). On the other hand, such a scenario can explain the black-body temperatures. Thus, the “atmospheric” temperatures (if the true temperatures turn out to be close to them) may provide more stringent test of the theory of neutron-star interior. For this reason let us adopt the “atmospheric” interpretation of observations and focus on lower error bars of Vela, PSR 0656+14, and Geminga.
Having fixed EOS we are left with three parameters which influence the cooling. These are the central stellar density $`\rho _\mathrm{c}`$ and the threshold values of mass density, $`\rho _{\mathrm{lp}}`$, and proton fraction, $`x_{p0}`$. In addition, we can vary $`\mathrm{\Delta }M`$, the mass of the surface envelope made of light elements. Note, that the presence of the hydrogen atmosphere is not equivalent to the presence of a light envelope as the amount of hydrogen needed to modify the emerging spectrum (1-10 g cm<sup>-2</sup> or $`10^{20}10^{19}M_{}`$) is much lower than that required to change the $`T_\mathrm{s}(T_\mathrm{b})`$ relationship. Let us fix the threshold proton fraction $`x_{p0}=5\%`$, and consider two models: without envelope, Fig. 2, and with a light envelope of mass $`\mathrm{\Delta }M=10^{10}M`$, Fig. 3. The strips bounded by the lines of various types in Figs. 2 and 3 correspond to the domain of $`\rho _{\mathrm{lp}}`$ and $`M`$, for which cooling curves cross the error bar of a given source. Let us stress that Geminga (whose parameter domain is shaded in Fig. 3) can be explained only assuming the light envelope. In the latter case and for realistic $`\rho _{\mathrm{lp}}`$ (presumably above $`4\rho _0`$) the mass of Geminga should be above 1.5 $`M_{}`$. The other two sources can be explained either with or without the envelope but for fixed $`\rho _{\mathrm{lp}}`$ their masses in the model with envelope should be higher. Finally let us mention that the strips are not very sensitive to the envelope mass for $`\mathrm{\Delta }M10^{10}M`$.
Fig. 1 illustrates the cooling curves for the fixed threshold density $`\rho _{\mathrm{lp}}=4.3\rho _0`$. The solid curve represents the model with the light envelope and $`M=1.61M_{}`$ ($`\rho _\mathrm{c}=5.4\rho _0`$). It crosses the “atmospheric” error bars of Vela, PSR 0656+14 and Geminga simultaneously. The dashed and dot-dashed curves are envelope-free models. The masses should be lower, for instance, $`1.44`$ and $`1.48M_{}`$ for Vela and PSR 0656+14, respectively.
We have performed cooling simulations of neutron stars, with realistic equation of state, assuming localization of protons above some threshold density. These results have been used for interpretation of effective surface temperatures of observed isolated, middle-aged neutron stars ($`t10^4`$ yr). We have shown that the available observational data are consistent with the proton localization and can be reproduced if one chooses model parameters within the specific domains in the $`\rho _{\mathrm{lp}}M`$ plane. Choosing a single threshold density for proton localization and adjusting neutron star mass, we reproduce the observational data for Vela and PSR 0656+14, with or without an accreted, hydrogen envelope. However, the presence of the hydrogen envelope is required in this model for explaining the observations of Geminga.
In conclusion, we note that the localization of protons is not necessary for the explanation of the available data on neutron star effective surface temperatures. Nevertheless, as is seen from the Fig. 1, it changes significantly the cooling rate and, if its existence is confirmed, it represents an important effect to be taken into account in realistic neutron star cooling calculations.
###### Acknowledgements.
We are grateful to D.G. Yakovlev and Yu.A. Shibanov for enlightening discussions. Special thanks go to K.P. Levenfish and O.Yu. Gnedin for assistance with the cooling code. The work was supported in part by RFBR (grant 99-02-18099), INTAS (96-0542), and KBN (2 P03D 014 13).
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# A NUMBER PROJECTED MODEL WITH GENERALIZED PAIRING INTERACTION
## I Introduction
The strongest effects associated with neutron-proton (np) pairing are expected in $`NZ`$ nuclei where valence protons and neutrons occupy the same shell-model orbits. The basic properties of the $`np`$-interaction are known from the studies of simple configurations of near closed shell nuclei . The isovector (T=1) interaction is dominated by the J=0<sup>+</sup> channel but the isoscalar (T=0) $`np`$-interaction is almost equally attractive in the J=1<sup>+</sup> and stretched J=(2$`j`$)<sup>+</sup> channels. The T=0 interaction is on the average stronger than the T=1 force. It may, therefore, lead to the appearance of a static T=0 $`np`$-pairing condensate, particularly in heavier $`NZ`$ nuclei where the large valence space allows for the creation of many np-pairs. However, it is not obvious whether these correlations are coherent enough to create this new type of collective mode nor what are the main building blocks or specific experimental fingerprints of such a condensate.
Theoretically, $`np`$-pairing is a challenging subject. It offers new opportunities to probe specific parts of the effective nucleon-nucleon interaction. The generalization of BCS (or HFB) techniques to incorporate and allow for unconstrained interplay of T=0 and T=1 pairs on equal footing is by itself non-trivial. Though the first steps to generalize the BCS theory as well as the first applications were done already in the sixties (for review of the early efforts see ) only recently the first symmetry unconstrained, self-consistent mean-field calculations have been performed. Extensions beyond mean-field, restoring either rigorously or approximately number symmetry and/or isospin symmetry are scarce.
The renaissance in the interest for $`np`$-pairing can be traced back to the fast progress in detection techniques and radioactive ion beam (RIB) programs. First experiments with RIB’s are soon to come and are targeted on heavy proton-rich nuclei in particular on $`NZ`$ nuclei. They are expected to provide important clues resolving the above mentioned, long standing difficulties in understanding $`np`$-pairing. The observables to look for are obviously those which are expected to be strongly modified by a static $`np`$-pair condensate like deuteron-transfer probability , $`\beta `$ and Gamow-Teller decay rates or ground-state and high-spin properties .
So far no clear, systematic experimental signature of the $`np`$-pairing condensate is known. There are, however, some indirect indications, for example, in recent spectroscopic data in $`{}_{36}{}^{}{}_{}{}^{72}`$Kr<sub>36</sub> and $`{}_{37}{}^{}{}_{}{}^{74}`$Rb<sub>37</sub> . In the ground state of $`{}_{37}{}^{}{}_{}{}^{74}`$Rb<sub>37</sub>, with T=1,T<sub>z</sub>=0, the $`\gamma `$ray energies of the collective $`4^+2^+0^+`$ transitions appear to be similar (isobaric analogues) to $`{}_{36}{}^{}{}_{}{}^{74}`$Kr<sub>38</sub>, the T=1,T<sub>z</sub>=1 nucleus in spite of the expected increase in the dynamical moment of inertia due to blocking of the like-particle superfluidity.<sup>*</sup><sup>*</sup>*Throughout the paper, the bold-faced symbols T and T<sub>z</sub>=(N-Z)/2 would refer to the total nuclear isospin and its z-component, respectively. The T and T<sub>z</sub> are reserved to distinguish between various two-body interaction channels. This phenomenon has been interpreted as a manifestation of T=1 $`np`$-pairing collectivity . At higher spins a transition from T=1 to T=0 band has also been observed. Calculations seem to confirm the T=1 $`np`$-collectivity at low-spins and predict an increasing role of the aligned T=0 pairs at higher spins, see Refs. . In $`{}_{36}{}^{}{}_{}{}^{72}`$Kr<sub>36</sub> a rather unexpected delay of the first crossing frequency has been measured . It again may have possible links to $`np`$-pairing, see discussion in , although more conventional explanations involving shape vibrations cannot be ruled out.
The strongest evidence for the enhancement of $`np`$-pairing effects seems to come from binding energies. The well known slope discontinuity of the isobaric mass parabola at $`NZ`$, see review and refs. therein, indicates an additional binding energy (Wigner energy) in $`NZ`$ nuclei. The Wigner energy is predominantly due to the T=0 interaction . However, the mechanism responsible for the extra binding energy seems to be rather complex when expressed in terms of $`np`$-pairs of given J,T . It cannot be solely explained in terms of J=1,T=0 $`np`$-pairs, at least not for $`sd`$ or $`pf`$ shell nuclei. A connection of the Wigner energy and the T=0 $`np`$-pairing condensate was suggested in our Letter based on deformed mean-field calculations with a schematic pairing interaction. Mass measurements of more heavy $`NZ`$ nuclei are needed to shed more light on this issue.
The aim of this paper is to further investigate basic features of $`np`$-pairing. The paper supplements the above mentioned letter explaining in more detail certain technical aspects of our model but also provides new numerical and analytical results. The paper is organized as follows: In Sect. II we introduce the basic concepts concerning the Bogolyubov transformation and the self-consistent symmetries (SCS) used here to simplify the calculations. Details concerning the model hamiltonian and implications of SCS on the structure and interpretation of the model can be found in Section III. Section IV presents the method used to restore approximatively the particle-number symmetry which is an extension of the so called Lipkin-Nogami technique for the case of a non-separable proton-neutron system. The ideas presented in this section are independent on the kind of two-body interaction used in the calculations. The results of numerical calculations, discussion and conclusions are given in Sections V and VI, respectively.
## II The Bogolyubov transformation and self-consistent symmetries
The starting point of our considerations are the eigenstates of a deformed phenomenological single-particle potential. The basis states can be divided into two groups with respect to the signature symmetry ($`\widehat{R}_x=e^{i\pi \widehat{j_x}}`$) quantum number $`r=i(+i)`$ which are later labeled as $`𝜶`$($`\stackrel{~}{𝜶}`$), respectively. Two different types of nucleonic pairs can therefore be formed, namely $`𝜶\stackrel{~}{𝜷}`$ and $`𝜶𝜷`$ pairs. A generalized BCS (gBCS) theory has to account for a scattering of these two types of nucleonic pairs simultaneously. In fact, this work is restricted to pairing in the same and signature reversed states i.e. for $`𝜶\stackrel{~}{𝜶}`$ and $`𝜶𝜶`$ modes only, see also .
The signature symmetry cannot be used as a self-consistent symmetry (SCS) in gBCS calculations. Indeed, in such a case the pairing tensor, $`𝜿\mathbf{=}𝑽^{}𝑼^T`$, connects only states of opposite signature . Consequently, the $`𝜶𝜶`$ $`np`$-pairing cannot be activated, see also . Therefore, to take into account simultaneously $`𝜶\stackrel{~}{𝜶}`$ and $`𝜶𝜶`$ pairing one needs to extend the Bogolyubov transformation. The most general Bogolyubov transformation can be written as:
$$\widehat{\alpha }_k^{}=\underset{\alpha >0}{}(U_{\alpha k}a_\alpha ^{}+V_{\stackrel{~}{\alpha }k}a_{\stackrel{~}{\alpha }}+U_{\stackrel{~}{\alpha }k}a_{\stackrel{~}{\alpha }}^{}+V_{\alpha k}a_\alpha )$$
(1)
where $`\alpha (\stackrel{~}{\alpha })`$ denote single particle states (including isospin indices) of signature $`r=i(+i)`$ respectively, while $`k`$ denotes quasiparticles. As discussed above, by superimposing any SCS one always excludes certain interaction channels in the mean-field approximation. Nevertheless, the use of SCS appear many times inherent to the nature of the physical problem and of course, make the theory more transparent and easier to handle numerically. Therefore, prior to construct a general theory involving the transformation (1) we simplify the problem by superimposing (beyond parity) the so called antilinear simplex symmetry, $`\widehat{S}_x^T=\widehat{P}\widehat{T}\widehat{R}_z`$, as SCS, see . One should bear in mind that due to the antilinearity of $`\widehat{S}_x^T`$ the transformation properties of creation and destruction operators with respect to $`\widehat{S}_x^T`$ will depend on the phases of the basis states, i.e. no new quantum number can by related directly to this symmetry.
When superimposing $`\widehat{S}_x^T`$ as the SCS it is rather convenient to choose the phase convention in such a way that the basis states will have exactly the same transformation properties with respect to both $`\widehat{R}_x`$ and $`\widehat{S}_x^T`$ and:
$$\begin{array}{ccc}& \widehat{S}_x^T\left(\begin{array}{c}a_\alpha ^{}\\ a_{\stackrel{~}{\alpha }}^{}\end{array}\right)(\widehat{S}_x^T)^1=i\left(\begin{array}{c}\hfill a_\alpha ^{}\\ \hfill a_{\stackrel{~}{\alpha }}^{}\end{array}\right)& \end{array}$$
(2)
Let us divide our quasiparticle states (1) into two families denoted as $`k`$ and $`\stackrel{~}{k}`$, respectively:
$$\begin{array}{ccc}\widehat{\alpha }_k^{}& =& \underset{\alpha >0}{}(U_{\alpha k}a_\alpha ^{}+V_{\stackrel{~}{\alpha }k}a_{\stackrel{~}{\alpha }}+U_{\stackrel{~}{\alpha }k}a_{\stackrel{~}{\alpha }}^{}+V_{\alpha k}a_\alpha )\\ \widehat{\alpha }_{\stackrel{~}{k}}^{}& =& \underset{\alpha >0}{}(U_{\stackrel{~}{\alpha }\stackrel{~}{k}}\widehat{a}_{\stackrel{~}{\alpha }}^{}+V_{\alpha \stackrel{~}{k}}\widehat{a}_\alpha +U_{\alpha \stackrel{~}{k}}\widehat{a}_\alpha ^{}+V_{\stackrel{~}{\alpha }\stackrel{~}{k}}\widehat{a}_{\stackrel{~}{\alpha }}).\end{array}$$
(3)
Enforcing $`\widehat{S}_x^T`$ symmetry as SCS requires that the quasiparticle operators of eq. (3) have the same transformation properties with respect to $`\widehat{S}_x^T`$ as the single particle operators (2. This leads to the following restrictions for the coefficients of the Bogolyubov transformation (3):
$$\left\{\begin{array}{c}U_{\alpha k}=U_{\alpha k}^{}\\ U_{\stackrel{~}{\alpha }\stackrel{~}{k}}=U_{\stackrel{~}{\alpha }\stackrel{~}{k}}^{}\\ V_{\stackrel{~}{\alpha }k}=V_{\stackrel{~}{\alpha }k}^{}\\ V_{\alpha \stackrel{~}{k}}=V_{\alpha \stackrel{~}{k}}^{}\\ 𝑹𝒆𝒂𝒍\end{array}\right\}\left\{\begin{array}{c}U_{\stackrel{~}{\alpha }k}=U_{\stackrel{~}{\alpha }k}^{}\\ U_{\alpha \stackrel{~}{k}}=U_{\alpha \stackrel{~}{k}}^{}\\ V_{\stackrel{~}{\alpha }\stackrel{~}{k}}=V_{\stackrel{~}{\alpha }\stackrel{~}{k}}^{}\\ V_{\alpha k}=V_{\alpha k}^{}\\ 𝑰𝒎𝒂𝒈𝒊𝒏𝒂𝒓𝒚\end{array}\right\}.$$
(4)
The formalism is complex but the density matrix, $`𝝆\mathbf{=}𝑽^{}𝑽^T`$, and pairing tensor $`𝜿\mathbf{=}𝑽^{}𝑼^T`$ take the relatively simple structure with real and imaginary blocks decoupled from each other:
$$𝝆=\left(\begin{array}{cc}\text{Re}(\rho _{\alpha \beta })& 0\\ 0& \text{Re}(\rho _{\stackrel{~}{\alpha }\stackrel{~}{\beta }})\end{array}\right)+i\left(\begin{array}{cc}0& \text{Im}(\rho _{\alpha \stackrel{~}{\beta }})\\ \text{Im}(\rho _{\stackrel{~}{\alpha }\beta })& 0\end{array}\right)$$
(5)
$$𝜿=\left(\begin{array}{cc}0& \text{Re}(\kappa _{\alpha \stackrel{~}{\beta }})\\ \text{Re}(\kappa _{\stackrel{~}{\alpha }\beta })& 0\end{array}\right)+i\left(\begin{array}{cc}\text{Im}(\kappa _{\alpha \beta })& 0\\ 0& \text{Im}(\kappa _{\stackrel{~}{\alpha }\stackrel{~}{\beta }})\end{array}\right)$$
(6)
Furthermore, the complex structure of the single particle potential $`𝚪`$ and the pairing potential $`𝚫`$
$$\mathrm{\Gamma }_{\alpha \beta }\underset{\gamma \delta }{}\overline{v}_{\alpha \gamma \beta \delta }\rho _{\delta \gamma }\text{and}\mathrm{\Delta }_{\alpha \beta }\frac{1}{2}\underset{\gamma \delta }{}\overline{v}_{\alpha \beta \gamma \delta }\kappa _{\gamma \delta }$$
(7)
and consequently the gBCS equations are fully determined by the $`𝝆`$ and $`𝜿`$ matrices, respectively.
In this work we define the two-body $`np`$-pairing interaction in terms of an extension of the standard seniority pairing interaction. It is separable in the particle-particle channel, $`\overline{v}_{\alpha \beta \gamma \delta }g_{\alpha \beta }g_{\gamma \delta }^{}`$, with $`g_{\alpha \beta }`$ proportional (up to a phase factor) to the overlap $`\alpha _\tau |\beta _\tau ^{}`$ between the single-particle wave functions.Subscript $`\tau `$ in $`\alpha _\tau `$ is necessary to distinguish between proton or neutron single-particle states. As already mentioned it takes essentially $`𝜶\stackrel{~}{𝜶}`$ and $`𝜶𝜶`$ types of pairing, see Sect. III for further detail. To further visualize the physical implications of the $`\widehat{S}_x^T`$ let us consider the limits of isospin ($`N=Z`$ case without Coulomb force) and time reversal symmetry (non-rotating case). Let us consider $`𝜶\overline{𝜶}`$ pairing which in principle consists of both T=1 and T=0 components of the $`np`$-pairing. By decomposing the pairing potential into the different isospin components T,T<sub>z</sub> one finds:
$$\mathrm{\Delta }_{\alpha t_z,\overline{\alpha t_z}}^{(1,0)}\kappa _{\alpha 1/2,\overline{\alpha 1/2}}+\kappa _{\alpha 1/2,\overline{\alpha 1/2}}\text{and}\mathrm{\Delta }_{\alpha t_z,\overline{\alpha t_z}}^{(0,0)}\kappa _{\alpha 1/2,\overline{\alpha 1/2}}\kappa _{\alpha 1/2,\overline{\alpha 1/2}}$$
(8)
i.e. the T=1 and T=0 components of $`np`$-pairing depend on the combinations of the same elements of the pairing tensor but with opposite sign. This is due to the phase relation for the Clebsh-Gordan coefficients which is (anti-)symmetric with respect to the interchange of a proton and neutron for T=(0)1, respectively. Time reversal symmetry further implies that :
$$\kappa _{\alpha 1/2,\overline{\alpha 1/2}}=\kappa _{\overline{\alpha 1/2},\alpha 1/2}^{}=\kappa _{\alpha 1/2,\overline{\alpha 1/2}}^{}$$
(9)
and therefore
$$\mathrm{\Delta }_{\alpha t_z,\overline{\alpha t_z}}^{(1,0)}\text{Re}(\kappa _{\alpha 1/2,\overline{\alpha 1/2}})\text{and}\mathrm{\Delta }_{\alpha t_z,\overline{\alpha t_z}}^{(0,0)}\text{Im}(\kappa _{\alpha 1/2,\overline{\alpha 1/2}})$$
(10)
Consequently, with the pairing tensor of the form of (6), the T=0 component of $`𝜶\overline{𝜶}`$ is ruled out through the $`\widehat{S}_x^T`$ symmetry. Similar analysis shows that T=1 component of the $`𝜶`$$`𝜶`$ pairing also vanishes due to symmetry reasons.
The latter is well justified due to the Pauli principle. The lack of T=0 $`𝜶\overline{𝜶}`$ is a deficiency of our model. However, for $`\mathrm{}\omega =0`$ the $`𝜶\stackrel{~}{𝜶}`$ and $`𝜶`$$`𝜶`$ $`np`$-pairing phases are in many applications indistinguishable due to time-reversal symmetry. Since our interaction is essentially structureless, based on pair counting we expect our results not to be sensitive to this restriction. For $`\mathrm{}\omega 0`$, on the other hand, one wants to associate the T=0, $`𝜶`$$`𝜶`$ ($`𝜶\stackrel{~}{𝜶}`$) $`np`$-pairing with the coupling to maximum (minimum) spin, respectively. The retained component is expected to be dominant but only at high spins. In order to probe the transition from low to high spin regime one has to explore all possible T=0 pairs and allow for an unconstrained interplay between the different pairing modes. Note, however, that again due to the simplicity of our interaction certain features of the transitional regime can be simulated to some extent with an isospin broken hamiltonian. Indeed, the missing T=0 $`𝜶\stackrel{~}{𝜶}`$ component is expected to respond to nuclear rotation in a similar way as T=1 $`𝜶\stackrel{~}{𝜶}`$. It was shown explicitely in Ref. for a single $`j`$-shell model.
In conclusion, in our model the $`𝜶\stackrel{\mathbf{~}}{𝜶}`$ pairing is equivalent to T=1 and $`𝜶𝜶`$ to T=0 and the isospin notation will be used in the following. This simple analysis reveals also the important role the self-consistent symmetries can play in theoretical description of $`np`$-pairing in the mean-field theory, see also .
## III The model Hamiltonian
The multipole-multipole expansion offers a rather good approximation to the pairing energy when neutrons and protons can be treated separably. In the following we extend this idea to the case of $`np`$-pairing by constructing a generalized pairing force separable in the particle-particle channel:
$$\widehat{V}_{pair}=\frac{1}{4}G\overline{v}_{\alpha \beta \gamma \delta }\widehat{a}_\alpha ^{}\widehat{a}_\beta ^{}\widehat{a}_\delta \widehat{a}_\gamma \frac{1}{4}\underset{\alpha \beta }{}g_{\alpha \beta }\widehat{a}_\alpha ^{}\widehat{a}_\beta ^{}\underset{\gamma \delta }{}g_{\gamma \delta }^{}\widehat{a}_\delta \widehat{a}_\gamma $$
(11)
where $`g_{\alpha \beta }\alpha |\widehat{G}|\beta `$ and $`\widehat{G}`$ is an auxiliary operator generating the specific pairing mode with strength $`G`$. The antisymmetry of the two-body matrix element $`\overline{v}_{\alpha \beta \gamma \delta }`$ implies that
$$\alpha |\widehat{G}|\beta =\beta |\widehat{G}|\alpha \alpha ,\beta .$$
(12)
and therefore, the generators $`\widehat{G}`$ must be antilinear and antihermitian. We assume here that the correlation energy of the nucleonic pair is proportional to the overlap $`\alpha _\tau |\beta _\tau ^{}`$ between single-particle states they occupy (extended seniority-type pairing interaction). The $`\widehat{G}`$ can then be chosen, for example, as:
$$\widehat{G}_{\tau \tau }=\widehat{T},\widehat{G}_{np}^{\alpha \stackrel{~}{\alpha }}=\widehat{G}_{np}^{T=1}=\frac{1}{\sqrt{2}}\widehat{\tau }_x\widehat{T},\widehat{G}_{np}^{\alpha \alpha }=\widehat{G}_{np}^{T=0}=\frac{1}{\sqrt{2}}\widehat{\tau }_y\widehat{S}_x^T$$
(13)
for pp(nn) pairing, $`𝜶\stackrel{~}{𝜶}`$ type of np-pairing and $`𝜶𝜶`$ of np-pairing. This choice is, however, not unique. In particular, it depends on the choice of the relative phases between neutron and proton states. The choice of phase convention induces strict transformation rules for the isospin Pauli operators $`\widehat{\tau }_i,i=x,y,z`$ with respect to time reversal symmetry . Constructing the generators (13) we assumed the same phases for proton and neutron states ( $`|\alpha _\tau ,r,\tau `$ denotes the basis state $`\alpha `$ of signature $`r=\pm i`$ and isospin $`\tau =1(1)`$ for neutrons(protons)):
$$\widehat{T}|\alpha _\tau ,r=i,\tau =|\alpha _\tau ,r=\pm i,\tau $$
(14)
Our phase convention further implies that:
$$\begin{array}{ccccc}\widehat{T}\widehat{\tau }_x|\alpha _\tau ,r=\pm i,\tau & =& \widehat{T}|\alpha _\tau ,r=\pm i,\tau & =& |\alpha _\tau ,r=i,\tau \\ \widehat{\tau }_x\widehat{T}|\alpha _\tau ,r=\pm i,\tau & =& \widehat{\tau }_x|\alpha _\tau ,r=i,\tau & =& |\alpha _\tau ,r=i,\tau \end{array}$$
(15)
leading to $`\widehat{T}\widehat{\tau }_x\widehat{T}^1=\widehat{\tau }_x`$. Similar considerations for $`\widehat{\tau }_y`$ and $`\widehat{\tau }_z`$ operators give:
$$\begin{array}{ccc}\widehat{T}\widehat{\tau }_x\widehat{T}^1& =& +\widehat{\tau }_x\\ \widehat{T}\widehat{\tau }_y\widehat{T}^1& =& \widehat{\tau }_y\\ \widehat{T}\widehat{\tau }_z\widehat{T}^1& =& +\widehat{\tau }_z\end{array}$$
(16)
It is straightforward to prove that, with the above phases, the operators (13) satisfy the antilinearity and antihermicity requirements formulated in Eq. (12).
As a main consequence of the separability of the pairing interaction there exists an average gap parameter:
$$\mathrm{\Delta }_{\alpha \beta }=\frac{1}{2}\underset{\gamma \delta }{}\overline{v}_{\alpha \beta \gamma \delta }\kappa _{\gamma \delta }=g_{\alpha \beta }\left\{\frac{G}{2}\underset{\gamma \delta }{}g_{\gamma \delta }^{}\kappa _{\gamma \delta }\right\}g_{\alpha \beta }\mathrm{\Delta }$$
(17)
Using the generators (13) we obtain:
$$\mathrm{\Delta }_{\alpha \tau \stackrel{~}{\beta \tau }}^{T=1}=\delta _{\alpha _\tau \beta _\tau }\mathrm{\Delta }_{\tau \tau }^{T=1}\text{where}\mathrm{\Delta }_{\tau \tau }^{T=1}=G_{\tau \tau }^{T=1}\underset{\alpha _\tau >0}{}\kappa _{\alpha _\tau \stackrel{~}{\alpha _\tau }}$$
(18)
$$\mathrm{\Delta }_{\alpha \tau \stackrel{~}{\beta \tau }}^{T=1}=\alpha _\tau |\beta _\tau \mathrm{\Delta }_{np}^{T=1}\text{where}\mathrm{\Delta }_{np}^{T=1}=\frac{1}{2}G_{np}^{T=1}\underset{\alpha _n\beta _p>0}{}\alpha _n|\beta _p\left\{\kappa _{\alpha _n\stackrel{~}{\beta _p}}+\kappa _{\beta _p\stackrel{~}{\alpha _n}}\right\}$$
(19)
$$\begin{array}{ccc}& \mathrm{\Delta }_{\alpha \tau \beta \tau }^{T=0}=i\tau \alpha _\tau |\beta _\tau \text{Im}(\mathrm{\Delta }_{np}^{T=0})\text{and}\mathrm{\Delta }_{\stackrel{~}{\alpha \tau }\stackrel{~}{\beta \tau }}^{T=0}=i\tau \alpha _\tau |\beta _\tau \text{Im}(\mathrm{\Delta }_{np}^{T=0})& \\ & \text{where}\mathrm{\Delta }_{np}^{T=0}=\frac{i}{2}G_{np}^{T=0}\underset{\alpha _n\beta _p>0}{}\alpha _n|\beta _p\{\text{Im}(\kappa _{\stackrel{~}{\alpha _n}\stackrel{~}{\beta _p}})\text{Im}(\kappa _{\alpha _n\beta _p}\})& \end{array}$$
(20)
for pp(nn) pairing, T=1, $`𝜶\stackrel{~}{𝜶}`$ type of np-pairing and T=0, $`𝜶𝜶`$ of np-pairing, respectively.
The single-particle potential $`𝒉`$ takes the following form
$$h_{\alpha \beta }=e_\alpha \delta _{\alpha \beta }\omega j_{\alpha \beta }^{(x)}+\mathrm{\Gamma }_{\alpha \beta }$$
(21)
where the single-particle energies, $`e_\alpha `$, are calculated using a deformed Woods-Saxon potential . Nuclear rotation is included using the cranking approximation and $`\mathrm{\Gamma }`$ originates from the contribution of the pairing interaction to the single-particle channel. For pp(nn) pairing we have:
$$\mathrm{\Gamma }_{\alpha _\tau \beta _\tau }^{T=1}=G_{\tau \tau }^{T=1}\rho _{\stackrel{~}{\beta _\tau }\stackrel{~}{\alpha _\tau }}\mathrm{\Gamma }_{\stackrel{~}{\alpha _\tau }\stackrel{~}{\beta _\tau }}^{T=1}=G_{\tau \tau }^{T=1}\rho _{\beta _\tau \alpha _\tau }\mathrm{\Gamma }_{\alpha _\tau \stackrel{~}{\beta _\tau }}^{T=1}=iG_{\tau \tau }^{T=1}\text{Im}(\rho _{\beta _\tau \stackrel{~}{\alpha _\tau }})$$
(22)
For T=1, $`𝜶\stackrel{~}{𝜶}`$ np-pairing we obtain:
$$\left\{\begin{array}{c}\mathrm{\Gamma }_{\alpha _\tau \beta _\tau ^{}}^{T=1}\\ \mathrm{\Gamma }_{\stackrel{~}{\alpha _\tau }\stackrel{~}{\beta _\tau ^{}}}^{T=1}\\ \mathrm{\Gamma }_{\alpha _\tau \stackrel{~}{\beta _\tau ^{}}}^{T=1}\end{array}\right\}=\frac{1}{2}G_{np}^{T=1}\underset{\gamma _\tau \delta _\tau ^{}>0}{}\alpha _\tau |\gamma _\tau \beta _\tau ^{}|\delta _\tau ^{}\left\{\begin{array}{c}\rho _{\stackrel{~}{\delta _\tau ^{}}\stackrel{~}{\gamma _\tau }}\\ \rho _{\delta _\tau ^{}\gamma _\tau }\\ i\text{Im}(\rho _{\delta _\tau ^{}\stackrel{~}{\gamma _\tau }})\end{array}\right\}$$
(23)
Finally, for T=0, $`𝜶𝜶`$ mode of np-pairing we get:
$$\left\{\begin{array}{c}\mathrm{\Gamma }_{\alpha _\tau \beta _\tau ^{}}^{T=0}\\ \mathrm{\Gamma }_{\stackrel{~}{\alpha _\tau }\stackrel{~}{\beta _\tau ^{}}}^{T=0}\\ \mathrm{\Gamma }_{\alpha _\tau \stackrel{~}{\beta _\tau ^{}}}^{T=0}\end{array}\right\}=\frac{1}{2}G_{np}^{T=0}\underset{\gamma _\tau \delta _\tau ^{}>0}{}\alpha _\tau |\gamma _\tau \beta _\tau ^{}|\delta _\tau ^{}\left\{\begin{array}{c}\rho _{\delta _\tau ^{}\gamma _\tau }\\ \rho _{\stackrel{~}{\delta _\tau ^{}}\stackrel{~}{\gamma _\tau }}\\ i\text{Im}(\rho _{\stackrel{~}{\delta _\tau ^{}}\gamma _\tau })\end{array}\right\}$$
(24)
## IV The Lipkin-Nogami method
The atomic nucleus is a mesoscopic system and as such never undergoes sharp phase transition. The fluctuations are always of importance in the transition region. The classical example is the mean-field prediction of a sharp superfluid-to-normal phase transition, induced by fast rotation, which is always smeared out in nature. The effect can be accounted for theoretically by restoring particle number and this motivated us to include approximate number-projection in our model.
In our letter we have demonstrated that number projection results in a mixing of the T=0 and T=1 pairing phases already in $`N=Z`$ nuclei. This suggests that the exclusiveness of these modes discussed in the literature is due to the mean-field approximation. For example, in the SO(8) model calculations the exact solutions allow for mixing of T=0 and T=1 phases while not the gBCS . Recently, it was shown by Goodman that the mean-field calculations with a $`G`$-matrix interaction performed in a rich model-space in fact do allow for coexistence of T=0 and T=1 pairing phases. The number-projection technique used in our calculations is known as the Lipkin-Nogami (LN) theory. Below, we will outline the main modifications necessary for applying the LN method to the non-separable proton-neutron system.
The LN theory attempts to construct the state $`|LN`$ where both linear and quadratic constraints ($`\tau \tau ^{}\{p,n\}`$):
$$LN|LN=1,LN|\mathrm{\Delta }\widehat{N}_\tau |LN=0,LN|\mathrm{\Delta }\widehat{N}_\tau \mathrm{\Delta }\widehat{N}_\tau ^{}|LN=0$$
(25)
are simultaneously fulfilled. Variation over the Lipkin-Nogami state $`|LN`$ is equivalent to a restricted HFB-type variation for the auxiliary Routhian:
$$\delta LN|\widehat{H}^\omega |LN\delta HFB|\widehat{}^\omega |HFB$$
(26)
where Obviously, $`\lambda _{np}^{(2)}=\lambda _{pn}^{(2)}`$ and the symmetric form of the auxiliary Routhian (27) is chosen for convenience.
$$\widehat{}^\omega =\widehat{H}^\omega \underset{\tau }{}\lambda _\tau ^{(1)}\mathrm{\Delta }\widehat{N}_\tau \underset{\tau \tau ^{}}{}\lambda _{\tau \tau ^{}}^{(2)}\mathrm{\Delta }\widehat{N}_\tau \mathrm{\Delta }\widehat{N}_\tau ^{}$$
(27)
The LN method is not a variational approximation. Only the standard linear constraints for the particle number are taken into account as Lagrange-type variational constraints. The parameters $`\lambda _{\tau \tau ^{}}^{(2)}`$ are kept constant during the variational procedure and eventually adjusted self-consistently using three additional subsidiary conditions<sup>§</sup><sup>§</sup>§In the following, all averages over $`|HFB`$ state will be denoted as $``$ for simplicity.:
$$\widehat{}^\omega (\mathrm{\Delta }\widehat{N}_\tau \mathrm{\Delta }\widehat{N}_\tau ^{}\mathrm{\Delta }\widehat{N}_\tau \mathrm{\Delta }\widehat{N}_\tau ^{})=0.$$
(28)
with $`\mathrm{\Delta }\widehat{N}_\tau \widehat{N}_\tau N_\tau `$.
At the point of self-consistency $`\widehat{}_{20}^\omega =0`$, and the equations (28) can be rewritten as
$$\widehat{}^\omega |44|\mathrm{\Delta }\widehat{N}_\tau \mathrm{\Delta }\widehat{N}_\tau ^{}=0$$
(29)
where $`|44|`$ denotes the projection onto the 4-quasiparticle space. The LN conditions (29) depend therefore only on the two-body residual interaction. These equations form a system of three linear equations:
$$\underset{\tau ^{\prime \prime }\tau ^{\prime \prime \prime }}{}A_{\tau \tau ^{}}^{\tau ^{\prime \prime }\tau ^{\prime \prime \prime }}x_{\tau ^{\prime \prime }\tau ^{\prime \prime \prime }}=B_{\tau \tau ^{}}$$
(30)
whereThe asymmetry is induced by the symmetric form of the routhian (27). Note, however, that the corrections to the $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ potentials due to LN terms comes out to be symmetric, see eq. (35).
$$x_{\tau \tau }=\lambda _{\tau \tau }^{(2)}\text{while}x_{pn}=2\lambda _{pn}^{(2)}$$
(31)
and
$$A_{\tau \tau ^{}}^{\tau ^{\prime \prime }\tau ^{\prime \prime \prime }}=\mathrm{\Delta }\widehat{N}_{\tau ^{\prime \prime }}\mathrm{\Delta }\widehat{N}_{\tau ^{\prime \prime \prime }}|44|\mathrm{\Delta }\widehat{N}_\tau \mathrm{\Delta }\widehat{N}_\tau ^{}$$
(32)
$$B_{\tau \tau ^{}}=\widehat{V}_{twobody}|44|\mathrm{\Delta }\widehat{N}_\tau \mathrm{\Delta }\widehat{N}_\tau ^{}$$
(33)
The $`x_{\tau \tau ^{}}`$ (or $`\lambda ^{(2)}`$) parameters are therefore equal to:
$$x_{\tau \tau ^{}}=\frac{Det𝑨^{\tau \tau ^{}}}{Det𝑨}$$
(34)
where $`𝑨^{\tau \tau ^{}}`$ denotes the matrix obtained from matrix $`𝑨`$ by replacing the $`\tau \tau ^{}`$-th column of $`𝑨`$ by the vector $`𝑩`$.
Obviously, the LN theory is technically an analog of the standard HFB theory but for the routhian (27). The resulting LN equations take therefore the form of HFB equations with the single particle field and pairing field renormalized as follows:
$$h_{\tau \tau ^{}}^{LN}h_{\tau \tau ^{}}+2\lambda _{\tau \tau ^{}}^{(2)}\rho _{\tau \tau ^{}}\text{and}\mathrm{\Delta }_{\tau \tau ^{}}^{LN}\mathrm{\Delta }_{\tau \tau ^{}}2\lambda _{\tau \tau ^{}}^{(2)}\kappa _{\tau \tau ^{}}$$
(35)
The contribution to the total energy/routhian due to the LN corrections can be written as:
$$\delta E_{LN}=\underset{\alpha \tau ,\beta \tau ^{}}{}\lambda _{\tau \tau ^{}}^{(2)}\{\rho _{\alpha \tau ,\beta \tau ^{}}(\delta _{\alpha \tau ,\beta \tau ^{}}\rho _{\beta \tau ^{},\alpha \tau })\kappa _{\alpha \tau ,\beta \tau ^{}}\kappa _{\beta \tau ^{},\alpha \tau }^{}\}$$
(36)
One should always bear in mind that the LN theory provides a sort of optimal HFB-like wave function being different then the $`|LN`$ state. Therefore the expectation values of all physical operators should be recalculated using the following formula:
$$LN|\widehat{Q}|LN\widehat{𝒬}=\widehat{Q}\underset{\tau \tau ^{}}{}\lambda _{\tau \tau ^{}}^{(2)}(\widehat{Q})\mathrm{\Delta }\widehat{N}_\tau \mathrm{\Delta }\widehat{N}_\tau ^{}$$
(37)
with $`\lambda _{\tau \tau ^{}}^{(2)}(\widehat{Q})`$ dependent on the operator $`\widehat{Q}`$ in question. These coefficients can be calculated from the following set of five equations:
$$\widehat{𝒬}\mathrm{\Delta }\widehat{N}_\tau =0,\widehat{𝒬}(\mathrm{\Delta }\widehat{N}_\tau \mathrm{\Delta }\widehat{N}_\tau ^{}\mathrm{\Delta }\widehat{N}_\tau \mathrm{\Delta }\widehat{N}_\tau ^{})=0.$$
(38)
## V Numerical calculations
An open question in mean-field calculations with $`np`$-pairing is related to the strength of the interaction. Whereas the strength of the $`pp`$\- and $`nn`$-seniority pairing force can be established by a fit to the odd-even mass differences, very little is known about the strength of the $`np`$-pairing force. Based on isospin invariance arguments, it seems well justified to assume that at the N$``$Z line $`G_{pp(nn)}^{T=1}G_{np}^{T=1}`$. Therefore, we will take in the following $`G_{np}^{T=1}=(G_{nn}^{T=1}+G_{pp}^{T=1})/2`$ with $`G_{nn(pp)}^{T=1}`$ calculated using the average gap method of Ref. . Concerning the strength of the isoscalar $`np`$-interaction we will present most of our results as a function of the parameter $`x^{T=0}`$ that scales the strength of T=0 $`np`$-pairing interaction with respect to $`G_{np}^{T=1}`$ i.e. $`x^{T=0}=G_{np}^{T=0}/G_{np}^{T=1}`$. In some cases, however, we will show examples of solutions for the hamiltonian with broken isospin symmetry.
The discussion is divided into two parts. In section V A we shall concentrate on the solutions for non-rotating cases. Subsections V A 1 and V A 2 briefly recall major properties of the solutions of both unprojected (lateron called BCS) and number-projected (lateron called BCSLN) versions of the model. In subsection V B, we will discuss a procedure allowing to estimate locally, i.e. in relatively narrow range of mass number A, a ’physical’ value of the scaling parameter $`x_0^{T=0}`$. In subsection V C we shall present the results of cranked, number projected calculations demonstrating different aspects of the interplay between mean-field, pairing forces and nuclear rotation.
### A The frequency zero case
#### 1 Basic properties of the BCS solutions
Let us first briefly recall the main properties of the BCS model solutions without number-projection . Let us start with the solutions for self-conjugated $`N=Z`$ nuclei. Disregarding the Coulomb interaction and assuming that $`G_{pp}^{T=1}=G_{nn}^{T=1}=G_{np}^{T=1}`$ leads to $`\mathrm{\Delta }_{pp}^{T=1}=\mathrm{\Delta }_{nn}^{T=1}=\mathrm{\Delta }_0`$. The solutions fall then into three classes depending on the value of $`x^{T=0}`$:
(i) For $`x^{T=0}<1`$ ($`G_{np}^{T=0}<G_{np}^{T=1}`$) the T=1 pairing is energetically favored over the T=0 pairing. The pairing energy depends only on $`\mathrm{\Delta }^22\mathrm{\Delta }_0^2+(\mathrm{\Delta }_{np}^{T=1})^2`$ and no energy is gained by activating the T=1 $`np`$-pairing. Note that in our model all pairing phases can be simultaneously present in contrast to the results presented in Ref. .
(ii) The solution at $`x^{T=0}=1`$ ($`G_{np}^{T=0}=G_{np}^{T=1}`$) is highly degenerate. The BCS energy depends only on $`\mathrm{\Delta }^22\mathrm{\Delta }_0^2+(\mathrm{\Delta }_{np}^{T=1})^2+|\mathrm{\Delta }_{np}^{T=0}|^2`$. Also in this limit, no energy is gained due to $`np`$-pairing.
(iii) The solution at $`x^{T=0}>1`$ ($`G_{np}^{T=0}>G_{np}^{T=1}`$) corresponds to a pure T=0 $`np`$-pairing phase.
For $`NZ`$ the main features of the BCS solutions can be characterized as follows:
(iv) The T=1, T<sub>z</sub>=0 $`np`$-pairing never forms a collective phase.
(v) There exists a critical value of the $`x^{T=0}`$ parameter, $`x_{crit}^{T=0}`$. For $`x^{T=0}<x_{crit}^{T=0}`$ only the T=1,$`|`$T$`{}_{z}{}^{}|`$=1 solution exist while for $`x^{T=0}x_{crit}^{T=0}`$ the T=1,$`|`$T$`{}_{z}{}^{}|`$=1 and T=0 phases are mixed.
(vi) The value of $`x_{crit}^{T=0}`$ sharply increases with increasing neutron/proton excess $`|\text{T}_z|`$. In consequence there exists a critical value $`|\text{T}_z^{crit}|`$ above which there is no collective solution for $`np`$-pairing. Our calculations show that this value is small ($`|\text{T}_z^{crit}|2`$), see next subsection.
Many of the above properties can be understood in terms of the variational approach, inherent to the (generalized) BCS-formalism. For the case of equal weight of each pair, the system will always choose only that kind of pairs that result in the lowest energy of the system, i.e. result in a sharp phase transition from T=0 to T=1 and vice versa. An interesting feature of the gBCS solution is the ’redundant’ role of the T=1 $`np`$-pairing in our model when applied to even-even nuclei. It does not mix with the other pairing phases and only in $`N=Z`$ nucleus it coexists with $`|`$T$`{}_{z}{}^{}|`$=1 pairing. Similar properties were found in Ref. . In the particular case of an $`N`$=$`Z`$ nucleus, the coexistence of T=1 pairing phases at $`x^{T=0}`$=1 is very fragile. After switching on the Coulomb force the T=1 $`np`$-pairing phase essentially vanishes. On the other hand, for odd-odd nuclei, one expects the T=1, $`|`$T$`{}_{z}{}^{}|`$=0 pairing to play a role.
The gBCS solutions for isosymmetric $`N=Z`$ nuclei have been analyzed analytically in Ref. assuming time-reversal invariance and spherical-symmetry in isospace i.e. $`\widehat{\text{T}}_i`$=0 for $`i=x,y,z`$. It has been demonstrated that coherent, $`np`$-paired, solutions can be obtained for this case provided that $`\rho _{\alpha p,\alpha p}=\rho _{\alpha n,\alpha n}`$ and $`\kappa _{\alpha p,\overline{\alpha p}}=\kappa _{\alpha n,\overline{\alpha n}}`$. The method of Ref. can be generalized to $`NZ`$ nucleus by extending symmetry condition to $`\widehat{\text{T}}_i`$=0 for $`i=x,y`$ and $`\widehat{\text{T}}_z`$=$`(NZ)/2`$. For $`NZ`$ nuclei, however, neither $`|\rho _{\alpha p,\alpha p}|=|\rho _{\alpha n,\alpha n}|`$ nor $`|\kappa _{\alpha p,\overline{\alpha }p}|=|\kappa _{\alpha n,\overline{\alpha }n}|`$. It is then rather straightforward to show, using the HFB conditions $`𝝆𝜿=𝜿𝝆^{}`$ and $`𝝆^2𝝆=𝜿𝜿^{}`$, that the gBCS theory does not allow for $`np`$ T=1 pairing solutions. For details, see Appendix A. This result relates the strong limitations of possible solutions of the generalized BCS theory to the underlying symmetries. In this particular example the limitations are due to the time-reversal symmetry and the symmetries (’deformations’) in isospace.
In our model, $`\widehat{\text{T}}_y`$=0, due to the particular form of the density matrix (5). We could verify numerically for a number of examples that the BCS solutions for an isospin invariant hamiltonian are realized at $`\widehat{\text{T}}_x`$=$`\widehat{\text{T}}_y`$=0 independent on $`x^{T=0}`$. This is not surprising since our interaction is symmetric in isospin space. For an isospin broken hamiltonian with $`G_{np}^{T=1}>G_{pp(nn)}^{T=1}`$ and with an overcritical value of $`G_{np}^{T=1}`$ the solutions for $`NZ`$ nuclei are triaxial in the isospace i.e. with $`\widehat{\text{T}}_x`$$``$0, see also discussion in Appendix A.
One may further extend the gBCS model and include an isospin cranking term, $`\mu \widehat{T}_x`$. Such a model can be considered as the lowest-order approximate isospin-projected gBCS theory. One-dimensional isospin rotations have been studied in Ref. in an exactly soluble model including T=1 pairing correlations only. It has been demonstrated that the isospin-cranking solutions rather poorly approximate the exact solutions but the particle-number projected isospin-cranking model offers a very good approximation to the exact solution. Note, that the isospin-cranking model may have several interesting analogies to the well-studied cranking approximation for spatial rotations. Due to the transformation rules of the isospin operators under time-reversal (16) which are different than those for ordinary angular momentum operators, there may also exist basic differences. Neither the possible analogies nor the differences were studied so far. The isospin cranking model in its simplest one-dimensional or more sophisticated three-dimensional version, including both T=0 and T=1 $`np`$-pairing correlations, may provide important clues on the role of $`np`$-pairing in excited configurations.
#### 2 Basic properties of the BCSLN solutions
The BCSLN solution is qualitatively similar to BCS. The major differences can be summarized as follow:
(i) Both the critical value of $`x_{crit}^{T=0}`$ as well as the ’physical’ value of $`x_0^{T=0}`$, extracted according to the prescription of subsection V B are larger in the BCSLN model than in BCS. For self-conjugated $`N=Z`$ nuclei, however, the ratio of $`x_{crit}^{T=0}/x_0^{T=0}`$ is very similar in both BCSLN and BCS models while for $`\text{T}_z0`$ this ratio is smaller in number-projected theory.
(ii) For self-conjugated $`N=Z`$ nuclei, the solutions of BCSLN allow for mixing of the T=1, $`|`$T$`{}_{z}{}^{}|`$=1 pairing phase with isoscalar $`np`$-pairing but not for the collective isovector $`np`$-pairing phase.
A shift of the critical value of $`x_{crit}^{T=0}`$ towards values that are larger than unity for $`N=Z`$ nuclei is due to the asymmetric way the Lipkin-Nogami corrections modify the different pairing channels. Whereas the ’normal’ constraints on proton and neutron number affect the normal and abnormal densities for protons and neutrons, respectively, the constraint on $`\mathrm{\Delta }NZ`$ involves an additional constrain on $`\rho _{np}`$ ($`\kappa _{np}`$). It appears (numerically) that $`\lambda _{pn}^{(2)}`$ is negative implying that the Lipkin-Nogami correction due to $`np`$-pairing, $`\delta E_{np}^{LN}`$, is positive, i.e. repulsive. Both quantities have opposite sign as compared to the analogous quantities for the like-particle channel, see Fig. 1. It means, that the effective gap parameters are weakened in the $`np`$\- and enhanced in $`nn`$\- and $`pp`$\- channels resulting in the above mentioned shift, see Eq. (35). It is interesting to note, however, that the onset of $`np`$-pairing suppresses like-particle correlations, increasing both $`\lambda _{p(n)}^{(2)}`$ and the (attractive) energy contribution $`\delta E_{p(n)}^{(LN)}`$ in such a way that the total LN correction $`\delta E_{tot}^{(LN)}`$ increases with increasing strength of the isoscalar $`np`$-interaction. Note also, that the LN corrections and $`\lambda ^{(2)}`$ parameters are strongly peaked for $`N=Z`$ nuclei and decrease rapidly with the value of $`|NZ|`$ (Fig. 1).
The $`\lambda ^{(2)}`$ parameters can be approximately estimated as:
$$\lambda _{p(n)}^{(2)}\frac{1}{2}\frac{^2E}{Z^2(N^2)}\text{and}\lambda _{pn}^{(2)}\frac{1}{2}\frac{^2E}{ZN}$$
(39)
where $`E(N,Z)`$ is the nuclear binding energy. Using the standard liquid-drop functional form for $`E(N,Z)`$ one can show that the leading order contributions to (39) for $`NZ`$ nuclei:
$$\lambda _{p(n)}^{(2)}\frac{1}{9}\frac{a_S}{A^{4/3}}+\frac{a_I}{A}+2\frac{a_W}{A}\delta _{N,Z}$$
(40)
$$\lambda _{pn}^{(2)}\frac{1}{9}\frac{a_S}{A^{4/3}}\frac{a_I}{A}2\frac{a_W}{A}\delta _{N,Z}$$
(41)
arise from the surface energy, symmetry energy and Wigner term, respectively. Obviously, these estimates are valid for complete HFBLN calculations including contributions to the $`\lambda ^{(2)}`$ parameters coming from the particle-hole channel. Nevertheless, the above estimates clearly show that (i) $`\lambda _{p(n)}^{(2)}`$ and $`\lambda _{pn}^{(2)}`$ have opposite signs and that (ii) an enhancement in both $`\lambda _{p(n)}^{(2)}`$ and $`\lambda _{pn}^{(2)}`$ is expected for $`N=Z`$ nuclei due to the singularity of the Wigner energy. These features are in nice qualitative agreement with our calculations, see Fig. 1. Whether or not this behavior is realistic; reflects deficiencies of the LN number projection scheme \[see discussion in Ref. \] or calls for isospin-projection remains to be studied.
### B Estimate of the isoscalar $`np`$-pairing strength
In our Letter we have demonstrated that the T=0 $`np`$-pairing field can yield a microscopic explanation of the Wigner energy in even-even nuclei. It naturally accounts for the singularity of nuclear masses in $`N=Z`$ nuclei. The primary mechanism leading to the Wig/-ner cusp can be viewed as a generalization of the well-known blocking effect with additional neutrons or protons outside an $`N=Z`$ core playing a similar role as the odd particle plays in the standard blocking phenomenon. This is visualized schematically in Fig. 2. It needs to be stressed that the generalized blocking is not a specific property of our schematic model but does apply for a more realistic pairing interactions as well, see discussion in . Moreover, while the standard blocking mechanism requires different trial wave functions for odd and even systems, the generalized blocking does not. The same trial wave function is used for $`N=Z`$ and $`NZ`$ nuclei in the case shown in Fig. 2.
The Wigner cusp can be obtained either by an isospin invariant model with $`G^{T=0}>G^{T=1}`$ or isospin broken model with $`G_{np}^{T=1}>G_{nn(pp)}^{T=1}`$ as shown in . We do not see, however, any particular reason to invoke an isospin-broken model, particularly in the light of the discussion in Refs. . In these works both experimental and theoretical arguments are given that the Wigner energy originates essentially from T=0 correlations. ¿From these studies it is however not clear whether it is due to $`np`$-pairing. In fact, in Ref. , the $`np`$-pairing mechanism was excluded. Note, however, that all these studies were performed using the nuclear shell-model. The shell-model Hamiltonian is usually written in the particle-particle representation but it can be rewritten in the particle-hole representation using the Pandya transformation . Therefore, in the shell-model there is no distinct division into the pairing- and single-particle field. This division is inherent to mean-field models only. Indeed, the shell-model definition of pairing in terms of $`L=0,S=1,T=1`$ and $`L=0,S=1,T=0`$ as derived from the $`G`$-matrix interaction in Ref. and used in Ref. to analyze the Wigner energy, is completely inappropriate from the point of mean-field model calculations, see also discussion in Ref. .
Ref. shows the technique to extract the strength, $`W(A)`$, of the Wigner energy \[$`E_W(A)=W(A)|NZ|`$\] from experimental data. Using this method and assuming that the Wigner energy is indeed due to the isoscalar $`np`$-pairing correlations allows us to determine a ’physical’ value of the scaling parameter $`x_0^{T=0}=G_{np}^{T=0}/G_{np}^{T=1}`$ simply by matching experimental and calculated values of $`W(A)`$. Note that this prescription is by no means limited to our schematic pairing interaction. Under the assumption that the Wigner energy is indeed due to isoscalar $`pn`$-pairing the method can be used for any pairing interaction to establish an overall scaling factor between isovector and isoscalar part.
Let us stress again at this point, that the ’physical’ value $`x_0^{T=0}`$ deduced in this subsection refers only to the frequency zero calculations because of the missing T=0 $`𝜶\overline{𝜶}`$ component. For a pairing force, that contains both $`𝜶\overline{𝜶}`$ and $`𝜶`$$`𝜶`$ components in the T=0 channel, the same prescription to determine the strength can be used. Note that also for this case the same value of $`G_{np}^{T=0}`$ will be obtained. Indeed, due to time reversal invariance, the BCS or BCSLN energy gain caused by the T=0 pairing will depend only on the modulus $`|\stackrel{}{\mathrm{\Delta }}^{T=0}|`$ of the total isoscalar gap, $`\stackrel{}{\mathrm{\Delta }}^{T=0}(\mathrm{\Delta }_{\alpha \alpha }^{T=0},\mathrm{\Delta }_{\alpha \overline{\alpha }}^{T=0}`$), but not on its direction (see discussion in Section V A).
Fig. 3 illustrates calculations for $`pf`$-shell $`A48`$ nuclei. In the calculations we have taken a constant deformation of $`\beta _2=0.25`$ and included 30 deformed neutron and proton states. The realistic values of $`x_0^{T=0}`$ are estimated to be $``$1.13 and $``$1.30 for BCS and BCSLN models, respectively, as shown in the figure. We performed similar BCSLN calculations for $`A76`$ nuclei assuming $`\beta _2=0.40`$ and taking 40 deformed neutron and proton states. In this case we have taken $`W(A)=47/A`$ MeV as a reference value because experimental masses are not available. The calculations yield $`x_0^{T=0}1.25`$ which would imply that $`G^{T=0}`$ has a different mass dependence than $`G^{T=1}`$.
In this context one should mention, that the deduced 1/$`A`$ dependence of $`W(A)`$ is essentially due to the fast increase of $`W(A)`$ for very light $`sd`$-shell nuclei. For $`sd`$-shell nuclei, however, we expect the T=0 pairing to have a pronounced $`L=0,S=1`$ component. This component becomes strongly quenched due to the spin-orbit interaction when entering the $`pf`$-shell. The properties of T=0 pairing may change strongly at the border of the $`pf`$ and $`sd`$-shells. The 1/$`A`$ dependence of $`W(A)`$ may therefore even be an artefact. With the present set of data one cannot extract a reliable $`A`$ dependence of the Wigner energy based only on $`pf`$-shell and heavier nuclei. No doubt, that data on masses along the $`N=Z`$ line are urgently required in order to resolve this important question.
The procedure to determine the strength is very sensitive to even small variations in $`x^{T=0}`$ particularly for BCS calculations because we are working in the region of the phase transition. Indeed, changing $`x_{|\text{BCS}}^{T=0}`$ by $`2\%`$ (in $`A48`$ nuclei) affects the Wigner energy by $`200`$ keV, as demonstrated in Fig. 3. Moreover, although $`x_{0|\text{BCSLN}}^{T=0}>x_{0|\text{BCS}}^{T=0}`$ in $`N=Z`$, $`A48`$ nuclei, the ratio $`yx_0^{T=0}/x_{crit}^{T=0}`$ appears to be very similar in both BCS and BCSLN models. However, for $`\text{T}_z0`$ we get $`y_{\text{BCS}}<y_{\text{BCSLN}}`$ as shown in Fig. 4. In consequence $`|\text{T}_{z|\text{BCSLN}}^{crit}|>|\text{T}_{z|\text{BCS}}^{crit}|`$ as expected. The difference is still rather small and $`|\text{T}_z^{crit}|2`$ independently on the model. Fig. 5 compares calculated (BCSLN model) mass excess $`\mathrm{\Delta }E=E(x^{T=0})E(x^{T=0}=1)`$ for Cr and Sr isotopes as a function of $`\text{T}_z`$. No pronounced differences are seen between Cr and Sr isotopes although there is clear tendency for $`\mathrm{\Delta }E(\text{T}_z)`$ to be more spread in heavier nuclei, see also Fig. 4.
The Wigner energy has an additional repulsive component in odd-odd $`N=Z`$ nuclei. In it has been demonstrated that the strength of this component $`d(A)_{T=0}W(A)`$ i.e. both components of the Wigner energy have, most likely, the same origin. This term can be understood within our pairing theory provided that the odd-spin, T=0 states in $`N=Z`$ odd-odd nuclei are not treated on the same footing like even-even nuclei but are interpreted as two-quasiparticle (2$`qp`$) configurations involving one broken $`np`$-pair. Indeed, treating even-even and odd-odd nuclei on the same footing within the generalized BCS(LN) approach gives no odd-odd versus even-even mass staggering, see Fig. 6. The blocking of a T=0 np quasiparticle, on the other side, would result in an odd-odd versus even-even mass (T=0) staggering in $`N=Z`$ nuclei due to the T=0 pairing energy, analogous to the standard odd-even mass staggering caused by $`T=1`$ pairing. Although the T=0 ground state in odd-odd nuclei is of $`2qp`$ character, it is still an open problem of how to construct a proper trail wave function appropriate for odd-odd nuclei. It seems that the isospin restoration (see discussion in Ref. ) may play a key role here.
### C The cranking calculations
As mentioned previously, our schematic, separable $`np`$-interaction does not allow for T=0 $`𝜶\stackrel{~}{𝜶}`$ $`np`$-pairing. In a condensate, that is dominated by T=0 $`𝜶`$$`𝜶`$ pairs, there is little resistance for these pairs to align their angular momenta. Note that this alignment occurs without invoking any pair-breaking mechanism. The alignment is a smooth function of the frequency, and no backbend occurs, see Figure 3 in our letter .
In the low spin regime, this is of course an unphysical behaviour. The presence of T=0 $`𝜶\stackrel{~}{𝜶}`$ pairing will invoke the pair-breaking mechanism (low-$`J`$ $`np`$-pairs dominate at low-spin ). However, as long as we do not use any physical measure sensitive to the isospin of the $`np`$-pairs the T=0 $`𝜶\stackrel{~}{𝜶}`$ $`np`$-pairing can be mocked up by T=1 $`𝜶\stackrel{~}{𝜶}`$ $`np`$-pairing and vice-versa via a simple readjustment of the strength (isospin-broken model). Indeed, both modes will respond to nuclear rotation in similar way, i.e. through the pair breaking mechanism which leads to a suppression of this field at high spins.
In contrast to the $`𝜶\stackrel{~}{𝜶}`$ pairing the $`𝜶`$$`𝜶`$ type pairing will become enhanced at high rotational frequencies as Coriolis and centrifugal forces align more and more pairs along the axis of collective rotation. Indeed, as shown in , the T=0 phase composed of high-$`J`$ $`np`$-pairs survives to much higher angular momentum. This is in accordance to recent shell model Monte-Carlo calculations . At very high spins, our schematic T=0 $`𝜶`$$`𝜶`$ pairing is probably somewhat too strong since there is no restriction in the $`J`$-values of the pairs that participate in the scattering. However, in accordance to early calculations, we expect to pick up the main essence of the high spin behaviour.
Fig. 7 shows an illustrative example phase-diagram of the critical frequency (the frequency corresponding to the onset of $`np`$-pairing of $`𝜶`$$`𝜶`$) versus the strength of this force. The calculations were performed using the BCSLN model for the $`N=Z`$ nucleus <sup>48</sup>Cr at constant deformation and for $`\mathrm{}\omega 2`$ MeV. Filled circles represent calculations with $`G_{pp}=G_{nn}=G_{np}^{\alpha \stackrel{~}{\alpha }}`$ while open circles correspond to an isospin symmetry broken Hamiltonian with $`G_{np}^{\alpha \stackrel{~}{\alpha }}=1.3G_{pp(nn)}`$. Two things are worth to be noticed: (i) $`\mathrm{}\omega _{crit}`$ is very sensitive to the strength parameter and (ii) the alignment (in this case the alignment of $`f_{7/2}`$ quasiparticles) triggers the onset of T=0 $`np`$-pairing setting a natural threshold ($`x_t^{\alpha \alpha }`$) for the process. Above this threshold i.e. for $`x^{\alpha \alpha }<x_t^{\alpha \alpha }`$ the value of $`\mathrm{}\omega _{crit}`$ increases very sharply. Therefore, even the optimistic scenario would involve the onset of $`np`$-pairing (if at all possible) around the crossing frequency i.e. in the region which is anyhow the most difficult to describe theoretically. This is in accordance to single-$`j`$ shell model calculations .
### D Terminating states
In the $`pf`$-shell, the state of the art $`0\mathrm{}\omega `$ shell-models are able to describe nuclear structure with excellent accuracy up to the maximum spin states (terminating states) which is built upon the pure $`f_{7/2}`$ sub-shell, see e.g. . For example, for the case of <sup>48</sup>Cr, the state of maximum spin of $`I=16^+\mathrm{}`$ can be reached within $`\pi f_{7/2}^4\nu f_{7/2}^4`$ configuration. Particle-hole excitations to $`p_{3/2}`$ and $`f_{5/2}`$ sub-shells allow further to built states of spin up to $`I=20^+\mathrm{}`$ within the full $`pf`$ space. However, already at these spins the states involving cross-shell particle-hole excitations are expected to compete. Unfortunately, with the present-day shell-model techniques it is not feasible to perform calculations within the model space needed to describe states above spin $`I=16\mathrm{}`$, which in principle should include the whole $`sd`$-shell, $`pf`$-shell and the $`g_{9/2}`$ sub-shell.
In the standard cranked mean-field model calculations the collectivity of the rotational band in <sup>48</sup>Cr is essentially exhausted after the simultaneous alignment of $`f_{7/2}`$ neutrons and protons and one enters the non-collective, unpaired regime around $`I_x=16\mathrm{}`$. Building up higher angular momentum states proceed further in discrete jumps, which are due to crossings between single-particle down-sloping and up-sloping routhians at spherical shape. This reoccupation process will finally lead again to the onset of collectivity.
The presence of T=0 correlations at high angular momenta may strongly alter the behaviour of high spin states discussed above. In our model, above the terminating $`I=16\mathrm{}`$ state we observe smooth increase of $`I_x`$ versus $`\mathrm{}\omega `$ rather than the step-like process expected in standard scenario, see Fig. 8. This is due to the pair scattering from $`d_{3/2}`$ and $`f_{7/2}`$ into the aligned $`g_{9/2}`$ and $`f_{5/2}`$ orbits which is entirely due to T=0 pairing. Partial occupation of these orbits triggers the onset of collectivity after the terminating state . New experiments targeted to measure the evolution of the rotational band in <sup>48</sup>Cr beyond $`I=16\mathrm{}`$ may therefore shed light into the nature of the mean-field T=0 pairing correlations at high spins. Particularly, the possibility of the T=0 pairing correlations to affect the electromagnetic transition rates. Indeed, without T=0 correlations, the wave functions of high-spin states will have the dominant component of 2p-2h and/or 4p-4h configuration (due to promotion of a $`\pi f_{7/2}\nu f_{7/2}`$ pair or $`\pi f_{7/2}^2\nu f_{7/2}^2`$ quartet to higher sub-shells). The E2-decay from such a configuration into the 0p-0h ($`\pi f_{7/2}^4\nu f_{7/2}^4`$) state would then be strongly hindered. In contrast, T=0 pairing will cause mixing of these (multi)particle-(multi)hole configurations (collectivity) enhancing the E2 transition amplitudes.
### E TRS-calculations at superdeformed shape
Nuclear superdeformation (SD) is an extreme case of the spontaneous breaking of spherical symmetry. In this case our T=0 pairing interaction, scattering uniformly $`np`$-pairs, may even be considered as quite a realistic approximation provided that the strength is properly adjusted. The standard total routhian surface (TRS) as well as Skyrme-Hartree-Fock calculations predict SD bands to appear at relatively low excitation energy in $`NZ`$ nuclei of A$``$80-90. Predicted deformations are $`\beta _20.55`$ for <sup>88</sup>Ru and even as large as $`\beta _20.75`$ for <sup>88</sup>Mo, see . To illustrate the influence of T=0 interaction on nuclear superdeformation we performed TRS calculations for SD <sup>88</sup>Ru. Three variants of the calculations different in the treatment of pairing have been performed: (i) unpaired single-particle (ii) including isovector T=1 pairing only, (iii) including both T=0 and T=1 pairing with $`x^{T=0}=0.9`$, $`x^{T=0}=1.1`$, and $`x^{T=0}=1.3`$. The dynamical moments of inertia ($`J^{(2)}dI_x/d\omega `$) versus rotational frequency resulting from this set of calculation are shown in Fig. 9.
As expected, the calculations (i) and (ii) differ strongly at low spins but converge at higher frequencies where static T=1 pairing correlations are essentially washed out by the Coriolis force. The pronounced kink at $`\mathrm{}\omega 0.75`$ MeV in calculations (ii) is due to the simultaneous alignment of the $`h_{11/2}`$ protons and neutrons. Choosing the strongly undercritical strength of $`G^{T=0}`$, corresponding to $`x^{T=0}=0.9`$, results of course in a pair field dominated by T=1 pairs. However, after the $`h_{11/2}`$ alignment and the disappearance of static T=1 pairing correlations, the T=0 correlations start to build up, resulting in a somewhat larger moment of inertia as compared to the unpaired case. Clearly, the onset of T=0 pairing is triggered by the alignment. The $`x^{T=0}=1.1`$ corresponds to a slightly undercritical $`G^{T=0}`$ strength at low $`\mathrm{}\omega `$. The low frequency hump in $`J^{(2)}`$ at $`\mathrm{}\omega 0.3`$ MeV is in this case due to the rotation-induced phase transition from predominantly T=1 to T=0 pairing . For larger values of $`x^{T=0}`$ the hump moves towards lower frequencies in accordance with the phase diagram in Fig. 7 discussed in Subsect. V C, and the $`h_{11/2}`$ alignment is smoothed out. Moreover, the high-frequency part of $`J^{(2)}`$ increases exceeding the ’rigid body’ value We use the notation of ’rigid body’ value as the one given by calculations without pairing. Note that this moment of inertia strongly differs from that of a rigid body, which would be a constant. by up to $``$30%. Although the increase depends very much on the strength of the T=0 interaction (and essentially vanishes for $`x^{T=0}`$$``$0.8) the observation of systematic deviations from the standard T=1 model predictions can be considered as a rather robust indicator of the importance of the stretched i.e. $`𝜶`$$`𝜶`$ type T=0 correlations at high spins.
## VI Summary and conclusions
Our work presents a number-projected, cranked mean-field formalism appropriate for the simultaneous treatment of T=1 and T=0 correlations. The number-projection technique is a generalization of the Lipkin-Nogami method. We use a schematic $`np`$-pairing interaction because our interest is to study general properties of $`np`$-paired solutions rather than to reproduce properties of specific nuclei.
Our work shows that the exclusiveness of different pairing phases is inherent to the BCS method. Already approximate number-projection leads to a mixing of the T=0 and T=1 pairing phases. Common for both approaches, though, is that different pairing phases counteract each other.
It is also demonstrated that neutron or proton excess block $`np`$-pairing. This generalized blocking mechanism causes $`np`$-pairing to vanish rapidly when departing from the $`N=Z`$ line. The calculations suggest that already at $`|𝑻_z|2`$ there is essentially no chance to observe any collective $`np`$-pairing. This mechanism serves as a possible microscopic explanation of the Wigner energy at the level of mean-field theory. The even-even versus odd-odd mass staggering in T=0 states of $`N=Z`$ nuclei indicates a $`2qp`$ structure (one broken $`np`$-pair) of the T=0 states in odd-odd $`N=Z`$ nuclei. This will require blocked calculations which are beyond the scope of our present paper.
Our model predicts that the T=0 correlations do not vanish at high spin, in accordance to previous calculations. In the high-spins regime stretched T=0 pairing is predicted to be dominant. The onset of this pairing phase seems to be triggered by the alignment of high-$`j`$ quasiparticles and the related quenching of isovector pairing correlations. The presence of T=0 pairing at high spins is predicted to increase the moments of inertia as well as to affect the evolution of rotational bands beyond their standard terminating states. The size of this effect depends critically on the strength of the T=0 force in the $`𝜶`$$`𝜶`$ channel. Below a certain limit, in our case below $`x_{\alpha ,\alpha }^{T=0}`$$``$0.8, the effect will not be seen.
This research was supported in part by the U.S. Department of Energy under Contract Nos. DE-FG02-96ER40963 (University of Tennessee), DE-FG05-87ER40361 (Joint Institute for Heavy Ion Research), DE-AC05-96OR22464 with Lockheed Martin Energy Research Corp. (Oak Ridge National Laboratory), by the Polish Committee for Scientific Research (KBN) under Contract No. 2 P03B 040 14, and by the Swedish Institute.
## A The time-reversal invariance, isospace deformations and the BCS solutions.
The time-reversal invariance implies the following relations for the hermitian density matrix and antisymmetric pairing tensor elements :
$$\left\{\begin{array}{c}\rho _{\alpha \tau ,\alpha \tau }=\rho _{\alpha \tau ,\alpha \tau }^{}\\ \rho _{\alpha \tau ,\alpha \tau ^{}}=\rho _{\overline{\alpha }\tau ,\overline{\alpha }\tau ^{}}^{}\\ \rho _{\alpha \tau ,\overline{\alpha }\tau }=0\\ \rho _{\alpha \tau ,\overline{\alpha }\tau }=\rho _{\overline{\alpha }\tau ,\alpha \tau }^{}\end{array}\right\}\left\{\begin{array}{c}\kappa _{\alpha \tau ,\overline{\alpha }\tau }=\kappa _{\alpha \tau ,\overline{\alpha }\tau }^{}\\ \kappa _{\alpha \tau ,\alpha \tau }=\kappa _{\overline{\alpha }\tau ,\overline{\alpha }\tau }^{}\\ \kappa _{\alpha \tau ,\alpha \tau }=0\\ \kappa _{\alpha \tau ,\overline{\alpha }\pm \tau }=\kappa _{\overline{\alpha }\tau ,\alpha \pm \tau }^{}\end{array}\right\}.$$
(A1)
The standard BCS constraints $`\mathrm{\Delta }\widehat{N}=\mathrm{\Delta }\widehat{Z}=0`$ automatically imply that $`\widehat{\text{T}}_z=(NZ)/2`$. Superimposing further symmetry constraint $`\widehat{\text{T}}_x=\widehat{\text{T}}_y=0`$ together with time-reversal symmetry induced relations (A1) implies:
$$\widehat{\text{T}}_x=\underset{\alpha >0}{}\text{Re}(\rho _{\alpha p,\alpha n})=\widehat{\text{T}}_y=\underset{\alpha >0}{}\text{Im}(\rho _{\alpha p,\alpha n})=0\rho _{\alpha \tau ,\alpha \tau }=0\alpha .$$
(A2)
The HFB solutions further demand that the generalized density matrix:
$$\mathbf{}=\left(\begin{array}{cc}𝝆& 𝜿\\ 𝜿^{}& \mathrm{𝟏}𝝆^{}\end{array}\right)$$
(A3)
obeys the idempotency condition $`\mathbf{}^2=\mathbf{}`$ or, equivalently, satisfies the following relations: $`𝝆𝜿=𝜿𝝆^{}`$ and $`𝝆^2𝝆=𝜿𝜿^{}`$. These relations impose six constraints on the density matrix and pairing tensor elements which all must be fulfilled simultaneously. They can be written as:
$$\kappa _{\alpha \tau ,\alpha \tau }\rho _{\alpha \tau ,\alpha \tau }+\kappa _{\alpha \tau ,\overline{\alpha }\tau }\rho _{\alpha \tau ,\overline{\alpha }\tau }=0$$
(A4)
$$\rho _{\alpha \tau ,\alpha \tau }\text{Re}(\kappa _{\alpha \tau ,\overline{\alpha }\tau })+i\text{Im}(\kappa _{\alpha \tau ,\alpha \tau }\rho _{\alpha \tau ,\overline{\alpha }\tau }^{})=0$$
(A5)
$$\kappa _{\alpha \tau ,\alpha \tau }(\rho _{\alpha \tau ,\alpha \tau }\rho _{\alpha \tau ,\alpha \tau })=\rho _{\alpha \tau ,\overline{\alpha }\tau }(\kappa _{\alpha \tau ,\overline{\alpha }\tau }\kappa _{\alpha \tau ,\overline{\alpha }\tau })$$
(A6)
$$\kappa _{\alpha \tau ,\overline{\alpha }\tau }(\rho _{\alpha \tau ,\alpha \tau }\rho _{\alpha \tau ,\alpha \tau })=\rho _{\alpha \tau ,\alpha \tau }(\kappa _{\alpha \tau ,\overline{\alpha }\tau }\kappa _{\alpha \tau ,\overline{\alpha }\tau })$$
(A7)
$$\rho _{\alpha \tau ,\alpha \tau }^2+|\rho _{\alpha \tau ,\alpha \tau }|^2+|\rho _{\alpha \tau ,\overline{\alpha }\tau }|^2\rho _{\alpha \tau ,\alpha \tau }=\kappa _{\alpha \tau ,\overline{\alpha }\tau }^2|\kappa _{\alpha \tau ,\alpha \tau }|^2|\kappa _{\alpha \tau ,\overline{\alpha }\tau }|^2$$
(A8)
$$\rho _{\alpha \tau ,\alpha \tau }(1\rho _{\alpha \tau ,\alpha \tau }\rho _{\alpha \tau ,\alpha \tau })=\kappa _{\alpha \tau ,\overline{\alpha }\tau }(\kappa _{\alpha \tau ,\overline{\alpha }\tau }+\kappa _{\alpha \tau ,\overline{\alpha }\tau })$$
(A9)
$$\rho _{\alpha \tau ,\overline{\alpha }\tau }(1\rho _{\alpha \tau ,\alpha \tau }\rho _{\alpha \tau ,\alpha \tau })=\kappa _{\alpha \tau ,\alpha \tau }(\kappa _{\alpha \tau ,\overline{\alpha }\tau }+\kappa _{\alpha \tau ,\overline{\alpha }\tau }).$$
(A10)
For $`N=Z`$ nuclei, the coherent $`np`$-paired solution can be obtained when $`\rho _{\alpha p,\alpha p}=\rho _{\alpha n,\alpha n}`$ and $`\kappa _{\alpha p,\overline{\alpha }p}=\kappa _{\alpha n,\overline{\alpha }n}`$. This is the solution found by Goodman and we refer reader to the Ref. for further details. For $`NZ`$ nucleus however, neither $`|\rho _{\alpha p,\alpha p}|=|\rho _{\alpha n,\alpha n}|`$ nor $`|\kappa _{\alpha p,\overline{\alpha }p}|=|\kappa _{\alpha n,\overline{\alpha }n}|`$. Using the relations (A4)-(A10) it is straightforward to show that the additional symmetry (axiality in isospace) (A2) rules out $`np`$-pairing of $`𝜶\overline{𝜶}`$ type. Indeed, Eqs. (A7) and (A9) give instantly $`\kappa _{\alpha \tau \overline{\alpha }\tau }=0`$.
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# 1 Introduction
## 1 Introduction
The grand unified model (GUT) of the strong and electroweak interactions is motivated not only from an esthetic point of view, but also due to the problem of the electric charge quantization in the standard model (SM) . The electric charge quantization is experimentally proved up to $`10^{21}`$ . Also, the unification of the gauge coupling constants, which is one of the GUT predictions, is shown to be valid in the supersymmetric extension of the GUT (SUSY GUT), and it is found that the GUT scale ($`M_{\mathrm{GUT}}`$) is about $`10^{16}`$ GeV .
The GUT has two aspects. One is the unification of the SM gauge groups to a simple group, such as SU(5) or SO(10). The charge quantization and the gauge coupling constant unification come from the unification of gauge groups. The other is unification of matters. In the SU(5) GUT, quarks and leptons are embedded in an economical and elegant way as
$`\psi (\mathrm{𝟏𝟎})=\left(\begin{array}{ccc}Q,& u^𝒞,& e^𝒞\end{array}\right),`$ (2)
$`\varphi (\mathrm{𝟓}^{})=\left(\begin{array}{cc}d^𝒞,& L\end{array}\right),`$ (4)
where $`Q(u,d)`$ and $`L(\nu ,e)`$. Here, we adopt the left-handed basis for fermions and omit the chirality indices of fermions as far as it does not lead to confusion in this article. In the SO(10) GUT the right-handed neutrinos are introduced and the embedding becomes simpler as
$`\phi (\mathrm{𝟏𝟔})=\left(\begin{array}{cccc}u^𝒞,& d^𝒞,& e^𝒞,& \nu ^𝒞\\ u,& d,& e,& \nu \end{array}\right).`$ (7)
The atmospheric neutrino observation by the Super-Kamiokande experiment suggests that the tau neutrino has a finite mass ($`m_{\nu _\tau }10^{(12)}`$eV). If the tiny neutrino mass comes from the see-saw mechanism , the right-handed tau neutrino mass is expected to be ($`10^2`$GeV)<sup>2</sup>/$`10^{(12)}`$eV$`10^{15}`$GeV. It is close to the GUT scale and supports the SO(10) SUSY GUT.
The unification of matters leads to violation of the global symmetries in the SM, such as baryon, lepton, and lepton flavor numbers, and then, proton decay and lepton-flavor violating processes are predicted. In this article we will concentrate on the proton decay in the SUSY SU(5) GUT. The proton decay is the direct prediction for the model, and the search is the most important for confirmation of the GUT. For the lepton-flavor violating processes in the SUSY GUT, see Refs. .
There are two sources to induce the proton decay in the SUSY SU(5) GUT. One is the $`X`$ boson, accompanied with the unification of gauge groups. The proton lifetime predicted by the $`X`$ boson exchange is proportional to $`M_X^4`$ with $`M_X`$ the $`X`$ boson mass, since the $`X`$ boson exchange induces the dimension-six operators. The dominant decay mode is $`p\pi ^0e^+`$. While this prediction is almost model-independent, the lifetime is sensitive to $`M_X`$. If $`M_X10^{16}`$GeV, the lifetime is $`10^{(3436)}`$ years, which is beyond the current experimental reach.
The second one is the colored Higgs, which is introduced for doublet Higgs in the SUSY SM to be embedded into the SU(5) multiplets. In the minimal SUSY SU(5) GUT, the colored Higgs predicts much shorter proton lifetime. The colored Higgs exchange leads to the baryon-number violating dimension-five operators, and then the proton lifetime is proportional to $`M_{H_C}^2`$ with $`M_{H_C}`$ the colored Higgs mass . While it is suppressed by the small Yukawa coupling of quarks and leptons, the minimal SUSY SU(5) GUT is excluded from the negative search of the proton decay .
This paper is organized as follows. In the next section, we show the status of the minimal SUSY SU(5) GUT from a point of view of the proton decay induced by the colored Higgs, and discuss the reason why the proton decay is not still discovered if the SUSY SU(5) GUT is valid. In Section 3 we show the proton decay rate induced by the $`X`$ boson, and discuss the model-dependence of the prediction. Section 4 is devoted to conclusion.
## 2 Proton decay induced by the colored Higgs exchange
In this section we show the status of the minimal SUSY SU(5) GUT from a point of view of the proton decay induced by the colored Higgs, and discuss the reason why the proton decay is not discovered if the SUSY SU(5) GUT is realistic.
First, let us introduce the minimal SUSY SU(5) GUT. In this model the doublet Higgs $`H_f`$ and $`\overline{H}_f`$ in the SUSY SM are embedded into the 5 and 5 dimensional multiplets as
$`H=(H_C,H_f),`$ $`\overline{H}=(\overline{H}_C,\overline{H}_f),`$ (8)
with the colored Higgs $`H_C`$ and $`\overline{H}_C`$. The superpotential of the Yukawa coupling for the doublet and colored Higgs is given as
$`W_Y`$ $`=`$ $`h^i(Q_iH_f)u_i^𝒞+V_{ij}^{}f^j(Q_i\overline{H}_f)d_j^𝒞+f^ie_i^𝒞(L_i\overline{H}_f)`$ (9)
$`+{\displaystyle \frac{1}{2}}h^ie^{i\phi _i}(Q_iQ_i)H_C+V_{ij}^{}f^j(Q_iL_j)\overline{H}_C`$
$`+h^iV_{ij}u_i^𝒞e_j^𝒞H_C+e^{i\phi _i}V_{ij}^{}f^ju_i^𝒞d_j^𝒞\overline{H}_C,`$
where the indices $`i`$ and $`j(=13)`$ are for the generations. $`V_{ij}`$ is the Kobayashi-Maskawa matrix element at the GUT scale, and $`\phi _i`$ are for additional degrees of freedom in the Yukawa coupling constants in the minimal SUSY SU(5) GUT. (See Ref. for notation and conversion.) Here, we use the same letters for the superfields as the component fields. The first-three terms correspond to the Yukawa coupling in the SUSY SM. As well-known, the charged leptons and the down-type quarks have common Yukawa coupling constants.
In the minimal SUSY SU(5) GUT $`H_C`$ and $`\overline{H}_C`$ have a common mass term by themselves, and then, the colored Higgs exchange gives following dimension-five operators,
$$W_5=\frac{1}{2M_{H_C}}h^ie^{i\phi _i}V_{kl}^{}f^l(Q_iQ_i)(Q_kL_l)+\frac{1}{M_{H_C}}h^iV_{ij}e^{i\phi _k}V_{kl}^{}f^lu_i^𝒞e_j^𝒞u_k^𝒞d_l^𝒞.$$
(10)
Since squarks or sleptons are on the external lines of these operators, Higgsino or gaugino is exchanged between them so that proton can decay (Fig. (1)).
The contractions of the indices for the gauge symmetries in Eq. (10) are understood as
$`(Q_iQ_i)(Q_k\dot{L}_l)`$ $`=`$ $`\epsilon _{\alpha \beta \gamma }(u_i^\alpha d_i^\beta d_i^\alpha u_i^\beta )(u_k^\gamma e_ld_k^\gamma \nu _l),`$ (11)
$`u_i^𝒞e_j^𝒞u_k^𝒞d_l^𝒞`$ $`=`$ $`\epsilon ^{\alpha \beta \gamma }u_{i\alpha }^𝒞e_j^𝒞u_{k\beta }^𝒞d_{l\gamma }^𝒞,`$ (12)
with $`\alpha `$, $`\beta `$, and $`\gamma `$ being color indices. Note that the total antisymmetry in the color indices requires that the operators are flavor non-diagonal $`(ik)`$. Therefore the dominant decay modes tend to have strangeness. In the minimal SUSY SU(5) GUT the dominant mode is $`pK^+\overline{\nu }`$ while $`pK^0\mu ^+`$ is not. The mode of $`K^0\mu ^+`$ is suppressed by $`(m_u/m_cV_{ud}V_{cd})^2`$.
The largest uncertainty to the decay rate predicted by the colored Higgs exchange comes from the colored Higgs mass itself in the minimal SUSY SU(5) GUT. However, the mass spectrum at the GUT scale can be constrained and the colored Higgs mass can be evaluated from the low energy parameters . Due to the gauge coupling unification in the minimal SUSY SU(5) GUT we can derive the following relations from the renormalization group equations of the gauge coupling constants at one-loop level,
$`(2\alpha _3^1+3\alpha _2^1{\displaystyle \frac{3}{5}}\alpha _Y^1)(m_Z)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\left\{{\displaystyle \frac{12}{5}}\mathrm{ln}{\displaystyle \frac{M_{H_C}}{m_Z}}2\mathrm{ln}{\displaystyle \frac{m_{SUSY}}{m_Z}}\right\},`$ (13)
$`(2\alpha _3^13\alpha _2^1+3\alpha _Y^1)(m_Z)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\left\{12\mathrm{ln}{\displaystyle \frac{M_X^2M_\mathrm{\Sigma }}{m_Z^3}}+8\mathrm{ln}{\displaystyle \frac{m_{SUSY}}{m_Z}}\right\}.`$ (14)
Here, we take the SUSY particle masses in the SUSY SM to be common ($`m_{SUSY}`$). $`M_\mathrm{\Sigma }`$ is the mass for the SU(3)<sub>C</sub> octet and SU(2)<sub>L</sub> triplet components in the 24-dimensional Higgs for SU(5) gauge symmetry breaking, which survive as the physical degrees of freedom after the Higgs mechanism. In Ref. the upperbound on $`M_{H_C}`$ is derived from Eq. (13), including the various correction to it, as
$$M_{H_C}2.4\times 10^{16}\mathrm{GeV}.$$
(15)
On the other hand, the recent study for the proton decay shows
$$M_{H_C}7.2\times 10^{16}\mathrm{GeV}$$
(16)
from the experimental bound on the proton lifetime $`\tau (pK^+\overline{\nu })>6.7\times 10^{32}`$ years , assuming the squark and slepton masses are lighter than 1TeV. Then, the minimal SUSY SU(5) GUT is excluded now.
While the minimal SUSY SU(5) GUT is excluded, the prediction of the proton decay induced by the colored Higgs exchange is itself highly model-dependent. Let us discuss the model-dependence. First, though it depends on the Yukawa coupling structure at the GUT scale, we cannot explain the mass relations between the down-type quarks and the charge leptons in the first and second generations in an SU(5) symmetric way as Eq. (9). We need to add some modification to it. One of the examples is introduction of the higher dimensional operators suppressed by the gravitational scale. If the Yukawa coupling of the colored Higgs is suppressed automatically or accidentally, the lowerbound on $`M_{H_C}`$ becomes looser. In this case, the mode of $`K^0\mu ^+`$ may be dominant. One of the explicit models is given in Ref. .
Second, in some models the effective mass for the colored Higgs, which suppressed the dimension-five operators, can be much heavier than the GUT scale . The dimension-five operators given in Eq. (10) are induced since $`H_C`$ and $`\overline{H}_C`$ have a common mass term. If $`H_C`$ and $`\overline{H}_C`$ have mass terms independent of each others by introduction of $`H_C^{}`$ and $`\overline{H}_C^{}`$, the dimension-five operators are not induced.
If the Peccei-Quinn (PQ) symmetry is introduced, this mechanism works . The PQ charges for matters and Higgs are assigned as
$`Q_{\mathrm{PQ}}(\psi _i)=1,`$ $`Q_{\mathrm{PQ}}(\varphi _i)=1,`$
$`Q_{\mathrm{PQ}}(H)=2,`$ $`Q_{\mathrm{PQ}}(\overline{H})=2,`$
$`Q_{\mathrm{PQ}}(H^{})=2,`$ $`Q_{\mathrm{PQ}}(\overline{H}^{})=2,`$
where $`H^{}((H_C^{},H_f^{}))`$ and $`\overline{H}^{}((\overline{H}_C^{},\overline{H}_f^{}))`$. This symmetry prohibits the mass term $`M_{H_C}\overline{H}_CH_C`$ and the dimension-five operators in Eq. (10). After the PQ symmetry is broken at $`M_{\mathrm{PQ}}10^{(1012)}`$ GeV, the dimension-five operators are generated. However, they are suppressed by $`M_{\mathrm{PQ}}/M_{\mathrm{GUT}}10^{(46)}`$ and they are harmless.
In order to confirm the SUSY GUT, we need more model-independent predictions, which come from the gauge sector. One of them is the proton decay induced by the $`X`$ boson. Then, we would like to discuss the model-dependence in next section.
## 3 Proton decay induced by the $`X`$ boson exchange
In this section we show the proton decay rate induced by the $`X`$ boson exchange, and discuss the model-dependence. The $`X`$ boson, is SU(3)<sub>C</sub> $`\mathrm{𝟑}^{}`$\- and SU(2)<sub>L</sub> 2\- dimensional, and the hypercharge is $`\mathrm{𝟓}/\mathrm{𝟔}`$. The interaction of the $`X`$ boson to matter fermions is given as follows;
$``$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}g_5V_{ij}\overline{d_j^𝒞}(X_\mu \gamma ^\mu L_i)+{\displaystyle \frac{1}{\sqrt{2}}}g_5\mathrm{e}^{i\phi _i}\overline{Q_i}X_\mu \gamma ^\mu u_i^𝒞`$ (17)
$`+{\displaystyle \frac{1}{\sqrt{2}}}g_5V_{ij}(\overline{e_j^𝒞}X_\mu )\gamma ^\mu Q_i+h.c..`$
This interaction leads to the following baryon-number violating operators, which contribute to $`p\pi ^0e^+`$,
$`_{eff}`$ $`=`$ $`A_R{\displaystyle \frac{g_5^2}{M_X^2}}\mathrm{e}^{i\phi _1}\times `$ (18)
$`ϵ_{\alpha \beta \gamma }((\overline{d_L^𝒞})^\alpha (u_R)^\beta (\overline{u_R^𝒞})^\gamma (e_L)+(1+|V_{ud}|^2)(\overline{d_R^𝒞})^\alpha (u_L)^\beta (\overline{u_L^𝒞})^\gamma (e_R))`$
where $`A_R`$ is the renormalization factor from the anomalous dimensions to these operators.
The renormalization factor $`A_R`$ has the short-distance contribution ($`A_R^{(SD)}`$) and the long-distance contribution ($`A_R^{(LD)}`$). $`A_R^{(SD)}`$ at one-loop level is given as
$`A_R^{(SD)}`$ $`=`$ $`\left({\displaystyle \frac{\alpha _3(m_Z)}{\alpha _5}}\right)^{\frac{4}{3b_3}}\left({\displaystyle \frac{\alpha _2(m_Z)}{\alpha _5}}\right)^{\frac{3}{2b_2}}`$ (19)
where $`b_3=92n_g`$ and $`b_2=52n_g`$ with $`n_g`$ the number of the generations . Here, we do not include the U(1)<sub>Y</sub> contribution. It is expected to be smaller while it has not been calculated completely. If the SUSY SM is the low energy effective theory below the GUT scale, $`A_R^{(SD)}=2.1`$ for $`\alpha _3(m_Z)=0.116`$, $`\mathrm{sin}^2\theta _W=0.2317`$, and $`\alpha ^1(m_Z)=127.9`$. The long-distance part $`A_R^{(LD)}`$ is
$`A_R^{(LD)}`$ $`=`$ $`\left({\displaystyle \frac{\alpha _3(m_b)}{\alpha _3(m_Z)}}\right)^{\frac{6}{23}}\left({\displaystyle \frac{\alpha _3(\mu _{had})}{\alpha _3(m_b)}}\right)^{\frac{6}{25}}=1.2.`$ (20)
From the effective lagrangian (18) the proton lifetime from $`p\pi ^0e^+`$ is
$`\mathrm{\Gamma }(p\pi ^0e^+)`$ $`=`$ $`\alpha _H^2{\displaystyle \frac{m_p}{64\pi f_\pi ^2}}(1+D+F)^2\left({\displaystyle \frac{g_5^2}{M_X^2}}A_R\right)^2(1+(1+|V_{ud}|^2)^2),`$ (21)
where $`m_p`$ is the proton mass, $`f_\pi `$ is the pion decay constant, and $`D`$ and $`F`$ are the chiral lagrangian parameters. The definition of $`\alpha _H`$ is
$`\alpha _Hu_L(𝐤)`$ $``$ $`ϵ_{\alpha \beta \gamma }0|(\overline{d_L^𝒞})^\alpha (u_R)^\beta (\overline{u_R^𝒞})^\gamma |p_𝐤.`$ (22)
$`\alpha _H`$ had one-order ambiguity in the old theoretical calculations and it led to a large uncertainty to the prediction. The latest evaluation by lattice calculation gives
$`\alpha _H`$ $`=`$ $`(0.015\pm 1)\mathrm{GeV}^3.`$ (23)
They adopt the naive dimensional renormalization scheme with the renormalization scale ($`\mu _{had}`$) 2.3GeV. This lattice calculation is the first realistic one, in which enough a large lattice spacing and a large statistics are prepared. While this calculation still has ambiguities from the quenched approximation and the $`a^1`$ correction, they are expected to be at most $`O(10)`$%. Finally, we get
$`1/\mathrm{\Gamma }(p\pi ^0e^+)`$ $`=`$ $`1.0\times 10^{35}\mathrm{years}`$ (24)
$`\times \left({\displaystyle \frac{\alpha _H}{0.015\mathrm{GeV}^3}}\right)^2\left({\displaystyle \frac{\alpha _5}{1/25}}\right)^2\left({\displaystyle \frac{A_R}{2.5}}\right)^2\left({\displaystyle \frac{M_X}{10^{16}\mathrm{GeV}}}\right)^4.`$
In Fig. (2) we show the lifetime of proton as a function of $`M_X`$. The shaded region has been excluded by Super-Kamiokande experiment. The dash-doted line is the reach of ten years run of the Super-Kamiokande experiment. In Fig. (2) we show the lowerbound on the $`X`$ boson mass in the minimal SUSY SU(5) GUT as a reference point ,
$$M_X1.1\times 10^{16}\mathrm{GeV}.$$
(25)
This bound comes from the constraint from the gauge coupling unification given in Eq. (14),
$$1.3\times 10^{16}\mathrm{GeV}(M_X^2M_\mathrm{\Sigma })^{\frac{1}{3}}3.2\times 10^{16}\mathrm{GeV}$$
(26)
and validity of perturbation below the gravitational scale,
$$M_\mathrm{\Sigma }/M_X=\lambda _\mathrm{\Sigma }/2\sqrt{2}g_51.8.$$
(27)
$`\lambda _\mathrm{\Sigma }`$ the self-coupling constant of the 24-dimensional Higgs multiplet, and if $`\lambda _\mathrm{\Sigma }`$ is much larger than $`g_5`$, it blows up below the gravitational scale. Similarly, we can derive the lowerbound on $`M_X`$ in the Modified Missing Doublet (MMD) model as $`M_X8.7\times 10^{15}`$GeV . From this figure, it is found that we need further effort to confirm or reject the SUSY GUT.
We have assumed that the SUSY SM is valid below the GUT scale. However, even if extra SU(5) complete multiplets have smaller masses than the GUT scale, the success of the gauge coupling unification is not spoiled. In fact, some models, such as the gauge mediated SUSY breaking model , the anomalous U(1) SUSY breaking model given in Ref. , and the E<sub>6</sub> model , predict existence of the extra matters at lower energy. In this case the the gauge coupling constant at the GUT scale becomes larger. Also, the renormalization factor $`A_R`$ is enhanced as Eq. (19). Then, for fixed $`M_X`$, the proton lifetime becomes shorter.
In Fig. (3) we show the proton lifetime as a function of contribution to the beta functions of the gauge coupling constants from the extra matters. The contribution from a pair of 5 and 5 is one, and that from a pair of 10 and 10 is three. Here, we fix $`M_X=1.0\times 10^{16}`$ GeV. and take the extra matter mass $`M`$ to be $`10^2`$, $`10^4`$, and $`10^6`$ GeV. From this figure, introduction of two extra generations (a pair of 5 and 5 and a pair of 10 and 10) with the masses $`10^2`$ GeV is slightly disfavored.
## 4 Conclusion
The proton decay search is the most important for confirmation of the GUT. Now the minimal SUSY SU(5) GUT is excluded by the negative search since the dimension-five operators induced by the charged Higgs exchange predict a large proton decay rate. However, the proton decay by the colored Higgs is highly model-dependent, and it is premature to conclude that the SUSY GUT is excluded. We hope that the search will be pushed as far as possible.
## Acknowledgement
We appreciate T. Onogi and N. Tsutsui for useful comments.
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# BACK REACTION OF COSMOLOGICAL PERTURBATIONS 11footnote 1BROWN-HET-1212, Jan. 2000; invited lecture at COSMO-99 (ICTP, Trieste, Sept. 27 - Oct. 2, 1999), to appear in the proceedings (World Scientific, Singapore, 2000).
## 1 Introduction
It is well known that gravitational waves propagating in some background space-time affect the dynamics of the background. This back-reaction can be described in terms of an effective energy-momentum tensor $`\tau _{\mu \nu }`$. In the short wave limit, when the typical wavelength of the waves is small compared with the curvature of the background space-time, $`\tau _{\mu \nu }`$ has the form of a radiative fluid with an equation of state $`p=\rho /3`$ (where $`p`$ and $`\rho `$ denote pressure and energy density, respectively).
In most models of the early Universe, scalar-type metric perturbations are more important than gravity waves. Here, we report on a study of the back reaction problem for both scalar and tensor gravitational perturbations . We derive the effective energy-momentum tensor $`\tau _{\mu \nu }`$ which describes the back-reaction and apply the results to inflationary cosmology.
A crucial issue to be addressed when studying scalar-type metric perturbations is the problem of gauge artifacts. As is well known (see e.g. Mukhanov et al. for a comprehensive review), cosmological perturbations transform non-trivially under coordinate transformations (gauge transformations). However, the answer to the question “how important are perturbations for the evolution of a background” must be independent of the choice of gauge, and hence the back-reaction problem must be formulated in a gauge invariant way. In recent work , we demonstrated that the back reaction problem can be set up in a way which is gauge-invariant (under linear coordinate transformations). This summary will not dwell on the issue of gauge-invariance, and the reader is referred to the original articles for details.
Here, we will briefly summarize the formalism and derive the effective energy-momentum tensor $`\tau _{\mu \nu }`$ which describes the back-reaction. We apply the results to a simple model of chaotic inflation and discuss the equation of state corresponding to the resulting $`\tau _{\mu \nu }`$. As we shall see, long wavelength fluctuations dominate, and in the limit where the short wavelength fluctuations can be neglected, the resulting equation of state of $`\tau _{\mu \nu }`$ is $`p\rho `$ with $`\rho <0`$, i.e. the equation of state corresponding to a negative cosmological constant. We speculate on the possible connection of this result with a dynamical relaxation mechanism for the cosmological constant (see also Tsamis & Woodard for related work).
## 2 Gravitational Back-Reaction
The analysis of gravitational back-reaction is related to early work by Brill, Hartle and Isaacson , among others. The idea is to expand the Einstein equations to second order in the perturbations, to assume that the first order terms satisfy the equations of motion for linearized cosmological perturbations (hence these terms cancel), to take the spatial average of the remaining terms, and to regard the resulting equations as equations for a new homogeneous metric $`g_{\mu \nu }^{(0,br)}`$ which includes the effect of the perturbations to quadratic order:
$$G_{\mu \nu }(g_{\alpha \beta }^{(0,br)})=\mathrm{\hspace{0.17em}8}\pi G\left[T_{\mu \nu }^{(0)}+\tau _{\mu \nu }\right]$$
(1)
where the effective energy-momentum tensor $`\tau _{\mu \nu }`$ of gravitational back-reaction contains the terms resulting from spatial averaging of the second order metric and matter perturbations:
$$\tau _{\mu \nu }=<T_{\mu \nu }^{(2)}\frac{1}{8\pi G}G_{\mu \nu }^{(2)}>,$$
(2)
where pointed brackets stand for spatial averaging, and the superscripts indicate the order in perturbations.
As formulated in (1) and (2), the back-reaction problem is not independent of the coordinatization of space-time and hence is not well defined. It is possible to take a homogeneous and isotropic space-time, choose different coordinates, and obtain a non-vanishing $`\tau _{\mu \nu }`$. This “gauge” problem is related to the fact that in the above prescription, the hypersurface over which the average is taken depends on the choice of coordinates.
The key to resolving the gauge problem is to realize that to second order in perturbations, the background variables change. This change can be calculated consistently, and given this change of background quantities it can be shown that the back-reaction problem is formulated in a covariant way by (1) and (2). A gauge independent form of the back-reaction equation (1) can be derived by defining background and perturbation variables $`Q=Q^{(0)}+\delta Q`$ which do not change under linear coordinate transformations. Here, $`Q`$ represents collectively both metric and matter variables. The gauge-invariant perturbation quantities are Bardeen’s gauge-invariant variables . The gauge-invariant form of the back-reaction equation then looks formally identical to (1), except that all variables are replaced by the corresponding gauge-invariant ones. We will follow the notation of , and use as gauge-invariant perturbation variables the Bardeen potentials $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ which in longitudinal gauge coincide with the actual metric perturbations $`\delta g_{\mu \nu }`$. Calculations hence simplify greatly if we work directly in longitudinal gauge. Recently, the back-reaction analysis of has been confirmed by working in a completely different gauge, making use of the covariant quantization approach.
For simplicity, we shall take matter to be described in terms of a single scalar field. In this case, there is only one independent metric perturbation variable, and in longitudinal gauge the perturbed metric can be written in the form
$$ds^2=(1+2\varphi )dt^2a(t)^2(12\varphi )\delta _{ij}dx^idx^j,$$
(3)
where $`a(t)`$ is the cosmological scale factor. The energy-momentum tensor for a scalar field is
$$T_{\mu \nu }=\phi _{,\mu }\phi _{,\nu }g_{\mu \nu }\left[\frac{1}{2}\phi ^{,\alpha }\phi _{,\alpha }V(\phi )\right].$$
(4)
By expanding the Einstein tensor and the above energy-momentum tensor to second order in the metric and matter fluctuations $`\varphi `$ and $`\delta \phi `$, respectively, it can be shown that the non-vanishing components of the effective back-reaction energy-momentum tensor $`\tau _{\mu \nu }`$ become
$`\tau _{00}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi G}}\left[+12H\varphi \dot{\varphi }3(\dot{\varphi })^2+9a^2(\varphi )^2\right]`$ (5)
$`+`$ $`{\displaystyle \frac{1}{2}}(\delta \dot{\phi })^2+{\displaystyle \frac{1}{2}}a^2(\delta \phi )^2`$
$`+`$ $`{\displaystyle \frac{1}{2}}V^{\prime \prime }(\phi _0)\delta \phi ^2+2V^{}(\phi _0)\varphi \delta \phi ,`$
and
$`\tau _{ij}`$ $`=`$ $`a^2\delta _{ij}\{{\displaystyle \frac{1}{8\pi G}}[(24H^2+16\dot{H})\varphi ^2+24H\dot{\varphi }\varphi `$ (6)
$`+`$ $`(\dot{\varphi })^2+4\varphi \ddot{\varphi }{\displaystyle \frac{4}{3}}a^2(\varphi )^2]+4\dot{\phi _0}^2\varphi ^2`$
$`+`$ $`{\displaystyle \frac{1}{2}}(\delta \dot{\phi })^2{\displaystyle \frac{1}{6}}a^2(\delta \phi )^24\dot{\phi _0}\delta \dot{\phi }\varphi `$
$``$ $`{\displaystyle \frac{1}{2}}V^{\prime \prime }(\phi _0)\delta \phi ^2+2V^{}(\phi _0)\varphi \delta \phi \},`$
where $`H`$ is the Hubble expansion rate.
## 3 Equation of State of Back-Reaction in Inflationary Cosmology
The metric and matter fluctuation variables $`\varphi `$ and $`\delta \phi `$ are linked via the Einstein constraint equations, and hence all terms in the above formulas for the components of $`\tau _{\mu \nu }`$ can be expressed in terms of two point functions of $`\varphi `$ and its derivatives. The two point functions, in turn, are obtained by integrating over all of the Fourier modes of $`\varphi `$, e.g.
$$\varphi ^2_{k_i}^{k_u}𝑑kk^2|\varphi _k|^2,$$
(7)
where $`\varphi _k`$ denotes the amplitude of the k’th Fourier mode. The above expression is divergent both in the infrared and in the ultraviolet. The ultraviolet divergence is the usual divergence of a free quantum field theory and can be “cured” by introducing an ultraviolet cutoff. In the infrared, we will discard all modes with wavelength larger than the Hubble radius at the beginning of inflation, since these modes are determined by the pre-inflationary physics. We take these modes to contribute to the background.
At any time $`t`$ we can separate the integral in (7) into the contribution of infrared and ultraviolet modes, the separation being defined by setting the physical wavelength equal to the Hubble radius. Thus, in an inflationary Universe the infrared phase space is continually increasing since comoving modes are stretched beyond the Hubble radius, while the ultraviolet phase space is either constant (if the ultraviolet cutoff corresponds to a fixed physical wavelength), or decreasing (if the ultraviolet cutoff corresponds to fixed comoving wavelength). In either case, unless the spectrum of the initial fluctuations is extremely blue, two point functions such as (7) will at later stages of an inflationary Universe be completely dominated by the infrared sector. In the following, we will therefore restrict our attention to this sector, i.e. to wavelengths larger than the Hubble radius.
In order to evaluate the two point functions which enter into the expressions for $`\tau _{\mu \nu }`$, we need to know the time evolution of the linear fluctuations $`\varphi _k`$, which is given by the linear theory of cosmological perturbations . On scales larger than the Hubble radius, and for a time-independent equation of state, the well-known result for $`\varphi _k`$ is
$$\varphi _k(t)\mathrm{const}.$$
(8)
The Einstein constraint equations yield a relation between the metric $`\varphi `$ and the matter fluctuations and $`\delta \phi `$.
$$\dot{\varphi }+H\varphi =4\pi G\dot{\phi }_0\delta \phi .$$
(9)
In most models of inflation, exponential expansion of the Universe results because $`\phi _0`$ is rolling slowly, i.e.
$$\dot{\phi }_0\frac{V^{}}{3H},$$
(10)
where a prime denotes the derivative with respect to the scalar matter field. Making use of (8), we can combine Eqs. (9) and (10) to obtain
$$\delta \phi =\frac{2V}{V^{}}\varphi .$$
(11)
Hence, in the expressions (5) and (6) for $`\tau _{\mu \nu }`$, all terms with space and time derivatives can be neglected, and we obtain
$$\rho _{br}\tau _0^0\left(2\frac{V^{\prime \prime }V^2}{V^{}^2}4V\right)<\varphi ^2>$$
(12)
and
$$p_{br}\frac{1}{3}\tau _i^i\rho _{br},$$
(13)
The main result which emerges from this analysis is that the equation of state of the dominant infrared contribution to the energy-momentum tensor $`\tau _{\mu \nu }`$ which describes back-reaction takes the form of a negative cosmological constant
$$p_{br}=\rho _{br}\mathrm{with}\rho _{br}<0.$$
(14)
The second crucial result is that the magnitude of $`\rho _{br}`$ increases as a function of time. This is due firstly to the fact that, in an inflationary Universe, as time increases more and more wavelengths become longer than the Hubble radius and begin to contribute to $`\rho _{br}`$. A second reason for the growth of the absolute value of $`\rho _{br}`$ is that the amplitude of the individual modes $`\varphi _k`$ changes as a consequence of the evolution and slow change in the equation of state of the background, as governed by the “conservation law” $`\zeta =\mathrm{const}`$ where
$$\zeta =\frac{2}{3}\frac{H^1\dot{\varphi }+\varphi }{1+w}+\varphi ,$$
(15)
where $`w=p/\rho `$. In the models studied, the effect of the growth of the individual fluctuation modes is larger than the effect due to the increase in the phase space of infrared modes.
## 4 Application to Chaotic Inflation
To study the magnitude of back-reaction, we will consider a single field chaotic inflation model with potential
$$V(\phi )=\frac{1}{2}m^2\phi ^2.$$
(16)
Furthermore, we specify an initial state at a time $`t_i`$ in which the homogeneous inflaton field has the value $`\phi _0(t_i)`$ and the fluctuations are minimal.
Using the values of $`\varphi _k`$ on long wavelengths which are well known , the (dominant) infrared contribution to the correlator becomes
$$\varphi ^2(t)=_{k_i}^{k_t}\frac{dk}{k}|\varphi _k|^2=\frac{m^2M_P^2}{32\pi ^4\phi _0^4(t)}_{k_i}^{k_t}\frac{dk}{k}\left[\mathrm{ln}\frac{H(t)a(t)}{k}\right]^2,$$
(17)
where $`M_P`$ is the Planck mass, $`k_i`$ corresponds to the physical infrared cutoff, and $`k_t`$ is the time-dependent comoving wavelength corresponding to the Hubble radius. Since the background evolution is known, the integral over $`k`$ can be done. The resulting back-reaction energy density $`\rho _{br}`$ can hence be calculated and compared with the background density $`\rho _0`$. The result is
$$\frac{\rho _{br}(t)}{\rho _0}\frac{3}{4\pi }\frac{m^2\phi _0^2(t_i)}{M_P^4}\left[\frac{\phi _0(t_i)}{\phi _0(t)}\right]^4.$$
(18)
As follows from (18), back-reaction may lead to a shortening of the period of inflation. Without back-reaction, inflation would end when $`\phi _0(t)M_P`$. Inserting this value into (18), we see that if
$$\phi _0(t_i)>\phi _{br}m^{1/3}M_P^{4/3},$$
(19)
then back-reaction will become important before the end of inflation and may shorten the period of inflation. It is interesting to compare this value with the scale
$$\phi _0(t_i)\phi _{sr}=m^{1/2}M_P^{3/2},$$
(20)
which emerges in the scenario of stochastic chaotic inflation as the “self-reproduction” scale beyond which quantum fluctuations dominate the evolution of $`\phi _0(t)`$. Notice that $`\phi _{sr}\phi _{br}`$ since $`mM_P`$. Hence, even in the case when self-reproduction does not take place, back-reaction effects can be very important.
## 5 Speculations
Since the back-reaction of cosmological fluctuations in an inflationary cosmology acts (see (14)) like a negative cosmological constant, and since the magnitude of the back-reaction effect increases in time, one may speculate that back-reaction will lead to a dynamical relaxation of the cosmological constant (see Tsamis & Woodard for similar considerations based on the back-reaction of long wavelength gravitational waves).
The background metric $`g_{\mu \nu }^{(0,br)}`$ including back-reaction evolves as if the cosmological constant at time $`t`$ were
$$\mathrm{\Lambda }_{\mathrm{eff}}(t)=\mathrm{\Lambda }_0+8\pi G\rho _{br}(t)$$
(21)
and not the bare cosmological constant $`\mathrm{\Lambda }_0`$. Hence we propose to identify (21) with a time dependent effective cosmological constant. Since $`|\rho _{br}(t)|`$ increases as $`t`$ grows, the effective cosmological constant will decay. Note that even if the initial magnitude of the perturbations is small, eventually (if inflation lasts a sufficiently long time) the back-reaction effect will become large enough to cancel any bare cosmological constant.
Furthermore, we speculate that this dynamical relaxation mechanism for $`\mathrm{\Lambda }`$ will be self-regulating. As long as $`\mathrm{\Lambda }_{\mathrm{eff}}(t)>\mathrm{\hspace{0.17em}8}\pi G\rho _m(t)`$, where $`\rho _m(t)`$ stands for the energy density in ordinary matter and radiation, the evolution of $`g_{\mu \nu }^{(0,br)}`$ is dominated by $`\mathrm{\Lambda }_{\mathrm{eff}}(t)`$. Hence, the Universe will be undergoing accelerated expansion, more scales will be leaving the Hubble radius and the magnitude of the back-reaction term will increase. However, once $`\mathrm{\Lambda }_{\mathrm{eff}}(t)`$ falls below $`\rho _m(t)`$, the background will start to decelerate, scales will enter the Hubble radius, and the number of modes contributing to the back-reaction will decrease, thus reducing the strength of back-reaction. Hence, it is likely that there will be a scaling solution to the effective equation of motion for $`\mathrm{\Lambda }_{\mathrm{eff}}(t)`$ of the form
$$\mathrm{\Lambda }_{\mathrm{eff}}(t)\mathrm{\hspace{0.17em}8}\pi G\rho _m(t).$$
(22)
Such a scaling solution would correspond to a contribution to the relative closure density of $`\mathrm{\Omega }_\mathrm{\Lambda }1`$.
## 6 Discussion
We have summarized a new formalism to study the effect of linear cosmological perturbations on the cosmological background. This effect is expressed in terms of an effective energy-momentum tensor. We have shown that the back-reaction can be described in a way which is invariant under linear coordinate transformations. The issue of invariance under higher order transformations is crucial but still unresolved . The most interesting result which emerges is that, in an inflationary background, the effective energy-momentum tensor which describes the back-reaction has the form of a negative cosmological constant. The absolute value of the induced effective energy density grows in time and, in a model with a long period of inflation, can become significant, which leads to the speculation that the effect may lead to a dynamical relaxation of the cosmological constant. However, the effective energy-momentum tensor defined in this work describes the effect of fluctuations on the homogeneous mode of the gravitational field. If the speculations in the previous section are to hold up, the analysis must be extended to give back-reaction effects on local quantities. Work on this issue is in progress.
## Acknowledgments
I thank my collaborators Raul Abramo and Slava Mukhanov for a very stimulating collaboration (they should not be blamed if the speculations of the final section do not hold up!). I also thank A. Guth, W. Unruh and R. Woodard for many lively discussions. This work is supported in part by DOE Contract DE-FG0291ER40688, Task A, and by the US NSF collaborative research award NSF-INT-9312335.
## References
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# Untitled Document
GALAXY SURFACE PHOTOMETRY
Bo Milvang-Jensen<sup>1</sup><sup>,</sup><sup>2</sup> and Inger Jørgensen<sup>3</sup> <sup>1</sup>Copenhagen University Observatory, 2100 Copenhagen Ø, Denmark, milvang@astro.ku.dk <sup>2</sup>School of Physics and Astronomy, University of Nottingham, University Park, NG7 2RD Nottingham, UK (postal address for BMJ) <sup>3</sup>Gemini Observatory, 670 N. Aohoku Pl., Hilo, Hawaii 96720, USA, ijorgensen@gemini.edu Received March 3, 2000
Abstract. We describe galaxy surface photometry based on fitting ellipses to the isophotes of the galaxies. Example galaxies with different isophotal shapes are used to illustrate the process, including how the deviations from elliptical isophotes are quantified using Fourier expansions. We show how the definitions of the Fourier coefficients employed by different authors are linked. As examples of applications of surface photometry we discuss the determination of the relative disk luminosities and the inclinations for E and S0 galaxies. We also describe the color-magnitude and color-color relations. When using both near-infrared and optical photometry, the age–metallicity degeneracy may be broken. Finally we discuss the Fundamental Plane where surface photometry is combined with spectroscopy. It is shown how the FP can be used as a sensitive tool to study galaxy evolution.
Key words: techniques: photometric – galaxies: photometry – galaxies: fundamental parameters – galaxies: elliptical and lenticular, cD – galaxies: stellar content – galaxies: evolution
1. INTRODUCTION
Surface photometry of galaxies is a technique to quantitatively describe the light distribution of the galaxies, as recorded in 2-dimensional images. This paper focuses on the techniques used to derive the surface photometry and presents some examples of scientific applications. The paper is a summary of the lectures given at the summer school, and is not intended as a complete review of the topic. Surface photometry of galaxies has previously been reviewed by Kormendy & Djorgovski (1989) and Okamura (1988).
The techniques and the software used for performing surface photometry of galaxies are described in Section 2. Surface photometry has many applications. In Section 3 we give an example of such an application, namely the determination of disk luminosities and inclinations of E and S0 galaxies. This determination is based on surface photometry – ellipticities and deviations from elliptical shape – alone.
From surface photometry in several passbands we can derive the colors of the galaxies. These colors provide information about the ages and the metal content of the stellar populations in the galaxies. In order to study this, stellar population models are needed. These models are the subject of Section 4. In Section 5, we then present examples of the color-magnitude and color-color relations of galaxies.
Although the topic of this summer school is photometry, we will nevertheless show an example of the science that can be carried out when surface photometry is combined with spectroscopy. The chosen example is the relation known as the Fundamental Plane for E and S0 galaxies, which we discuss in Section 6.
We end this description of surface photometry with a few suggestions for future projects that can be carried out based on surface photometry only, see Section 7.
Throughout the paper we use $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$.
2. SURFACE PHOTOMETRY
Surface photometry is used to study extended objects, such as galaxies, as opposed to point sources, such as stars. From an image of a point source, only its total magnitude can be derived. From an image of a galaxy, it is possible to determine a number of quantities. Some of these quantities are derived from surface photometry, such as how the intensity and ellipticity vary with radius. Other quantities are determined by other methods – the morphological type, for example, is determined from visual inspection of the image. There also exist schemes to do automated morphological classification, e.g. Abraham et al. (1994, 1996); Naim et al. (1995).
In this context, we are talking about galaxies where the individual stars cannot be resolved in the images. This is for example the case for the HydraI (Abell 1060) cluster, which is a nearby cluster at a distance of $`80`$ Mpc. At that distance, even the Hubble Space Telescope (HST) cannot resolve individual stars in the galaxies.
Photometry of crowded stellar fields, such as globular clusters, is not considered surface photometry. Surface photometry is used when the magnitudes of the individual stars (typically in a galaxy) cannot be measured, but only their smooth integrated light.
2.1 Ellipse fitting
Surface photometry of galaxies is usually done by fitting ellipses to the isophotes. This choice is motivated by that fact that the isophotes of galaxies are not far from ellipses. This is especially the case for elliptical (E) and lenticular (S0) galaxies. In this paper we will concentrate on E and S0 galaxies. We will also limit our discussion to data obtained with CCDs (Charge-Coupled Devices).
There exist several software packages for deriving surface photometry. For the exercise related to the lectures given at the summer school we have chosen the ELLIPSE task in the ISOPHOTE package (see Busko 1996). This task is based on the ellipse fitting algorithm used in the GASP package by Cawson (1983; see also Davis et al. 1985), the code for which was later rewritten by Jedrzejewski (1987). ISOPHOTE is a part of the external IRAF IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. See also http://iraf.noao.edu/ package STSDAS<sup>∗∗</sup> STSDAS is distributed by the Space Telescope Science Institute, operated by AURA, Inc., under NASA contract NAS 5-26555. See also http://ra.stsci.edu/STSDAS.html. We chose this software for the exercise since it is publicly available, and because it includes some documentation. The examples of surface photometry presented in this section are also based on the ISOPHOTE package. Our aim is to describe the basic principles of surface photometry based on ellipse fitting, and the choice of software package is not critical for that purpose. A comparison of the results obtained with ELLIPSE with those obtained with other software packages is beyond the scope of this paper.
To illustrate how surface photometry on E and S0 galaxies is derived we have chosen three example galaxies, see Figure 1 and Table 1. We will refer to these galaxies as the ‘pure E’, the ‘S0’ and the ‘boxy E’, respectively. The name ‘pure E’ refers to the fact that the isophotes of this galaxy are almost perfect ellipses. For the ‘boxy E’, the isophotes are box shaped. The pure E has been morphologically classified as E3/S0. The boxy E has been morphologically classified as SB(rs)0(0) (barred S0 with rings). For our purpose of illustrating surface photometry these morphological types are not important.
Fig. 1. The three example galaxies – pure E (top), S0 (left), boxy E (right). The width of the montage is 32 kpc. North is down and east is to the right.
| Table 1. | The three example galaxies | | | |
| --- | --- | --- | --- | --- |
| Description | Name | Morph. type | $`M_{\mathrm{r}_\mathrm{T}}`$ | $`r_\mathrm{e}`$ |
| Pure E | R347 / IC2597 | E3/S0 | $`23.35`$ mag | 8.9 kpc |
| S0 | R338 | S0(5) | $`20.74`$ mag | 1.9 kpc |
| Boxy E | R245 | SB(rs)0(0) | $`20.94`$ mag | 4.8 kpc |
| Note: Name and morphological type are from Richter (1989). Total absolute Gunn $`r`$ magnitude ($`M_{\mathrm{r}_\mathrm{T}}`$) and effective radius ($`r_\mathrm{e}`$) are from Milvang-Jensen & Jørgensen (2000), based on $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$. | | | | |
Fig. 2. The three example galaxies, with best-fitting ellipses, and with pixels contaminated by other objects flagged (the hatched areas).
Fig. 3. The three example galaxies: residual images (i.e. the original image minus the model image). The intensity scaling is linear, and the cuts are symmetrical around zero. Black represents positive residuals.
The example galaxies are all members of the HydraI (Abell 1060) cluster. The images used are 300 sec Gunn $`r`$ exposures ($`\lambda _{\mathrm{eff}}=6550\mathrm{\AA }`$, Thuan & Gunn 1976) obtained with the Danish 1.5-meter Telescope, La Silla, Chile. The spatial scales, as well as the grey-scales of the images of the three galaxies in the figures are identical, allowing a direct visual comparison between the three galaxies. The intensity scaling in the images shown on Figure 1 is logarithmic. We have given lengths in kpc rather than arcsec. To calculate lengths in arcsec, use $`\mathrm{log}(\mathrm{}/\mathrm{arcsec})=\mathrm{log}(\mathrm{}/\mathrm{kpc})+0.407`$. The corresponding distance modulus for the cluster is $`(mM)=34.60`$ mag.
In doing the surface photometry, the first step is to identify the other objects in the image (other galaxies, stars, and cosmic ray events) and mask (flag) these. This is shown in Figure 2. The masked pixels are not used in the surface photometry of the galaxy. The masking shown in Figure 2 uses squares for the masking of objects. A masking using circles is of course also possible, and that would mask fewer ‘uncontaminated’ pixels.
Once the masking is done, we need only provide the surface photometry task the approximate location of the center of the galaxy. The task then fits ellipses to the intensities in the image. This is done at a number of discrete radii, as also shown in Figure 2. For the ELLIPSE task, these discrete radii are specified by the rule that the different semi-major axis lengths $`a`$ are spaced by a factor of $`1.1`$.
The center ($`x`$,$`y`$) and the shape (ellipticity $`ϵ`$ and position angle) of the ellipses are kept as free parameters in the fit for semi-major axes at which the signal-to-noise ratio is sufficiently high. For larger semi-major axes, where the signal-to-noise ratio is lower, the center and shape of the ellipses are fixed.
Fig. 4. Radial profiles of intensity, ellipticity and position angle for the three example galaxies. The plotted range in equivalent radii $`r`$ is 0.3–40 kpc (0.8<sup>′′</sup>–102<sup>′′</sup>).
We get two types of output from the ellipse fit. One type is the residual image, which is the difference between the original image and the model image based on the best-fitting ellipses. The residual images for the three example galaxies are shown in Figure 3.
For the pure E, the residuals are fairly small. For the S0 in particular, and for the boxy E, the residuals are larger. The pure E is much brighter than the other two galaxies, so in relative terms the residuals for the pure E are much smaller than for the other two galaxies. For the pure E, very little structure is seen in the residual image. The residual images of the S0 and the boxy E show clear structures. In Section 2.2 we will discuss how to quantify these structures.
The other type of output from the ellipse fit is the radial profiles of a number of quantities. I.e., for each ellipse we get the intensity, center, ellipticity, position angle, and measures of the deviations from perfect elliptical isophotes (see Section 2.2), as well as the uncertainties for all these quantities. In Figure 4 we show the radial profiles of intensity, ellipticity and position angle for the three example galaxies. The position angles shown in Figure 4 are measured counter-clockwise from the $`y`$-axis of the images. The ELLIPSE task adopts this definition of position angles. The standard (astronomical) definition of position angles is from north through east. Our images have north down and east to the right, and thus need to be rotated by 180 to have the $`y`$-axis poiting towards north. However, position angles are only unique to within 180 since the major axis of a galaxy does not have a direction (as opposed to a coordinate axis). Therefore, the position angles shown in Figure 4 are expressed in the standard way. It also follows, for example, that the position angle of the inner isophotes of the boxy galaxy could be said to be 150 as well as $`30^{}`$.
In Figure 4, we have plotted the different quantities against the logarithm of the equivalent radius $`r`$. The equivalent radius is defined by $`r=\sqrt{ab}`$, where $`a`$ and $`b`$ are the semi-major and semi-minor axis of the ellipse, respectively. A circle with radius $`r`$ has an area equivalent to an ellipse with axes $`a`$ and $`b`$, hence the name equivalent radius. The program used for the illustrations in this section, ELLIPSE, uses $`a`$ to characterize the size of the ellipses. However, since the ellipticity $`ϵ`$ of the ellipse is defined as $`ϵ=1b/a`$, it follows that $`r`$ can be calculated from $`a`$ and $`ϵ`$ as $`r=a\sqrt{1ϵ}`$.
From Figure 4 it is seen how the ellipticity and position angle are free parameters until a certain radius were their values are fixed. The ellipticity is seen to vary with radius for all three galaxies. The position angle for the boxy galaxy is seen to vary rapidly in the outer parts, which is also seen in Figure 2. This behavior is known as an isophote twist.
2.2 Quantifying the deviations from elliptical shapes
As we saw from the residual images in Figure 3, the isophotes of E and S0 galaxies are not always perfectly elliptical. We wish to quantify these deviations from elliptical shapes. To illustrate how this is done, we have chosen an example ellipse for each of the three example galaxies. The three example ellipses are shown in Figure 5, 6 and 7 – they are overlayed both on the original and on the residual images of the example galaxies.
In Figure 5, 6 and 7, we also plot the quantity
$$\mathrm{\Delta }I_{\mathrm{norm}}\frac{II_0}{r\left|\frac{\mathrm{d}I}{\mathrm{d}r}\right|}.$$
$`(1)`$
Fig. 5. Top: Original and residual images of the pure E galaxy, with a single ellipse at $`a=2.6\mathrm{kpc}`$ ($`r=2.3\mathrm{kpc}`$) shown. The images are 13 kpc on the side. The ‘Data’ plot shows $`\mathrm{\Delta }I_{\mathrm{norm}}`$ versus $`\theta `$ (see text) along the shown ellipse. The ‘Data(2)’ plot shows the same, except that the intensity $`I`$ is used rather than $`\mathrm{\Delta }I_{\mathrm{norm}}`$. Little structure is seen, consistent with the fact that all the Fourier coefficients are close to zero, see Table 2.
Fig. 6. Top: Original and residual images of the S0 galaxy, with a single ellipse at $`a=2.2\mathrm{kpc}`$ ($`r=1.4\mathrm{kpc}`$) shown. The images are 13 kpc on the side. The ‘Data’ plot shows $`\mathrm{\Delta }I_{\mathrm{norm}}`$ versus $`\theta `$ along the shown ellipse. Substantial structure is seen. The starting point of the curve ($`\theta =0^{}`$) corresponds to the apogee of the ellipse that is closest to the top of the figure. At that point the residual is positive (black). The angle $`\theta `$ is measured counter-clockwise. The two most dominant Fourier modes, the 4th and 6th order cosine terms, cf. Table 2, are illustrated in the three other plots. This example galaxy is a disky galaxy and as such has a positive $`c_4`$ (4th order cosine coefficient).
Fig. 7. Top: Original and residual images of the boxy galaxy, with a single ellipse at $`a=3.8\mathrm{kpc}`$ ($`r=2.8\mathrm{kpc}`$) shown. The images are 13 kpc on the side. The ‘Data’ plot shows $`\mathrm{\Delta }I_{\mathrm{norm}}`$ versus $`\theta `$ along the shown ellipse. Substantial structure is seen. The starting point of the curve ($`\theta =0^{}`$) corresponds to the apogee of the ellipse that is closest to the top of the figure. At that point the residual is negative (white). The angle $`\theta `$ is measured counter-clockwise. The two most dominant Fourier modes, the 4th order cosine and sine terms, cf. Table 2, are illustrated in the three other plots. Note how a boxy galaxy has a negative $`c_4`$ (4th order cosine coefficient).
As can be seen, $`\mathrm{\Delta }I_{\mathrm{norm}}`$ is the deviation in intensity $`I`$ from the mean intensity at the given ellipse, $`I_0`$, normalized by the equivalent radius $`r`$, and by the absolute value of the intensity gradient, $`|\mathrm{d}I/\mathrm{d}r|`$. For comparison, we also show the intensity $`I`$ in Figure 5. The reason for choosing this particular normalization will be described in the following. One advantage of using the quantity $`\mathrm{\Delta }I_{\mathrm{norm}}`$ in the plots is that allows a direct comparison between the plots for the three example galaxies.
To quantify how the intensity deviates from being constant along the fitted ellipse, the following Fourier series is fitted to the intensity $`I(\theta )`$
$$I(\theta )=I_0+\underset{n=1}{\overset{N}{}}\left[A_n\mathrm{sin}(n\theta )+B_n\mathrm{cos}(n\theta )\right].$$
$`(2)`$
$`N`$ is the highest order fitted. $`\theta `$ is the angle measured counter-clockwise from the major axis of the ellipse. The different Fourier modes, e.g. $`\mathrm{cos}4\theta `$, will be discussed below.
We are more interested in how the isophote deviates (in the radial direction) from the fitted ellipse. Let $`R_{\mathrm{iso}}(\theta )`$ denote the distance from the center of the ellipse to the isophote, and let $`R_{\mathrm{ell}}(\theta )`$ denote the distance from the center of the ellipse to the ellipse itself. For a perfectly elliptical isophote, the difference between $`R_{\mathrm{iso}}(\theta )`$ and $`R_{\mathrm{ell}}(\theta )`$ would be zero for all values of $`\theta `$. We can Fourier expand the difference as
$$\mathrm{\Delta }R(\theta )R_{\mathrm{iso}}(\theta )R_{\mathrm{ell}}(\theta )=\underset{n=1}{\overset{N}{}}\left[A_n^{}\mathrm{sin}(n\theta )+B_n^{}\mathrm{cos}(n\theta )\right].$$
$`(3)`$
$`\mathrm{\Delta }R(\theta )`$ is the radial deviation of the isophote from elliptical shape. The relative deviation is more interesting, so we take $`\mathrm{\Delta }R(\theta )`$ relative to the size of the ellipse, given by the equivalent radius $`r`$. This relative radial deviation of the isophote from elliptical shape, $`\mathrm{\Delta }R(\theta )/r`$, is described by the Fourier coefficients $`A_n^{}/r`$ and $`B_n^{}/r`$. We will allocate new symbols for these quantities: $`s_nA_n^{}/r`$ and $`c_nB_n^{}/r`$.
The ELLIPSE task calculates $`A_n`$ and $`B_n`$, but not $`A_n^{}`$ and $`B_n^{}`$. However, we are able to link the two sets of coefficients. Consider a Taylor expansion to first order of $`I(R)`$ around $`R_0`$
$$I(R)=I(R_0)+\frac{\mathrm{d}I}{\mathrm{d}R}(RR_0).$$
$`(4)`$
Let $`R`$ be a point on the isophote and let $`R_0`$ be a point on the ellipse. $`I(R)`$ is constant since $`R`$ is on the isophote. The intensity on the ellipse $`I(R_0)`$ is not necessarily constant. The intensity on the ellipse is also given by $`I(\theta )`$ (Equation 2). We can identify the difference $`(RR_0)`$ with $`\mathrm{\Delta }R(\theta )`$ (Equation 3). A suitable mean value of the gradient $`\mathrm{d}I/\mathrm{d}R`$ is $`\mathrm{d}I/\mathrm{d}r`$. (The ‘effective intensity gradient’ $`\mathrm{d}I/\mathrm{d}r`$ can be calculated as the difference in intensity divided by the difference in equivalent radius for two adjacent ellipses.) By inserting Equation (2) and (3) in Equation (4) we find the following relations hold for all $`n`$
$$A_n^{}=\frac{A_n}{|\frac{\mathrm{d}I}{\mathrm{d}r}|},B_n^{}=\frac{B_n}{|\frac{\mathrm{d}I}{\mathrm{d}r}|},$$
$`(5)`$
where we have used $`\mathrm{d}I/\mathrm{d}r=|\mathrm{d}I/\mathrm{d}r|`$ since the intensity gradient is negative. From the definitions of $`s_n`$ and $`c_n`$ given above, we finally get
$$s_n=\frac{A_n}{r\left|\frac{\mathrm{d}I}{\mathrm{d}r}\right|},c_n=\frac{B_n}{r\left|\frac{\mathrm{d}I}{\mathrm{d}r}\right|}.$$
$`(6)`$
These definitions of the Fourier coefficients $`s_n`$ and $`c_n`$ are used in the literature by e.g. Franx, Illingworth & Heckman (1989b); Jørgensen, Franx & Kjærgaard (1992); Jørgensen & Franx (1994); and Jørgensen, Franx & Kjærgaard (1995).
Slightly different definitions are also in use. Some authors, e.g. Bender & Möllenhoff (1987); Bender, Döbereiner & Möllenhoff (1988); Bender et al. (1989); and Nieto & Bender (1989), use $`a_4/a`$ for the 4th order cosine coefficient (still for the radial deviation). In our notation, this is equal to $`B_4^{}/a`$. The only difference between $`c_4`$ and $`a_4/a`$ is that $`c_4B_4^{}/r`$ is taken relative to the equivalent radius $`r`$, whereas $`a_4/a=B_4^{}/a`$ is taken relative to the semi-major axis $`a`$. Thus, the two are related by
$$a_4/a=c_4r/a=c_4\sqrt{1ϵ}=c_4\sqrt{b/a},$$
$`(7)`$
where we have used the known relations between $`r`$, $`a`$, $`b`$ and $`ϵ`$. For an apparently round galaxy (i.e. for $`ϵ=0`$), $`a_4/a`$ is equal to $`c_4`$.
Yet another definition is used by e.g. Peletier et al. (1990). These authors expand the intensity along the ellipse as
$$I(\theta )=I_0\left(1+\underset{n=1}{\overset{N}{}}\left[S_n\mathrm{sin}(n\theta )+C_n\mathrm{cos}(n\theta )\right]\right)$$
$`(8)`$
(note the upper case $`C_n`$). When comparing with Equation (2) it is seen that $`I_0S_n=A_n`$ and $`I_0C_n=B_n`$. This means, for example, that $`c_4`$ and $`C_4`$ are related as
$$C_4=c_4\frac{r}{I_0}\left|\frac{\mathrm{d}I}{\mathrm{d}r}\right|=c_4\left|\frac{\mathrm{d}\mathrm{log}I}{\mathrm{d}\mathrm{log}r}\right|.$$
$`(9)`$
It also follows that $`a_4/a`$ is related to $`C_4`$ by
$$a_4/a=C_4\left|\frac{\mathrm{d}\mathrm{log}I}{\mathrm{d}\mathrm{log}r}\right|^1\sqrt{b/a},$$
$`(10)`$
a relation used by Faber et al. (1997) (but note the upper case $`C_4`$).
| Table 2. | Fourier coefficients for the three example galaxies at | | |
| --- | --- | --- | --- |
| | the particular ellipses shown in Figure 5, 6 and 7. | | |
| | The semi-major axis $`a`$ for these ellipses is given. | | |
| Coeff. | pure E | S0 | boxy E |
| | $`a=2.6\mathrm{kpc}`$ | $`a=2.2\mathrm{kpc}`$ | $`a=3.8\mathrm{kpc}`$ |
| $`c_3`$ | $`0.001\pm 0.002`$ | $`0.003\pm 0.013`$ | $`0.002\pm 0.008`$ |
| $`s_3`$ | $`0.001\pm 0.002`$ | $`0.000\pm 0.014`$ | $`0.002\pm 0.008`$ |
| $`c_4`$ | $`0.000\pm 0.002`$ | $`0.065\pm 0.013`$ | $`0.058\pm 0.006`$ |
| $`s_4`$ | $`0.004\pm 0.002`$ | $`0.001\pm 0.008`$ | $`0.032\pm 0.005`$ |
| $`c_5`$ | $`0.000\pm 0.002`$ | $`0.003\pm 0.008`$ | $`0.002\pm 0.004`$ |
| $`s_5`$ | $`0.001\pm 0.002`$ | $`0.007\pm 0.008`$ | $`0.006\pm 0.004`$ |
| $`c_6`$ | $`0.003\pm 0.001`$ | $`0.037\pm 0.006`$ | $`0.018\pm 0.004`$ |
| $`s_6`$ | $`0.002\pm 0.001`$ | $`0.000\pm 0.006`$ | $`0.012\pm 0.004`$ |
| $`c_7`$ | $`0.000\pm 0.001`$ | $`0.001\pm 0.006`$ | $`0.010\pm 0.003`$ |
| $`s_7`$ | $`0.002\pm 0.001`$ | $`0.000\pm 0.006`$ | $`0.005\pm 0.003`$ |
| $`c_8`$ | $`0.004\pm 0.001`$ | $`0.021\pm 0.005`$ | $`0.004\pm 0.003`$ |
| $`s_8`$ | $`0.001\pm 0.001`$ | $`0.003\pm 0.005`$ | $`0.005\pm 0.003`$ |
The fitted values of the Fourier coefficients $`c_n`$ and $`s_n`$ for the example ellipses for the three example galaxies are listed in Table 2. For the S0 galaxy, it is seen that the two numerically largest coefficients are $`c_4=0.065`$ and $`c_6=0.037`$. For the boxy galaxy the two numerically largest coefficients are $`c_4=0.058`$ and $`s_4=0.032`$.
From the definition of $`\mathrm{\Delta }I_{\mathrm{norm}}`$ (Equation 1), and from Equation (3) and (4) it is seen that $`\mathrm{\Delta }I_{\mathrm{norm}}=\mathrm{\Delta }R(\theta )/r`$, i.e. $`\mathrm{\Delta }I_{\mathrm{norm}}`$ measures the relative radial deviation of the isophote from elliptical shape. This is also the quantity measured by $`s_n`$ and $`c_n`$ (indeed, they are the Fourier coefficients of $`\mathrm{\Delta }R(\theta )/r`$), so plots of the Fourier modes $`s_n\mathrm{sin}(n\theta )`$ and $`c_n\mathrm{cos}(n\theta )`$ can be directly compared with the $`\mathrm{\Delta }I_{\mathrm{norm}}`$ versus $`\theta `$ plot. In Figure 6 and 7 we show the two most important Fourier modes and their sum for the S0 galaxy and the boxy E galaxy, respectively. It is seen in both cases how the two most important modes account for most of the structure.
Fig. 8. Radial profiles of the Fourier coefficients $`c_4`$, $`s_4`$, $`c_6`$ and $`s_6`$ for the three example galaxies. The plotted range in equivalent radii $`r`$ is 0.3–40 kpc (0.8<sup>′′</sup>–102<sup>′′</sup>). Points with uncertainty larger than 0.025 are not plotted.
An ellipse can be described by the first and second order Fourier coefficients. Since the Fourier expansion is along the best-fitting ellipses, the first and second order Fourier coefficients obtained there will be zero.
The cosine modes are symmetrical around the major axis ($`\theta =0^{}`$ or $`180^{}`$), whereas the sine modes are not, see Figure 7. The odd cosine modes (e.g. 3rd order) are not symmetrical around the minor axis (see Figure 1a in Peletier et al. 1990). The mode that is dominating for most E and S0 galaxies is the 4th order cosine. Note from Figure 6 and 7 how $`c_4`$ is an indicator of disky ($`c_4>0`$) or boxy ($`c_4<0`$) isophotes aligned with the major axis (e.g. Carter 1987; Bender et al. 1989; Peletier et al. 1990).
The radial profiles of the 4th and 6th order Fourier coefficients are shown in Figure 8.
2.2.1 The output from the ELLIPSE task
The radial profiles determined by the ELLIPSE task are output to an STSDAS table. The most important columns are listed in Table 3, along with their meaning in our notation.
The equivalent radius is not in the table, but it can be calculated as SMA\*sqrt(1-ELLIP). The 5th to 8th order Fourier coefficients $`s_n`$ and $`c_n`$ can be calculated as e.g. $`c_6`$ = BI6/(SMA\*abs(GRAD)). This follows from Equation (6) since $`a\mathrm{d}I/\mathrm{d}a=r\mathrm{d}I/\mathrm{d}r`$. Note that the GRAD column is $`\mathrm{d}I/\mathrm{d}a`$ (I. Busko 2000, private communication).
| Table 3. Output columns from ELLIPSE | |
| --- | --- |
| Column name | Content in our notation |
| SMA | $`a`$ |
| INTENS | $`I`$ |
| ELLIP | $`ϵ`$ |
| PA | PA (position angle) |
| X0 | $`x`$ (center of ellipse) |
| Y0 | $`y`$ (center of ellipse) |
| GRAD | $`\mathrm{d}I/\mathrm{d}a`$ (not $`\mathrm{d}I/\mathrm{d}r`$) |
| A$`n`$ | $`s_n`$, $`n`$ = 3, 4 |
| B$`n`$ | $`c_n`$, $`n`$ = 3, 4 |
| AI$`n`$ | $`A_n`$, $`n`$ = 5, 6, 7, 8<sup>a</sup> |
| BI$`n`$ | $`B_n`$, $`n`$ = 5, 6, 7, 8<sup>a</sup> |
| <sup>a</sup> Only calculated with the following option set: | |
| harmonics="5 6 7 8". | |
2.3 Determination of magnitudes
The intensity $`I`$ (in ADU) contains the signal from the galaxy plus the sky background. The sky background level can be determined in several ways. One way is to identity empty regions in the image and measure the level there. If the galaxy fills most of the image, this can be difficult. Another way is to fit a suitable analytical expression to outer part of galaxy plus sky intensity profile obtained from the surface photometry. This method has been used by e.g. Jørgensen et al. (1992), fitting $`I_{\mathrm{galaxy}+\mathrm{sky}}(r)=I_{\mathrm{sky}}+I_{\mathrm{galaxy},0}r^\alpha `$, with $`\alpha `$ = 2 or 3.
Magnitudes can be calculated from the sky subtracted intensity. This can either be integrated magnitudes within a certain aperture (elliptical or circular), or the surface brightness at a given ellipse, $`\mu (r)`$. By knowing the pixel scale of the CCD (in arcsec/pixel), $`\mu (r)`$ can be expressed in units of mag/arcsec<sup>2</sup>. With the use of observed standard stars, the magnitudes and surface brightnesses can be transformed to a standard photometric system.
2.4 Global parameters
The surface photometry has produced radial profiles of a number of quantities. It is desirable to condense these radial profiles to a few characteristic numbers, the global parameters.
2.4.1 Effective parameters
Elliptical galaxies have surface brightness profiles that are well approximated by the $`r^{1/4}`$ law (de Vaucouleurs 1948). By fitting the aperture magnitudes to an $`r^{1/4}`$ growth curve, the following two parameters can be derived:
$``$$`r_\mathrm{e}`$: Effective radius, in arcsec
$``$$`\mu _\mathrm{e}`$: Mean surface brightness within $`r_\mathrm{e}`$, in mag/arcsec<sup>2</sup> The seeing needs to be taken into account (Saglia et al. 1993).
For galaxies with perfect $`r^{1/4}`$ profiles, the effective radius $`r_\mathrm{e}`$ is the half-light radius, i.e. the radius that encloses half of the light from the galaxy. Spiral galaxies are better described by an exponential surface brightness profile than by an $`r^{1/4}`$ profile.
We can express the mean surface brightness in units of $`L_{}/\mathrm{pc}^2`$, where $`L_{}`$ is the luminosity of the Sun in the given passband (e.g. Gunn $`r`$). We will call this quantity $`I_\mathrm{e}`$. The relation is $`\mathrm{log}I_\mathrm{e}=0.4(\mu _\mathrm{e}k)`$, where the constant $`k`$ is given by $`k=M_{}+5\mathrm{log}(206265\mathrm{pc}/10\mathrm{pc})`$. As is seen, the calculation does not involve the distance to the galaxy, but only the absolute magnitude of the Sun in the given passband. The calculation of $`r_\mathrm{e}`$ in kpc from $`r_\mathrm{e}`$ in arcsec, however, does involve the distance to the galaxy.
With $`r_\mathrm{e}`$ as the half-light radius, it follows that the total luminosity is given by $`L=2\pi I_\mathrm{e}r_\mathrm{e}^2`$.
2.4.2 Global Fourier parameters
As is seen from Figure 8, also the Fourier parameters vary with radius. One way of getting a characteristic value of e.g. the $`c_4(r)`$ profile is to take the extremum value. In case the profile does not have a clear extremum, we can take the value at the effective radius. We will use the symbol $`c_4`$ for this characteristic value of $`c_4(r)`$.
Another way to get a global Fourier coefficient is to calculate an intensity weighted mean value as
$$s_n\frac{_{r_{\mathrm{min}}}^{r_{\mathrm{max}}}I(r)s_n(r)dr}{_{r_{\mathrm{min}}}^{r_{\mathrm{max}}}I(r)dr},c_n\frac{_{r_{\mathrm{min}}}^{r_{\mathrm{max}}}I(r)c_n(r)dr}{_{r_{\mathrm{min}}}^{r_{\mathrm{max}}}I(r)dr},$$
$`(11)`$
where $`r_{\mathrm{min}}`$ is the radius where seeing effects are no longer important, and $`r_{\mathrm{max}}`$ is the radius where the Fourier coefficients can no longer reliably be determined (see Jørgensen & Franx 1994).
2.4.3 Global ellipticities and colors
As global ellipticity can be taken the extremum of the $`ϵ(r)`$ profile, the value of $`ϵ(r)`$ at the effective radius, or the value of $`ϵ(r)`$ at a certain isophote level, e.g. $`\mu =21.85{}_{}{}^{\mathrm{m}}/\mathrm{arcsec}^2`$ in Gunn $`r`$ as used by Jørgensen & Franx (1994).
As global color, the color within the effective radius can be used. By color is meant the difference between the magnitudes in two different passbands, such as (B$``$V). The color is always calculated as the magnitude in the passband with the shortest effective wavelength minus the magnitude in the passband with the longest effective wavelength. Thus, a large value of the color means that the galaxy is red, and a small value means that it is blue. Another global parameter related to the color is the color gradient, defined as $`\mathrm{\Delta }\mathrm{color}/\mathrm{\Delta }\mathrm{log}r`$, i.e. as the slope of the color versus $`\mathrm{log}r`$ plot.
3. QUANTITATIVE MORPHOLOGY FOR E AND S0 GALAXIES
CCD surface photometry as described in the previous section offers the possibility of carrying out quantitative morphologic studies of E and S0 galaxies. The traditional classification of these galaxies distinguishes between E and S0 galaxies based on the presense of a disk (S0 galaxies) or no disk (E galaxies). E and S0 galaxies are considered to belong to two separate classes of galaxies.
In the following we summarize the methods and results presented by Jørgensen & Franx (1994, hereafter JF94). The results lead to the conclusion that E and S0 galaxies fainter than an absolute blue magnitude of $`M_{\mathrm{B}_\mathrm{T}}=22`$ mag form one class of galaxies with a broad and continuous distribution of the relative luminosity of the disks.
3.1. Sample properties and the data
The sample used in the study by JF94 is a magnitude limited sample of galaxies within the central square degree of the Coma cluster. The sample is complete to an apparent magnitude in Gunn $`r`$ of 15.3 mag. 171 galaxies are included in the sample. Because the sample is well-defined and complete it is possible based on these data to draw conclusions regarding the E and S0 galaxies as a class.
CCD photometry was obtained of the full sample. Surface photometry for the galaxies was derived using GALPHOT (Jørgensen et al. 1992; Franx et al. 1989b). From the surface photometry, global parameters were derived. The important parameters used in this study are summarized in Table 4. $`c_4`$ and $`c_6`$ are intensity weighted mean values of $`c_4`$ and $`c_6`$, respectively (see JF94 and Section 2.4.2), while $`c_4`$ represents the extremum the $`c_4`$ radial profile. The intensity weighted parameters are less sensitive to small scale features in the radial profiles and are therefore to be preferred for studies of global properties.
| Table 4. Surface photometry parameters | |
| --- | --- |
| Parameter | Description |
| $`m_\mathrm{T}`$ | Total magnitude in Gunn $`r`$ |
| $`ϵ_{21.85}`$ | Ellipticity at $`\mu =21.85{}_{}{}^{\mathrm{m}}/\mathrm{arcsec}^2`$ in Gunn $`r`$ |
| $`c_4`$ | Extremum of the $`c_4\left(r\right)`$ radial profile |
| $`c_4`$ | Intensity weighted mean value of $`c_4\left(r\right)`$ |
| $`c_6`$ | Intensity weighted mean value of $`c_6\left(r\right)`$ |
3.2. Morphology of the E and S0 galaxies
Figure 9 shows the morphological parameters $`ϵ_{21.85}`$, $`c_4`$ and $`c_6`$ as functions of the total magnitudes $`m_\mathrm{T}`$ in Gunn $`r`$. The spiral galaxies are also shown on these figures, but we will omit them from the following discussion. From Figure 9 it is clear that the galaxies fainter than about 12.7 mag span the full range in morphological parameters. The E and S0 galaxies cannot be separated into two classes of galaxies based on one of these morphological parameters. The difference in properties of E and S0 galaxies fainter and brighter than $`m_\mathrm{T}=12.7`$ mag is striking. This demarcation magnitude corresponds to a blue absolute magnitude of $`M_{\mathrm{B}_\mathrm{T}}=22`$ mag.
In Figure 10 we show the distribution of the ellipticities, $`ϵ_{21.85}`$, of the galaxies. Figure 10b shows the cumulative frequency of $`ϵ_{21.85}`$ for the E galaxies and the S0 galaxies separately, as well as the cumulative frequency of $`ϵ_{21.85}`$ for the E and S0 galaxies as one class.
Fig. 9. Morphological parameters versus the total magnitudes. Open boxes – E galaxies with dominating regular $`c_4`$-profiles; filled boxes – S0 galaxies with dominating regular $`c_4`$-profiles; crosses – E galaxies with irregular or non-dominating $`c_4`$-profiles; stars – S0 galaxies with irregular or non-dominating $`c_4`$-profiles; triangles – spirals. The $`c_4`$-profiles are considered non-dominating if $`|c_4|`$ is more than one sigma smaller than the absolute value of one of the other intensity weighted mean coefficients. Typical measurement errors are given on the panels. Galaxies with uncertainty on $`c_4`$ respectively $`c_4`$ larger than 0.02 are not plotted. The six brightest galaxies have small Fourier coefficients and $`ϵ_{21.85}<0.4`$. Other dependence on $`m_\mathrm{T}`$ is not seen. (From JF94.)
If the E galaxies and the S0 galaxies form two separate classes, then we expect that their $`ϵ_{21.85}`$ distributions can be modeled independently. The simplest assumption is that the galaxies are randomly oriented in space and have some simple distribution of the intrinsic ellipticities. JF94 attempted to fit the $`ϵ_{21.85}`$ distributions with intrinsically uniform ellipticity distributions and with intrinsic ellipticity distribution that were Gaussians. The fit was done by maximizing the probability that the data was drawn from the model as reflected by a Kolmogorov–Smirnov (K–S) test (e.g. Press et al. 1992). The K–S test gives the probability that a data distribution is drawn from a model distribution (or another data distribution) based on the maximum difference between the cumulative frequencies of the two distributions, as it is illustrated in Figure 10b. JF94 found that the $`ϵ_{21.85}`$ distribution of the E galaxies could be fitted satisfactory with either a uniform or a Gaussian intrinsic distribution of the ellipticities. Both of these models resulted in probabilities of 84 per cent or larger. However, the $`ϵ_{21.85}`$ distribution of the S0 galaxies could not be fitted satisfactory. Both intrinsic distributions have probabilities of only 11 per cent. The $`ϵ_{21.85}`$ distribution of the S0 galaxies lack apparently round galaxies, see Figure 10b. JF94 also find that the $`ϵ_{21.85}`$ distribution of the E and S0 galaxies treated as one class can be fitted satisfactory by either a uniform or a Gaussian intrinsic distribution of the ellipticities.
Fig. 10. Distributions of the apparent ellipticities. (a) Solid line – E galaxies, dashed line – S0 galaxies, Dashed-dotted line – spirals. (b) Solid lines – E, S0 galaxies, and E and S0 galaxies together. The six brightest E galaxies are excluded. Dotted lines – best fitting uniform distributions. Dashed lines – best fitting Gaussian distributions. (From JF94.)
These results show that some (maybe all) of the E galaxies fainter than $`m_\mathrm{T}=12.7`$ mag must be S0 galaxies seen face-on. When the galaxies are seen face-on a disk is more difficult to detect (as also noted by Rix & White 1990), and the galaxies are mostly classified as E galaxies even when a disk is present.
3.3 The relative disk luminosities
The next task is to determine the relative disk luminosities, $`L_\mathrm{D}/L_{\mathrm{tot}}`$, of the galaxies. $`L_\mathrm{D}/L_{\mathrm{tot}}`$ is the luminosity of the disk relative to the total luminosity of the galaxy. JF94 showed that two of the morphological parameters can be used together to derive $`L_\mathrm{D}/L_{\mathrm{tot}}`$ if a simple model for the disk and the bulge is assumed.
Fig. 11. Fourth and 6th order Fourier coefficients plotted against ellipticity and versus each other. Data symbols as in Figure 9. Galaxies with uncertainty on $`c_4`$ respectively $`c_4`$ larger than 0.02 are not plotted. Typical measurement errors are given on the panels. The models with $`a_{\mathrm{eB}}=a_{\mathrm{eD}}`$ are overplotted. The curves are labeled with $`L_\mathrm{D}/L_{\mathrm{tot}}`$. The dashed lines on (b) mark the inclinations where $`\mathrm{cos}i=0.1`$, 0.2, 0.3, 0.4, and 0.5 (i.e. $`i=84^{}`$, $`78^{}`$, $`73^{}`$, $`66^{}`$, and $`60^{}`$), from top to bottom. No dashed line is shown for the $`i=90^{}`$ (i.e. edge-on) models, but this line would be above the $`i=84^{}`$ dashed line, connecting the end points of the solid lines. It is seen that for a given $`L_\mathrm{D}/L_{\mathrm{tot}}`$ the highest inclination gives the largest $`c_4`$. The non-zero coefficients for the pure disk-model are due to the inclusion of seeing effects in the models. (From JF94.)
Figure 11 illustrates this technique for the Coma cluster sample. JF94 constructed models consisting of a bulge with an $`r^{1/4}`$ profile and a disk with an exponential profile. The bulge is assumed to have an intrinsic ellipticity of 0.3, while the intrinsic ellipticity of the disk is assumed to be 0.85. JF94 tested models for both equal major axis of the two components, $`a_{\mathrm{eB}}=a_{\mathrm{eD}}`$, and for $`a_{\mathrm{eB}}=0.5a_{\mathrm{eD}}`$. The results are not significantly different, so here we will concentrate on the $`a_{\mathrm{eB}}=a_{\mathrm{eD}}`$ models. JF94 convolved the model images of bulge plus disk with a representative seeing and then analyzed the model images in the same way as done for the data. Models with relative disk luminosities $`L_\mathrm{D}/L_{\mathrm{tot}}`$ between zero (no disk) and one (all disk) were constructed. Further, the inclination was varied between face-on (small inclination, cos $`i`$ = 1) and edge-on (large inclination, cos $`i`$ = 0) in steps of 0.1 in cos $`i`$.
The results in terms of morphological parameters are shown on Figure 11. The models reproduce the general variation of $`c_4`$, $`c_4`$ and $`c_6`$ with ellipticity. They also span reasonably well the section of the $`c_4`$$`c_6`$ diagram covered by the data. Models of the kind used by JF94 cannot reproduce boxy isophotes of the galaxies. However, the Coma cluster sample contains only two galaxies fainter than $`m_\mathrm{T}=12.^\mathrm{m}7`$ that have $`c_4`$ significantly smaller than zero. For galaxies with ellipticities larger than 0.3 and $`c_4`$ larger than 0.007, the models are well separated in $`c_4`$ versus $`ϵ_{21.85}`$, see Figure 11b. Thus, these two parameters can be used to derive the relative disk luminosities, $`L_\mathrm{D}/L_{\mathrm{tot}}`$, and the inclinations, $`i`$, of the galaxies.
Fig. 12. Relative (a) and cumulative (b) frequency for the relative disk luminosities. Solid line – determinations from the $`ϵ_{21.85}`$$`c_4`$ diagram. The distribution has been normalized with the total number of galaxies fainter than $`m_\mathrm{T}=12.^\mathrm{m}7`$. Dashed line – model prediction for a uniform intrinsic distribution, corrected for the incompleteness due to the limits enforced on $`ϵ_{21.85}`$ and $`c_4`$. Dotted lines – model predictions for a uniform intrinsic distribution plus a fraction of diskless galaxies. Models for fractions of 0.1, 0.2, and 0.5 are shown. A higher fraction of diskless galaxies moves the curve for the normalized distribution downwards in both (a) and (b). (From JF94.)
Figure 12 shows the resulting distribution of $`L_\mathrm{D}/L_{\mathrm{tot}}`$. It was possible to derive $`L_\mathrm{D}/L_{\mathrm{tot}}`$ for 52 of the E and S0 galaxies in the sample. The cumulative frequency shown on Figure 12b is normalized to the 140 E and S0 galaxies fainter than $`m_\mathrm{T}=12.7`$ mag. Overplotted on Figure 12 are models of the distribution of $`L_\mathrm{D}/L_{\mathrm{tot}}`$. These models are uniform distributions with some fraction of diskless galaxies added. The best fitting model is a uniform distribution of $`L_\mathrm{D}/L_{\mathrm{tot}}`$ between zero and one with an additional 10 per cent diskless galaxies. JF94 find that the resulting distributions of inclinations, $`ϵ_{21.85}`$ and $`c_4`$ for this model also fit the observed distributions of these parameters.
3.4 Conclusions from JF94
The results presented by JF94 illustrate the strengths of studies of quantitative morphology based on surface photometry for statistically well-defined and complete samples. JF94 were able to show that the E and S0 galaxies in the Coma cluster (fainter than $`m_\mathrm{T}=12.7`$ mag in Gunn $`r`$) form one class of galaxies with a broad (most likely uniform) distribution of $`L_\mathrm{D}/L_{\mathrm{tot}}`$. This result contradicts the traditional classification of these galaxies into two separate classes. It also provides constraints for models for morphological evolution of galaxies, since the end-result for the Coma cluster represents one scenario that the models need to reproduce.
4. STELLAR POPULATION MODELS
Stellar population models are tools for interpreting the integrated light, such as the colors, observed from galaxies. Ideally, we want to determine the mix of stars that give rise to the observations. This problem is usually underconstrained, so it is necessary to make some assumptions regarding how the numbers of different types of stars are related. Here we will consider so-called single-age single-metallicity models, also known as single stellar population (SSP) models. In these models, all the stars are formed at the same time, with distribution in mass given by the chosen initial mass function (IMF), and with identical chemical composition.
SSP models are based on the following ingredients. First, theoretical stellar isochrones are needed. Isochrones are loci in the theoretical HR-diagram $`(\mathrm{log}T_{\mathrm{eff}},\mathrm{log}L)`$ for a stellar population of a given age and chemical composition. Second, a conversion is needed between the theoretical parameters of $`T_{\mathrm{eff}}`$, $`L`$, $`\mathrm{log}g`$ (surface gravity) and the metallicity \[M/H\] to the observable parameters such as colors. This conversion can be either empirical or theoretical. The empirical conversion is based on observations of individual stars in our Galaxy for which the ‘theoretical parameters’ can be inferred, and the observable parameters measured. The theoretical conversion is based on model atmospheres and synthetic spectra. Third, the IMF has to be specified.
An example of SSP models are those presented by Vazdekis et al. (1996). These models use the isochrones from the Padova group (Bertelli et al. 1994). The conversion from the theoretical to the observable parameters is empirical. Models are presented for several different IMFs. One IMF is a constant below 0.2 $`M_{}`$, a Salpeter (1955) IMF above 0.6 $`M_{}`$, and a spline in the interval 0.2–0.6 $`M_{}`$. The models that we use in Section 5 are based on this IMF.
When the IMF has been specified, the models have only two parameters: age and metallicity. The metallicity can be expressed either as the mass fraction in heavier elements, $`Z`$, or as the metal abundance \[M/H\] $`\mathrm{log}(Z/Z_{})`$ (with $`Z_{}=0.02`$). The models have solar abundance ratios, while this may not be the case for all galaxies.
Fig. 13. Example of the predictions from the Vazdekis et al. (1996) models: (B$``$R<sub>c</sub>) color as function of age (the $`x`$-axis) and metallicity (the different symbols). Triangles – $`[\mathrm{M}/\mathrm{H}]=0.4`$; crosses – $`[\mathrm{M}/\mathrm{H}]=0.0`$; boxes – $`[\mathrm{M}/\mathrm{H}]=0.4`$. It is seen that the stellar populations get redder with higher age and higher metallicity.
As an example, the model predictions for the (B$``$R<sub>c</sub>) color is given in Figure 13. It is seen that the color depends both on age and metallicity. As will be illustrated in the next section, optical colors depend roughly in the same way on age and metallicity. Therefore, in a plot of one optical color versus another optical color the model lines of constant age will be almost on top of the model lines of constant metallicity. The effects of age and metallicity cannot be disentangled in such a diagram. This is known as the age–metallicity degeneracy (e.g. Worthey 1994).
It should be noted that real galaxies are not necessarily single stellar populations. For example, a galaxy could have experienced a second star formation event. Therefore, when SSP model predictions are compared with data for real galaxies to determine the age and the metallicity, the resulting ages and metallicities are luminosity weighted mean values. Further, dust can also cause a stellar population to appear red. This can be an additional complication in determining ages and metallicities.
5. COLOR RELATIONS
The E and S0 galaxies follow very well-defined relations between the optical colors and the total magnitudes. Figure 14 shows the optical color-magnitude relation for the Coma cluster. The figure includes all objects detected in a field covering the central 75 arcmin $`\times `$ 80 arcmin of the cluster. The data were obtained with the McDonald Observatory 0.8-meter Telescope equipped with the Prime Focus Camera.
The color-magnitude relation is well-defined and has a very low scatter for (E and S0) galaxies brighter than about $`R_\mathrm{c}=17`$ mag, see Figure 14. For galaxies fainter than $`R_\mathrm{c}=17.5`$ mag only sparse redshift information is available, and many of these faint galaxies may be background galaxies.
The color-magnitude relation is thought to be primarily a result of differences in metallicity as a function of luminosity. However, recent results based on spectroscopy show that both the mean ages and the mean metallicities varies for E and S0 galaxies at low redshifts (e.g. Worthey, Trager & Faber 1995; Jørgensen 1999). Thus, there is a need for interpreting the color-magnitude relation within these recent results in order to achieve an self-consistent interpretation of the spectroscopic and the photometric results. The work by Kauffmann & Charlot (1998) represents one of the only attempts to model the color-magnitude relation and relations involving spectroscopic information in a consistent manner. In the models by Kauffmann & Charlot both age and metallicity varies with the luminosity.
Fig. 14. Color-magnitude relation for the central 75 arcmin $`\times `$ 80 arcmin of the Coma cluster. Both stars and galaxies are included on this figure, a total of 15370 objects. The lack of objects in the upper right hand corner is due to the combination of the magnitude limits in B and $`R_\mathrm{c}`$. Small black points – stars; black \[blue\] boxes – spirals and irregulars, confirmed members; dark grey \[red\] boxes – E and S0 galaxies, confirmed members; light grey \[green\] boxes – unclassified confirmed members; small grey \[orange\] points – other galaxies in the field, members and non-members. The photometric data are from Jørgensen (2000). The redshift data are from Jørgensen & Hill (2000).
Color-color relations in the optical are in Figure 15 shown for the same sample of objects as shown in Figure 14. The confirmed members of the cluster form a tight relation in these two color-color relations. This is in agreement with predictions from stellar population models (see Section 4), which predicts the optical colors to be degenerate in age and metallicity. This is seen on Figure 15, left panel, where the same models as shown on Figure 13 are overplotted. The lines of constant age fall right on top of the lines of constant metallicity, forming a single line along the ridge of the location of the E and S0 galaxies. This means that the optical colors alone cannot be used to derive the mean ages and the mean metallicities. A younger age will lead to bluer colors, but a similar change in colors can be caused by a lower metallicity.
Fig. 15. Optical color-color relations for the Coma cluster. Symbols as in Figure 14. Solid line on the left panel – SPP models from Vazdekis et al. (1996), see text.
While the optical color-color diagrams are degenerate in age and metallicity, the combination of one optical color and one optical-infrared color may be used to break the degeneracy. An example of this is shown in Figure 16 for a small sample of E and S0 galaxies in the Coma cluster. The optical color (U$``$B) and the optical-infrared color (V$``$K) are both sensitive to both age differences and metallicity differences. However, (U$``$B) is more sensitive to the age differences than to metallicity differences, while the opposite is the case for (V$``$K). The stellar population models in the near infrared (near-IR), JHK, in this case the K-band, are still rather uncertain, but the technique is promising for studies of faint high redshift galaxies for which spectroscopy would come at a very high cost of telescope time at 8-meter class telescopes. The apparently rather low mean age of the Coma cluster galaxies as estimated from Figure 16 is in fact in agreement with results based on spectroscopic information (see Jørgensen 1999).
Fig. 16. Optical-infrared color-color relation for the Coma cluster. The data are from Jørgensen (2000) and Mobasher et al. (1999). The lines represent stellar population models by Vazdekis et al. (1996). Solid lines – metallicities of \[Fe/H\]=$``$0.4, 0.0 and 0.4; the lowest metallicity leads to the smallest (V$``$K). Dashed lines – ages of 1, 2, 5, 8, 12, 15, and 17 Gyr; the largest ages lead to the largest (U$``$B).
6. THE FUNDAMENTAL PLANE
The Fundamental Plane (FP) is a relation that combines surface photometry with spectroscopy. We will discuss this relation both at low and at high redshift.
The FP (Djorgovski & Davis 1987; Dressler et al. 1987) is the relation
$$\mathrm{log}r_\mathrm{e}=\alpha \mathrm{log}\sigma +\beta \mathrm{log}I_\mathrm{e}+\gamma ,$$
$`(12)`$
where $`\sigma `$ is the line-of-sight stellar velocity dispersion for the galaxy in question. In other words, the measured values of $`\mathrm{log}r_\mathrm{e}`$, $`\mathrm{log}I_\mathrm{e}`$ and $`\mathrm{log}\sigma `$ for a sample of E and S0 galaxies do not populate this 3-parameter space evenly, but are limited to a thin plane.
The velocity dispersion $`\sigma `$ is determined from spectroscopy, see Figure 17. The absorption lines in galaxies are broadened due to the internal motions of the stars in the galaxy. To determine how much the stars are moving, it is necessary to know what the spectrum of the galaxy would be if all the stars were at rest with respect to each other. This is approximated by the spectrum of af K giant star. This template star spectrum is broadened by a Gaussian broadening function until it matches the galaxy spectrum. The velocity dispersion of the galaxy is then the dispersion $`\sigma `$ of the broadening function. In more precise terms, this determination of $`\sigma `$ can be done using the Fourier fitting method (Franx, Illingworth & Heckman 1989a) or the Fourier quotient method (Sargent et al. 1977).
Fig. 17. Illustration of the effect of the velocity dispersion $`\sigma `$. The top panel shows a K giant star in our own Galaxy. This star is representative of the stellar populations in E or S0 galaxies. The two lower panels show E or S0 galaxies in the HydraI cluster. The $`\mathrm{Mg}\mathrm{b}`$ absorption line triplet at 5177 Å (rest frame) is broadened in the galaxy spectra. The instrumental resolution is 79 km/s.
6.1 The interpretation of the Fundamental Plane
The physics behind the FP can be illuminated by some simple arguments (Djorgovski, de Carvalho & Han 1988; Faber et al. 1987). Consider the virial theorem for the stars in the galaxy
$$\frac{GM}{R}=2\frac{V^2}{2},$$
$`(13)`$
We relate the observable quantities $`r_\mathrm{e}`$, $`\sigma `$ and $`I_\mathrm{e}`$ to the ‘physical’ quantities $`R`$, $`V^2`$ and luminosity $`L`$ through
$$r_\mathrm{e}=k_\mathrm{R}R,\sigma ^2=k_\mathrm{V}V^2,L=k_\mathrm{L}I_\mathrm{e}r_\mathrm{e}^2,$$
$`(14)`$
The parameters $`k_\mathrm{R}`$, $`k_\mathrm{V}`$, and $`k_\mathrm{L}`$ reflect the density structure, kinematical structure, and luminosity structure of the given galaxy. If these parameters are constant, the galaxies constitute a homologous family. Homology means that structure of small and big galaxies is the same.
Combining Equation (13) and (14) gives
$$r_\mathrm{e}=k_\mathrm{S}(M/L)^1\sigma ^2I_\mathrm{e}^1,k_\mathrm{S}=(Gk_\mathrm{R}k_\mathrm{V}k_\mathrm{L})^1.$$
$`(15)`$
For homology $`k_\mathrm{S}`$ will be constant. When this relation is compared to the observed FP,
$$r_\mathrm{e}=\mathrm{constant}\sigma ^{1.24\pm 0.07}I_\mathrm{e}^{0.82\pm 0.02}$$
$`(16)`$
(Jørgensen, Franx & Kjærgaard 1996, in Gunn $`r`$), it is seen that the coefficients of the FP are not 2 and $`1`$ as expected from homology and constant mass-to-light ratios. The product $`k_\mathrm{S}(M/L)^1`$ cannot be constant, but has to be a function of $`\sigma `$ and $`I_\mathrm{e}`$. A non-constant $`k_\mathrm{S}(M/L)^1`$ can be explained by a systematic deviation from homology ($`k_\mathrm{S}`$ varies), or a systematic variation of the $`M/L`$ ratios, or both. When homology is assumed, the observed FP coefficients give the relation
$$M/L_\mathrm{r}M^{0.24\pm 0.03},$$
$`(17)`$
(Jørgensen et al. 1996). The interpretation of the FP is still a matter of debate. The $`M/LM^b`$ interpretation seems to be the most favored one, although there is some evidence that non-homology may play a role too (e.g. Hjorth & Madsen 1995; Pahre, de Carvalho & Djorgovski 1998).
6.2 The evolution of the Fundamental Plane as a function of redshift
The Fundamental Plane can be used to study the evolution of galaxies as a function of redshift. As explained above, the FP may be interpreted as a relation between the masses and the $`M/L`$ ratios of the galaxies. Under the assumption that the masses do not change with redshift, e.g. no merging takes place, the evolution of the FP zero point with redshift can be interpreted as the evolution of the $`M/L`$ ratios.
Several authors have studied the FP for clusters at redshifts higher than 0.1, see Table 5 for clusters and references. Additional studies by Pahre, Djorgovski & de Carvalho (1999) and Kelson et al. (1999) are soon to be published in refereed journals.
Table 5. Fundamental Plane studies of cluster with $`z>0.1`$ Cluster $`z`$ $`N_{\mathrm{gal}}`$ Reference A2218 0.18 9 Jørgensen & Hjorth (1997), Jørgensen et al. (1999) A665 0.18 6 Jørgensen & Hjorth (1997), Jørgensen et al. (1999) CL1358+62 0.33 10 Kelson et al. (1997) MS1512+36 0.37 2 Bender et al. (1998) A370 0.37 7 Bender et al. (1998) CL0024+16 0.39 8 van Dokkum & Franx (1996) MS2053$``$04 0.58 5 Kelson et al. (1997) MS1054$``$03 0.83 8 van Dokkum et al. (1998)
As examples of high redshift studies of the FP we show in Figure 18 the FP for the Coma cluster and for five clusters with redshift larger than 0.1. The data for the Coma are from Jørgensen (1999) and Jørgensen et al. (1995). Abell 2218 and Abell 665 are discussed by Jørgensen et al. (1999). The sources for the rest of the clusters are given on the figure.
Fig. 18. The FP edge-on for Coma, A2218, A665, CL1358+62, CL0024+16, and MS2053$``$04. The sources of the data are given on the panels (‘This paper’ refers to Jørgensen et al. 1999). The skeletal symbols on panel (c) and (d) are the E+A galaxies. The photometry is calibrated to Gunn $`r`$ in the rest frames of the clusters. The mean surface brightness $`\mathrm{log}I_\mathrm{e}=0.4(\mu _\mathrm{e}26.4)`$ is in units of $`\mathrm{L}_{}/\mathrm{pc}^2`$ (cf. Section 2.4.1). The photometry is not corrected for the dimming due to the expansion of the Universe. The effective radii are in kpc ($`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and $`q_\mathrm{o}=0.5`$). The solid lines are the FPs with coefficients adopted from Jørgensen et al. (1996), and with zero points derived from the data presented in the figure. Typical error bars are given on the panels; the thin and thick error bars show the random and systematic uncertainties, respectively. (From Jørgensen et al. 1999.)
From the change in the zero point of the FP as a function of redshift Jørgensen et al. (1999), in agreement with other studies, find that the $`M/L`$ ratios of the E and S0 galaxies change very slowly with redshift. The sample of clusters shown on Figure 18 span about half of the current age of the Universe (for $`q_0=0.5`$). Under the assumption that the galaxies evolve passively over this time interval, e.g. no merging and no formation of new E and S0 galaxies, then it is possible to put limits of the redshift at which the majority of the stars were formed. This redshift is called the formation redshift. The study by Jørgensen et al. (1999) as well as other studies conclude that the formation redshift is larger than about 2.5 for $`q_0=0.5`$ and larger than about 1.5 for $`q_0=0.15`$.
It is important to keep in mind that the assumption regarding passive evolution represents a very simplified view of the galaxy evolution. Most likely the real evolution since $`z0.6`$ cannot be modeled with passive evolution. If the galaxies experience on-going star formation, then the observed evolution will appear smaller than for passive evolution because the galaxies are continuously forming young bright stars. A similar effect can be caused by a series of smaller bursts of star formation. Finally, the interactions, the possible merging and the morphological evolution of the galaxies over the last half of the age of the Universe cannot be ignored. Some of the E and S0 galaxies that we observe in low redshift clusters may not have ended up in the samples if we could have observed those clusters at a much earlier stage in their evolution, simply because some of the E and S0 galaxies may have formed recently by merging of spiral galaxies.
7. SUGGESTED FUTURE PROJECTS
We end this paper by a brief summary of some of the projects that may be carried out building on the techniques and results discussed in this paper. We concentrate on projects that involve photometry only. Some of the projects may be carried out using existing archive data from HST.
Evolution of morphology as a function of redshift
Dressler et al. (1997) have recently used HST/WFPC2 data of clusters to study the morphology-density relation as a function of redshift. Dressler et al. found that the fraction of S0 galaxies is lower at high redshift than at low redshift, while the fraction of spirals is higher at high redshift than at low redshift. However, this study is based on the traditional method of classifying galaxies. We suggest that a quantitative approach is taken to any study of the morphology. Of special interest would be to study how the relative disk luminosities $`L_\mathrm{D}/L_{\mathrm{tot}}`$ for E and S0 galaxies evolve with redshift, both in clusters and in the field. The first study of this kind was done for the cluster CL0024+16 ($`z=0.39`$) by Bergmann & Jørgensen (1999) who found that the $`L_\mathrm{D}/L_{\mathrm{tot}}`$ distribution for the E and S0 galaxies in CL0024+16 shows a paucity of disk-dominated galaxies when compared to the Coma cluster.
It would be valuable to apply the same technique and derive $`L_\mathrm{D}/L_{\mathrm{tot}}`$ for a larger sample of cluster and field galaxies at redshifts larger than 0.1 to establish the possible evolution of the distribution of $`L_\mathrm{D}/L_{\mathrm{tot}}`$. This may be done using the HST/WFPC2 archive data.
Studies of global colors
There have been many studies of the global colors of galaxies as a function of redshift (e.g. Stanford, Eisenhardt & Dickinson 1998; Bower, Kodama & Terlevich 1998; Kodama et al. 1998 and references in these papers). However, most studies concentrate on the optical colors of the galaxies, while the near-IR (JHK) data is very sparse. The study by Stanford et al. includes near-IR data and addresses the question of how the color-magnitude relations for the near-IR colors evolve with redshift.
The combination of optical and optical-infrared colors may be used to break the age-metallicity degeneracy (see Section 5). Observations of low redshift E and S0 galaxies in the near-IR may be used to establish the zero redshift properties and the methods and models needed to break the age-metallicity degeneracy. For high redshift galaxies ($`z>0.5`$), the near-IR photometry may be obtained with 8-meter class telescopes with superior spatial resolution. Such data will give the possibility of studying the mean ages and mean metallicities as functions of redshift for significantly fainter galaxies than it is currently possible by obtaining spectroscopy with 8-meter class telescopes.
Color gradients
While we have not discussed color gradients in this paper, color gradients provide an alternative method of studying galaxy evolution. The color gradients in E and S0 galaxies reflect underlying radial gradients in the metallicity (and maybe the age) of the stellar populations. Models for galaxy formation predict the sizes of these gradients. In general, the predicted gradients are steeper for models based on a monolithic collapse (Carlberg 1984) than for models based on the merger hypothesis (White 1980). Determination of color gradients for high redshift galaxies requires high signal-to-noise data with very good spatial resolution. Several of the rich galaxy clusters observed with HST/WFPC2 have sufficiently high signal-to-noise data that a study may be carried out using the available archive data.
ACKNOWLEDGEMENTS. It is a pleasure to thank the organizers for a successful and stimulating school. Support from the Nordic Research Academy (REF 99.10.003-B) for the course is kindly acknowledged, as well as support from Nato Scientific and Environmental affairs division linkage grant, Computer network supplement 97 46622 Re CRG.LG 972172 for the Internet connection.
The data used in this paper were obtained at the Nordic Optical Telescope, the Danish 1.5-meter Telescope LaSilla, the Kitt Peak National Observatory 4-meter Telescope, the Multi-Mirror Telescope, the McDonald Observatory 0.8-meter and 2.7-meter Telescopes, and the Hubble Space Telescope. We thank the telescope allocation committees for granting time to these project, and the staff at NOT, ESO/LaSilla, KPNO, MMT and McDonald Observatory for assistance during the observations.
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# Technique for measuring the parameters of polarization of an ultrasonic wave
Among physical phenomena consisting in variation of the polarization of a shear ulrasonic wave the acoustic analogs of the Faraday and the Cotton-Mouton effects are investigated at present (see and – the first theorectical papers, – discovery of rotation of the polarization, and – some experiments). They are observed when initially linearly polarized ultrasonic wave propagates inside a bulk specimen and are due to interaction between elastic and magnetic subsystems or conduction electrons. Quantitative characteristics of the effects are polarization parameters: $`\epsilon `$ – the ellipticity which modulus is the ratio of the minor and major ellipse axes and $`\varphi `$ – the angle of rotation of the polarization plane or, more correctly, of the major ellipse axis if $`\epsilon 0`$. Most of recent experiments on polarization phenomena were performed with the use of phase-amplitude methods. A review of them is given in Ref. .
Besides, a phenomenon is considered as the acoustic analog of magneto-optic Kerr effect if variation of the polarization occurs while reflection of the wave from an interface between magnetic medium and isotropic non-magnetic one. It was predicted by Vlasov and Kuleev in 1968, however, there was no papers yet about experiments in which both the parameters characterizing the polarization, $`\epsilon `$ and $`\varphi `$, were measured.
We have completed such an experiment and the results will be published soon. While performing it we found that a very small variations of a high level signal took place and came to a conclusion that amplitude variant of a technique should be more suitable here.
First amplitude technique for a precise measurement of $`\varphi `$ was introduced by Boyd and Gavenda . Its aplicability was limited to the case where $`\epsilon 0`$. Though, we developed an amplitude method free of this restriction for measuring $`\varphi `$ as well as $`\epsilon `$. A description of the technique is the subject of this paper.
The method consists of measuring the amplitude of the voltage, $`V(H)`$, on the receiving transducer at a certain $`B_1`$ relative to an initial $`B=B_0`$ using three different angles for the receiving transducer $`\psi `$ with futher processing of the data with the formulas (2), (22), and (23) presented below. It can be used for investigating the acoustic analogs of the Faraday and the Cotton-Mouton effects as well.
A periodic motion of the volume element over an elliptic trajectory can be repersented with the help of amplitudes, $`u^\pm `$, and phases, $`\phi ^\pm `$, of circular elastic vibrations. Introducing a parameter
$$p=(u^{}/u^+)e^{i(\phi ^{}\phi ^+)},$$
(1)
expressions for $`\epsilon `$ and $`\varphi `$ have the form:
$`\epsilon ={\displaystyle \frac{1|p|}{1+|p|}},\varphi ={\displaystyle \frac{1}{2}}\mathrm{Im}\left[\mathrm{ln}(p)\right].`$ (2)
Projection of the elastic vibrations to polarization direction of the receiving transducer can be written as follows:
$`u_r(t)=𝐮𝐞_r=\mathrm{Re}\{u^+\mathrm{exp}\left[i(\omega t\phi ^+\psi )\right]`$
$`+u^{}\mathrm{exp}[i(\omega t\phi ^{}+\psi )]\},`$ (3)
where * designates the complex conjugate, $`𝐞_r`$ is unit vector of the direction of the polarization of the receiving transducer, $`\psi `$ is the angle between this direction and the plane of incidence, $`\omega `$ is frequency, and $`t`$ is time. $`u_r`$ excite an ac voltage $`V\mathrm{cos}(\omega t\alpha )=\eta u_r`$ (where $`\eta ^2`$ is the coefficient of transformation of elastic vibration energy into electric field energy, and $`\alpha `$ is a phase constant). Using Eq. (3) we have
$`{\displaystyle \frac{V}{\eta }}\left[\mathrm{cos}\omega t\mathrm{cos}\alpha +\mathrm{sin}\omega t\mathrm{sin}\alpha \right]`$ (4)
$`=`$ $`\left[u^+\mathrm{cos}(\phi ^++\psi )+u^{}\mathrm{cos}(\phi ^{}\psi )\right]\mathrm{cos}\omega t`$
$`+\left[u^+\mathrm{sin}(\phi ^++\psi )+u^{}\mathrm{sin}(\phi ^{}\psi )\right]\mathrm{sin}\omega t.`$
Since Eq. (4) is valid for arbitrary $`t`$, it may be transformed into two equations:
$`{\displaystyle \frac{V}{\eta }}\mathrm{cos}\alpha `$ $`=`$ $`u^+\mathrm{cos}(\phi ^++\psi )+u^{}\mathrm{cos}(\phi ^{}\psi ),`$ (5)
$`{\displaystyle \frac{V}{\eta }}\mathrm{sin}\alpha `$ $`=`$ $`u^+\mathrm{sin}(\phi ^++\psi )+u^{}\mathrm{sin}(\phi ^{}\psi ).`$ (6)
Multiplying Eq. (6) by $`i`$ and adding the result to Eq. (5) we obtain
$`{\displaystyle \frac{V}{\eta }}e^{i\alpha }=u^+e^{i(\phi ^++\psi )}+u^{}e^{i(\phi ^{}\psi )}.`$ (7)
The method suggested here for determining the polarization of the reflected wave consists of measuring the amplitude of the signal at a certain $`B_1`$ relative to an initial $`B=B_0`$ using three different angles for the receiving transducer: $`\psi _1,\psi _2`$, and $`\psi _3`$. We assume that $`\epsilon (B_0)=0`$ and $`\varphi (B_0)=0`$. Relevant equations for the two different values of $`B`$ and three of $`\psi `$ may be obtained by making the appropriate substitutions into Eq. (7). Introducing indexes $`j=0,1`$ for the two values of $`B`$ and $`k=1,2,3`$ for the three values of $`\psi `$ for $`V_{kj}`$, $`\alpha _{kj}`$, $`u_j^\pm `$, and $`\phi _j^\pm `$ we have:
$`{\displaystyle \frac{V_{10}}{\eta }}e^{i\alpha _{10}}=u_0^+e^{i\left(\phi _0^++\psi _1\right)}+u_0^{}e^{i\left(\phi _0^{}\psi _1\right)},`$ (8)
$`{\displaystyle \frac{V_{11}}{\eta }}e^{i\alpha _{11}}=u_1^+e^{i\left(\phi _1^++\psi _1\right)}+u_1^{}e^{i\left(\phi _1^{}\psi _1\right)}.`$ (9)
Dividing Eq. (9) by (8), we obtain
$`{\displaystyle \frac{V_{11}}{V_{10}}}e^{i\left(\alpha _{11}\alpha _{10}\right)}=F_1^+e^{i\psi _1}+F_1^{}e^{i\psi _1},`$ (10)
where
$`F_1^\pm {\displaystyle \frac{u_1^\pm e^{i\phi _1^\pm }}{u_0^+\mathrm{exp}[i(\phi _0^++\psi _1)]+u_0^{}\mathrm{exp}[i(\phi _0^{}\psi _1)]}}.`$ (11)
Similar equations for $`\psi =\psi _2`$ have the form
$`{\displaystyle \frac{V_{21}}{V_{10}}}e^{i\left(\alpha _{21}\alpha _{10}\right)}\delta _2e^{i\lambda _2}=F_1^+e^{i\psi _2}+F_1^{}e^{i\psi _2},`$ (12)
where $`\lambda _2`$ and $`\delta _2`$ describe variations in phase and amplitude of the signal, respectively, caused by differences in transducer coupling to the sample while changing $`\psi `$ from $`\psi _1`$ to $`\psi _2`$.
One more change in $`\psi `$ gives the following equations in addition to (10) and (12):
$`{\displaystyle \frac{V_{31}}{V_{10}}}e^{i\left(\alpha _{31}\alpha _{10}\right)}\delta _3e^{i\lambda _3}=F_1^+e^{i\psi _3}+F_1^{}e^{i\psi _3}.`$ (13)
Here $`\delta _3`$ and $`\lambda _3`$ have the same origin as $`\delta _2`$ and $`\lambda _2`$, but correspond to changing $`\psi `$ from $`\psi _1`$ to $`\psi _3`$.
After multiplying the left and right sides of Eqs. (10), (12), and (13) by their complex conjugates we obtain
$`\left({\displaystyle \frac{V_{11}}{V_{10}}}\right)^2=\left|F_1^+\right|^2+\left|F_1^{}\right|^2+2\left|F_1^+\right|\left|F_1^{}\right|\mathrm{cos}\left(\mathrm{\Delta }\phi _1+2\psi _1\right),`$ (14)
$`\left({\displaystyle \frac{V_{21}\delta _2}{V_{10}}}\right)^2=\left|F_1^+\right|^2+\left|F_1^{}\right|^2+2\left|F_1^+\right|\left|F_1^{}\right|\mathrm{cos}\left(\mathrm{\Delta }\phi _1+2\psi _2\right),`$ (15)
$`\left({\displaystyle \frac{V_{31}\delta _3}{V_{10}}}\right)^2=\left|F_1^+\right|^2+\left|F_1^{}\right|^2+2\left|F_1^+\right|\left|F_1^{}\right|\mathrm{cos}\left(\mathrm{\Delta }\phi _1+2\psi _3\right),`$ (16)
where
$`\mathrm{\Delta }\phi _1=\phi ^+(B_1)\phi ^{}(B_1),`$ (17)
and, due to the assumption of $`\epsilon (B_0)=0`$ and $`\varphi (B_0)=0`$,
$`\delta _i={\displaystyle \frac{V_{10}\mathrm{cos}\left(\psi _i\right)}{V_{i0}\mathrm{cos}\left(\psi _1\right)}}.`$ (18)
These operations are necessary to remove the phase $`\alpha _{kj}`$ from our equations since amplitude is the only parameter measured in this variant of a technique. We divide both sides of Eqs. (14)–(16) by $`\left|F_1^+\right|\left|F_1^{}\right|`$ to obtain
$`\left|p_1\right|^1+\left|p_1\right|+2\mathrm{cos}\left[2(\varphi _1\psi _1)\right]`$ $`=`$ $`{\displaystyle \frac{\left(V_{11}/V_{10}\right)^2}{\left|F_1^+\right|\left|F_1^{}\right|}},`$ (19)
$`\left|p_1\right|^1+\left|p_1\right|+2\mathrm{cos}\left[2(\varphi _1\psi _2)\right]`$ $`=`$ $`{\displaystyle \frac{\left(V_{21}\delta _2/V_{10}\right)^2}{\left|F_1^+\right|\left|F_1^{}\right|}},`$ (20)
$`\left|p_1\right|^1+\left|p_1\right|+2\mathrm{cos}\left[2(\varphi _1\psi _3)\right]`$ $`=`$ $`{\displaystyle \frac{\left(V_{31}\delta _3/V_{10}\right)^2}{\left|F_1^+\right|\left|F_1^{}\right|}},`$ (21)
where $`p_1p(B_1)`$ .
Thus we have three equations with three unknowns, namely $`\left|F_1^+\right|\left|F_1^{}\right|`$, $`\left|p_1\right|`$, and $`\varphi _1`$. The latter two are the parameters we are interested in and corresponding solutions of the system have the form
$`\varphi _1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{tan}^1\{[(V_{21}^2\delta _2^2V_{31}^2\delta _3^2)\mathrm{cos}2\psi _1`$ (22)
$`\text{ }+(V_{31}^2\delta _3^2V_{11}^2)\mathrm{cos}2\psi _2+(V_{11}^2V_{21}^2\delta _2^2)\mathrm{cos}2\psi _3]`$
$`\times [(V_{21}^2\delta _2^2V_{31}^2\delta _3^2)\mathrm{sin}2\psi _1+(V_{31}^2\delta _3^2V_{11}^2)\mathrm{sin}2\psi _2`$
$`\text{ }+(V_{11}^2V_{21}^2\delta _2^2)\mathrm{sin}2\psi _3]^1\}`$
and
$`\left|p_1\right|={\displaystyle \frac{a_1}{c_1}}\pm \left[\left({\displaystyle \frac{a_1}{c_1}}\right)^21\right]^{1/2},`$ (23)
where
$`a_1`$ $`=`$ $`V_{11}^2\mathrm{sin}\left[2(\psi _2\psi _3)\right]+V_{21}^2\delta _2^2\mathrm{sin}[(2(\psi _3\psi _1)]`$
$`\text{ }+V_{31}^2\delta _3^2\mathrm{sin}\left[2(\psi _1\psi _2)\right]\mathrm{cos}2\varphi _1,`$
$`c_1`$ $`=`$ $`\left(V_{21}^2\delta _2^2V_{31}^2\delta _3^2\right)\mathrm{sin}2\psi _1+\left(V_{31}^2\delta _3^2V_{11}^2\right)\mathrm{sin}2\psi _2`$
$`\text{ }+\left(V_{11}^2V_{21}^2\delta _2^2\right)\mathrm{sin}2\psi _3.`$
The $`()`$ sign should be taken before the square root in Eq. (23), since it alone allows $`\left|p_1\right|=0`$ and therefore $`\epsilon =1`$.
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# Dirac monopole with Feynman brackets
## I Introduction
In 1990, Dyson published a proof due to Feynman of the Maxwell equations, assuming only commutation relations between position and velocity. In this article we don’t use the commutation relations explicitly. In fact what we call a commutation law is a structure of algebra between position and velocity called in this letter Feynman’s brackets. With this minimal assumption Feynman never supposed the existence of an Hamiltonian or Lagrangian formalism and didn’t need the not gauge invariant momentum. Tanimura extended Feynman’s derivation to the case of the relativistic particle.
In this letter one concentrates only on the following point: the study of a nonrelativistic particle using Feynman brackets. We show that Poincare’s magnetic angular momentum is the consequence of the structure of the sO(3) Lie algebra defined by Feynman’s brackets.
## II Feynman brackets
Assume a particle of mass m moving in a three dimensional Euclidean space with position: $`x_i(t)`$ ($`i=1,2,3`$) depending on time. As Feynman we consider a non associative internal structure (Feynman brackets) between the position and the velocity. The starting point is the bracket between the various components of the coordinate:
$$[x_i,x_j]=0$$
(1)
We suppose that the brackets have the same properties than in Tanimura’s article , that is:
$$[A,B]=[A,B]$$
(2)
$$[A,BC]=[A,B]C+[A,C]B$$
(3)
$$\frac{d}{dt}[A,B]=[\stackrel{.}{A},B]+[A,\stackrel{.}{B}]$$
(4)
where the arguments $`A`$, $`B`$ and $`C`$ are the positions or the velocities.
The following Jacobi identity between positions is also trivially satisfied:
$$[x_i,[x_j,x_k]]+[x_j,[x_k,x_i]]+[x_k,[x_i,x_j]]=0$$
(5)
In addition we will need also a “Jacobi identity” mixing position and velocity such that:
$$[\stackrel{.}{x_i},[\stackrel{.}{x_j},x_k]]+[\stackrel{.}{x_j},[x_k,\stackrel{.}{x_i}]]+[x_k,[\stackrel{.}{x_i},\stackrel{.}{x_j}]]=0$$
(6)
Deriving (1) gives:
$$[\stackrel{.}{x_i},x_j]+[x_i,\stackrel{.}{x_j}]=0$$
(7)
This implies:
$$[x_i,\stackrel{.}{x_j}]=g_{ij}(x_k),$$
(8)
where $`g_{ij}(x_k)`$ is a symmetric tensor. We consider here only the case where:
$$g_{ij}=\frac{\delta _{ij}}{m}$$
(9)
this gives the following relations:
$$[x_i,f(x_j)]=0$$
(10)
$$[x_i,f(x_j,\stackrel{.}{x_j})]=\frac{1}{m}\frac{f(\stackrel{.}{x_j})}{\stackrel{.}{x_i}}$$
(11)
$$[\underset{i}{\overset{.}{x}},f(x_j)]=\frac{1}{m}\frac{f(x_j)}{x_i}$$
(12)
## III Angular momentum
Suppose first the following relation:
$$[\underset{i}{\overset{.}{x}},\stackrel{.}{x_j}]=0$$
(13)
which permits to say that the force law is velocity independent:
$$\underset{i}{\overset{..}{x}}=\underset{i}{\overset{..}{x}}(x_j)$$
(14)
By definition the orbital angular momentum is:
$$L_i=m\epsilon _{ijk}x_j\stackrel{.}{x_k}$$
(15)
which satisfies the standard sO(3) Lie algebra for Feynman’s brackets:
$$[L_{i,}L_j]=\epsilon _{ijk}L_k$$
(16)
The transformation law of the position and velocity under this symmetry is:
$$[x_{i,}L_j]=\epsilon _{ijk}x_k$$
(17)
$$[\stackrel{.}{x_i},L_j]=\epsilon _{ijk}\stackrel{.}{x_k}$$
(18)
We consider as Feynman , the case with a ”gauge curvature”:
$$[\underset{i}{\overset{.}{x}},\stackrel{.}{x_j}]=\frac{\alpha }{m^2}F_{ij}$$
(19)
where $`F`$ must be an antisymmetric tensor (electromagnetic tensor for our example) and $`\alpha `$ a constant. The goal of our work is to see what happens if we keep the structure of the Lie algebra of the angular momentum and the transformation law of the position and velocity. Using (6) we get the relations:
$`\alpha {\displaystyle \frac{F_{jk}}{\underset{i}{\overset{.}{x}}}}`$ $`=`$ $`m^2[x_i[\stackrel{.}{x_j}\stackrel{.}{,x_k}]]`$ (20)
$`=`$ $`m^2[\stackrel{.}{x_j},[x_{{}_{i}{}^{},}\underset{k}{\overset{.}{x}}]]+[\underset{k}{\overset{.}{x}},[\underset{j,}{\overset{.}{x}}x_i]]=0`$ (21)
then the electromagnetic tensor is independent of the velocity:
$$F_{jk}=F_{jk}(x_i)$$
(22)
By deriving (8) we have:
$$[x_i,\underset{j}{\overset{..}{x}}]=[\underset{i}{\overset{.}{x}},\stackrel{.}{x_j}]=\frac{\alpha F_{ij}}{m^2}$$
(23)
then:
$$m\frac{\underset{j}{\overset{..}{x}}}{\stackrel{.}{x_i}}=\alpha F_{ji}(x_k)$$
(24)
or:
$$m\underset{i}{\overset{..}{x}}=\alpha (E_i(x_k)+F_{ij}(x_k)\underset{j}{\overset{.}{x}})$$
(25)
We get the ” Lorentz force’s law”, where the electric field appears as a constant of integration (this is not the case for the relativistic problem, see ). Now the force law is velocity dependent:
$$\underset{i}{\overset{..}{x}}=\underset{i}{\overset{..}{x}}(x_j,\underset{j}{\overset{.}{x}})$$
(26)
For the case (19), the equations (16), (17)and (18) become :
$$[x_{i,}L_j]=\epsilon _{ijk}x_k$$
(27)
$$[\underset{i,}{\overset{.}{x}}L_j]=\epsilon _{ijk}\underset{k}{\overset{.}{x}}+\alpha \epsilon _{jkl}x_k\frac{F_{il}}{m}$$
(28)
$$[L_{i,}L_j]=\epsilon _{ijk}L_k+\alpha \epsilon _{ikl}\epsilon _{jms}x_kx_mF_{ls}$$
(29)
Introducing the magnetic field we write $`F`$ in the following form:
$$F_{ij}=\epsilon _{ijk}B_k,$$
(30)
We get then the new relations:
$$[\underset{i,}{\overset{.}{x}}L_j]=\epsilon _{ijk}\underset{k}{\overset{.}{x}}+\frac{\alpha }{m}\{x_iB\delta _{ij}(\stackrel{}{r}.\stackrel{}{B})\}$$
(31)
$$[L_{i,}L_j]=\epsilon _{ijk}\{L_k+\alpha x_k(\stackrel{}{r}.\stackrel{}{B})\}$$
(32)
To keep the standard relations we introduce a generalized angular momentum:
$$_i=L_i+M_i$$
(33)
We call $`M_i`$ the magnetic angular momentum because it depends on the field $`\stackrel{}{B}`$. It has no connection with the spin of the particle, which can be introduced by looking at the spinorial representations of the sO(3) algebra. Now we impose for the $`\left\{\alpha _j\right\}`$’s the following commutation relations:
$$[\underset{i,}{\overset{.}{x}}_j]=\epsilon _{ijk}x_k$$
(34)
$$[\underset{i,}{\overset{.}{x}}_j]=\epsilon _{ijk}\underset{k}{\overset{.}{x}}$$
(35)
$$[_{i,}_j]=\epsilon _{ijk}_k$$
(36)
This first relation gives:
$$M_i=M_i(x_j)$$
(37)
and the second:
$$[\underset{i,}{\overset{.}{x}}M_j]=\frac{\alpha }{m}[\delta _{ij}(\stackrel{}{r}.\stackrel{}{B})x_iB_j]$$
(38)
If we replace it in (36) we deduce:
$$M_i=\alpha (\stackrel{}{r}.\stackrel{}{B})x_i$$
(39)
Putting this result in (35) gives the following equation of constraint for the field $`\stackrel{}{B}:`$
$$x_iB_j+x_jB_i=x_jx_k\frac{B_k}{x_i}$$
(40)
One solution has the form of a radial vector field centered at the origin:
$$\stackrel{}{B}=\beta \frac{\stackrel{}{r}}{r^3}$$
(41)
The generalized angular momentum then becomes:
$$\stackrel{}{}=m(\stackrel{}{r}\stackrel{\stackrel{.}{}}{r})\alpha (\stackrel{}{r}.\stackrel{}{B})\stackrel{}{r}$$
(42)
We can check the conservation of the total angular momentum:
$$\frac{d\stackrel{}{}}{dt}=m(\stackrel{}{r}\stackrel{\stackrel{..}{}}{r})\alpha \{\stackrel{}{r}(\stackrel{.}{\stackrel{}{r}}\stackrel{}{B})\}=0$$
(43)
because the particle satisfies the usual equation of motion:
$$m\frac{d^2\stackrel{\stackrel{..}{}}{r}}{dt^2}=\alpha (\stackrel{\stackrel{.}{}}{r}\stackrel{}{B})$$
(44)
If we choose: $`\alpha =q`$ and $`\beta =g`$, where $`q`$ and $`g`$ are the electric and magnetic charges, we obtain as a the special case the Poincaré magnetic angular momentum:
$$\stackrel{}{M}=\frac{qg}{4\pi }\frac{\stackrel{}{r}}{r}$$
(45)
and the Dirac magnetic monopole:
$$\stackrel{}{B}=\frac{g}{4\pi }\frac{\stackrel{}{r}}{r^3}$$
(46)
In addition we find that for the Dirac monopole the source of the field is localized at the origin:
$$div\stackrel{}{B}=[\stackrel{.}{x_i},[\stackrel{.}{x_j},\underset{k}{\overset{.}{x}}]]+[\stackrel{.}{x_j},[\underset{k}{\overset{.}{x}},\stackrel{.}{x_i}]]+[\underset{k}{\overset{.}{x}},[\stackrel{.}{x_i},\stackrel{.}{x_j}]]=\frac{g}{4\pi }[\underset{i}{\overset{.}{x}},\frac{x_i}{r^3}]=g\delta (\stackrel{}{r})$$
(47)
We see that in the construction of the Feynman’s brackets algebra the fact that we didn’t impose the Jacobi identity between the velocities is a necessary condition to obtain a monopole solution.
In summary, we used the Feynman’s algebra between position and velocity to compute the algebra of the angular momentum of a non relativistic particle in a electromagnetic field. The Dirac monopole and magnetic angular momentum is a direct consequence of the conservation of the form of the standard sO(3) Lie algebra.
## IV Casimir Operator
In the same spirit, it is interesting to introduce $`L^2,`$ the Casimir operator of sO(3) Lie algebra. Again we want to keep the same commutation relations in the two cases corresponding to zero and non zero curvature.
In the first case, we easily see that:
$$[x_{i,}L^2]=2(\stackrel{}{L}\stackrel{}{r})_i$$
(48)
$$[\underset{i,}{\overset{.}{x}}L^2]=2(\stackrel{}{L}\stackrel{\stackrel{.}{}}{r})_i$$
(49)
$$[L_{i,}L^2]=0$$
(50)
and in presence of a curvature:
$$[x_{i,}L^2]=2(\stackrel{}{L}\stackrel{}{r})_i$$
(51)
$$[\underset{i,}{\overset{.}{x}}L^2]=2[(\stackrel{}{L}\stackrel{\stackrel{.}{}}{r})_i+\alpha (\stackrel{}{L}\stackrel{}{r})_lF_{il}]$$
(52)
$$[L_{i,}L^2]=2\alpha (\stackrel{}{L}\stackrel{}{r})_i(\stackrel{}{r}.\stackrel{}{B})$$
(53)
then we want:
$$[x_{i,}^2]=2(\stackrel{}{}\stackrel{}{r})_i$$
(54)
$$[\underset{i,}{\overset{.}{x}}^2]=2(\stackrel{}{}\stackrel{\stackrel{.}{}}{r})_i$$
(55)
$$[_{i,}^2]=0$$
(56)
and we can deduce:
$$[x_{i,}M^2]=2(\stackrel{}{M}\stackrel{}{r})_i$$
(57)
$$[\underset{i,}{\overset{.}{x}}M^2]=2[(\stackrel{}{M}\stackrel{}{r})_i\alpha (\stackrel{}{L}\stackrel{}{r})_lF_{il}$$
(58)
$$2\alpha (\stackrel{}{L}\stackrel{}{r})_i(\stackrel{}{r}.\stackrel{}{B})+[L_i,M^2]+[M_{i,}L^2]=0$$
(59)
The last equation becomes after a straightforward computation:
$$(\stackrel{}{M}\stackrel{}{r})(\stackrel{}{L}\stackrel{\stackrel{.}{}}{r})(\stackrel{}{L}\stackrel{}{r})(\stackrel{}{M}\stackrel{\stackrel{.}{}}{r})(\stackrel{}{M}\stackrel{\stackrel{.}{}}{r})(\stackrel{}{L}\stackrel{}{r})+(\stackrel{}{M}\stackrel{}{r})(\stackrel{}{L}\stackrel{\stackrel{.}{}}{r})=0$$
(60)
We can check that this equation of constraint is in particular satisfied for the Poincaré angular momentum.
## V Conclusion
We find that the structure of Feynman’s brackets (without an Hamiltonian or Lagrangian), illuminates the connections between the spaces with gauge curvature, the sO(3) Lie algebra and the existence of the Poincaré magnetic angular momentum. It seems that more than the phase space formalism, the Feynman’s one is a good approach of the mechanics in a space with gauge symmetry, because it avoids the introduction of the not gauge invariant momentum. Further, other applications of this method, for example, the case of the Minkowski space with Lorentz Lie algebra, will be consider in the future.
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# The dynamics of zeros of the elliptic solutions to the Schrödinger equation
## 1 Introduction
In the recent paper it was noticed that the dynamics of zeros of $`n`$-solitonic solutions to the Schrödinger equation with the reflectionless potential is governed by the rational Ruijsenaars – Schneider system with the harmonic term . This result appears to be surprising since the aforementioned dynamics was described long ago, though in a different form. In it was shown that the Bloch solution to the Schrödinger equation
$$(_x^2u(x))\psi (x,E)=E\psi (x,E)$$
with the finite-gap potential $`u(x)`$ is a well-defined function on the hyperelliptic curve
$$y^2=\underset{i=1}{\overset{2g+1}{}}(EE_i).$$
The zeros of this function satisfy the Dubrovin equations :
$$\frac{\widehat{\gamma }_s}{x}=\frac{2y(\gamma _s)}{_{js}(\widehat{\gamma }_s\widehat{\gamma }_j)},$$
where by $`\widehat{\gamma }`$ we denote the projection of a point $`\gamma `$ on the $`E`$-plane. Notice that these equations contain parameters of the curve. An analog of the Dubrovin equations holds also for degenerate hyperelliptic curves (these are described by the same equation where not all $`E_i`$’s are distinct) and in particular for fully degenerate hyperelliptic curves which can be thought of as a Riemann sphere with $`n`$ couples of pairwise identified points. As it was shown in in the latter case parameters of the curve can be excluded from the system. The modified system then is a system of second-order differential equations written solely in terms of zeros of the corresponding function. It coincides with the Ruijsenaars-Schneider system and therefore is Hamiltonian, the expressions for the parameters of the curve being the integrals of motion.
In this paper we exploit the algebraic-geometrical approach developed in to apply these ideas to the case of the reflectionless potentials on a background of finite zone potentials, corresponding to the elliptic curves with self-intersections. The dynamics of zeros of the corresponding solutions resembles the elliptic Ruijsenaars – Schneider system . We show that the system describing these dynamics is Hamiltonian and completely integrable, the angle-type variables being the analogs of the components of the Abel map.
We hope to come up with the general system describing the case of the hyperelliptic curve with arbitrary degree of degeneracy shortly.
## 2 Algebraic-geometrical data
In this section we present some basic facts from the finite-gap theory.
Consider an elliptic curve $`\mathrm{\Gamma }`$, given by the equation
$$y^2=E^3g_2Eg_3.$$
(1)
It’s compactified at infinity by one point which we denote by $`\mathrm{}`$. The only (up to multiplication by constant) holomorphic differential on $`\mathrm{\Gamma }`$ has the following form: $`\omega ^h={\displaystyle \frac{dE}{y}}`$. It defines the map from $`\mathrm{\Gamma }`$ to the torus $`\widehat{\mathrm{\Gamma }}=/[2\omega _1,2\omega _2]`$, where $`2\omega _1`$ and $`2\omega _2`$ are $`a`$\- and $`b`$-periods of $`\omega ^h`$, respectively. This map, given by
$$A:Pz=_{\mathrm{}}^P\omega ^h$$
and known as the Abel map, allows us to identify $`\mathrm{\Gamma }`$ and $`\widehat{\mathrm{\Gamma }}`$.
Corresponding to the torus $`\widehat{\mathrm{\Gamma }}`$ are the standard Weierstrass functions
$$\sigma (z|\omega _1,\omega _2),\zeta (z|\omega _1,\omega _2)=\frac{\sigma ^{}(z|\omega _1,\omega _2)}{\sigma (z|\omega _1,\omega _2)},\mathrm{}(z|\omega _1,\omega _2)=\zeta ^{}(z|\omega _1,\omega _2)$$
(see for reference). The function $`\sigma (z)`$ has the following properties:
i) in the neighborhood of zero $`\sigma (z)=z+O(z^5);`$
ii) $`\sigma (z+2\omega _j)=e^{2\eta _j(z+\omega _j)}\sigma (z)`$, where $`\eta _j=\zeta (\omega _j)`$.
Notice that $`\mathrm{}(z)`$ is an elliptic function with the only (double) pole at $`z=0`$ and $`\zeta (z)`$ has the simple pole at $`z=0`$ and satisfies the following monodromy conditions:
$$\zeta (z+2\omega _1)=\zeta (z)+\eta _1,\zeta (z+2\omega _2)=\zeta (z)+\eta _2.$$
The map $`z(E=\mathrm{}(z),y=\mathrm{}^{}(z))`$ is inverse to the Abel map.
Let us fix $`n1`$ points $`\kappa _1,\mathrm{},\kappa _{n1}`$ on $`\widehat{\mathrm{\Gamma }}`$.
###### Proposition.
** For generic divisor $`D=\gamma _1+\mathrm{}+\gamma _n`$ on the curve $`\widehat{\mathrm{\Gamma }}`$ there exists a unique function $`\psi (x,z|D)`$ satisfying the following conditions:
1. It’s meromorphic on the curve $`\widehat{\mathrm{\Gamma }}`$ outside the point $`z=0`$ and has poles of at most first order at the points $`\gamma _i`$, $`i=1,\mathrm{},n`$.
2. In the neigborhood of $`z=0`$ it has a form
$$\psi (x,z)=e^{xz^1}(1+\underset{s=1}{\overset{\mathrm{}}{}}\xi _s(x)z^s).$$
3. $`\psi (x,\kappa _i)=\psi (x,\kappa _i).`$
###### Remark.
In general the function $`\psi (x,z|D)`$ is defined on the curve $`\mathrm{\Gamma }`$ itself, but here for the sake of brevity we use the identification between $`\mathrm{\Gamma }`$ and $`\widehat{\mathrm{\Gamma }}`$.
###### Proof.
The uniqueness of such a function follows immediately from the Riemann – Roch Theorem. To show the existence we shall consider the following function
$$\psi (x,z|D)=e^{\zeta (z)x}\frac{_{i=1}^n\sigma (zz_i(x))}{_{s=1}^n\sigma (z\gamma _s)_{i=1}^n\sigma (z_i(x))}.$$
(2)
The set of conditions $`\psi (x,\kappa _i)=\psi (x,\kappa _i)`$ and the constraint $`_{i=1}^nz_i(x)=x`$ (the latter means that $`\psi `$ is an elliptic function) form the system of $`n`$ equations on the functions $`z_i(x).`$ For generic data this system is non-degenerate. Then it has the only solution (up to the permutations) and therefore defines the function $`\psi (x,z|D)`$ uniquely. ∎
###### Corollary.
The above-constructed function $`\psi (x,z|D)`$ is a solution to the Schrödinger equation
$$(_x^2+u(x))\psi (x,z)=\mathrm{}(z)\psi (x,z),$$
(3)
where $`u(x)=2_{i=1}^n\mathrm{}(z_i(x))z_i^{}(x)`$.
###### Proof.
Consider a function $`\psi _0(x,z)=(_x^2+u(x)\mathrm{}(z))\psi (x,z)`$. It’s straightforward to check that the function $`\psi +\psi _0`$ satisfy all defining properties of the function $`\psi `$. The uniqueness of $`\psi `$ implies that $`\psi _0=0`$. ∎
## 3 Main results
###### Theorem 1.
The zeros of the function $`\psi (x,z|D)`$ satisfy the following dynamics:
$$z_i^{\prime \prime }=\underset{ki}{}z_i^{}z_k^{}\frac{\mathrm{}^{}(z_i)+\mathrm{}^{}(z_k)}{\mathrm{}(z_i)\mathrm{}(z_k)},i=1,\mathrm{},n.$$
(4)
###### Proof.
To obtain these equations one has to divide (3) by $`\psi (x,z)`$ and compare the residues of the both sides of the obtained equation at the points $`z_i(x).`$
###### Remark.
Theorem 1 provides us with a wide class of solutions to system (4) coming from the algebraic-geometrical data. The simple ”dimensional” argument shows that in fact these are all solutions. We could reverse the whole reasoning starting with the solution to (4) and showing that the corresponding elliptic function (2) solves the Shrödinger equation.
From now on we shall study system (4). Let us introduce the variables $`\xi _i=\mathrm{ln}z_i^{}`$, $`i=1,\mathrm{},n`$. In the variables $`z_i`$, $`\xi _i`$ system (4) has the following form:
$`z_i^{}`$ $`=e^{\xi _i},`$ (5)
$`\xi _i^{}`$ $`={\displaystyle \underset{ki}{}}e^{\xi _k}{\displaystyle \frac{\mathrm{}^{}(z_i)+\mathrm{}^{}(z_k)}{\mathrm{}(z_i)\mathrm{}(z_k)}},i=1,\mathrm{},n.`$
###### Proposition.
System (5) is Hamiltonian with respect to the Hamiltonian $`H={\displaystyle \underset{i=1}{\overset{n}{}}}e^{\xi _j}`$ and a 2-form
$$\omega =\underset{i=1}{\overset{n}{}}dz_id\xi _i\frac{1}{2}\underset{ij}{}\frac{\mathrm{}^{}(z_i)+\mathrm{}^{}(z_j)}{\mathrm{}(z_i)\mathrm{}(z_j)}dz_idz_j.$$
(6)
The proof is a straigtforward calculation.
Note that
$`\omega `$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}dz_id\xi _i{\displaystyle \underset{ji}{}}{\displaystyle \frac{\mathrm{}^{}(z_i)}{\mathrm{}(z_i)\mathrm{}(z_j)}}dz_idz_j=`$
$`={\displaystyle \underset{i=1}{\overset{n}{}}}dz_id\xi _i+{\displaystyle \underset{ji}{}}dz_i{\displaystyle \frac{\mathrm{}^{}(z_j)dz_j\mathrm{}^{}(z_i)dz_i}{\mathrm{}(z_j)\mathrm{}(z_i)}}=`$
$`={\displaystyle \underset{i=1}{\overset{n}{}}}dz_id\xi _i+{\displaystyle \underset{ij}{}}dz_id\left(\mathrm{ln}(\mathrm{}(z_j)\mathrm{}(z_i))\right)={\displaystyle \underset{i=1}{\overset{n}{}}}dz_id\rho _i,`$
where
$$\rho _i=\xi _i+\underset{ji}{}\mathrm{ln}(\mathrm{}(z_j)\mathrm{}(z_i)).$$
The algebraic-geometrical construction from the previous section provides us with a hint on how the first integrals of system (4) should look like. The constraints $`\psi (x,\kappa _s)=\psi (x,\kappa _s)`$ imply the equations
$$\underset{j=1}{\overset{n}{}}\frac{z_j^{}}{\mathrm{}(\kappa _s)\mathrm{}(z_j)}=0,$$
which can be rewritten in the following form
$$\underset{i=1}{\overset{n}{}}z_i^{}\underset{ji}{}(\mathrm{}(z_j)\mathrm{}(\kappa _s))=0.$$
These considerations motivate
###### Theorem 2.
The coefficients $`H_k`$ of the polynomial
$$L(\lambda |z,z^{})=\underset{k=0}{\overset{n1}{}}H_k(z,z^{})\lambda ^k=\underset{i=1}{\overset{n}{}}z_i^{}\underset{ji}{}(\mathrm{}(z_j)\lambda )$$
(7)
are the integrals of motion of system (4).
###### Remark.
Note that the leading cofficient $`H_{n1}(z,z^{})`$ of $`L`$ is equal up to the sign to the Hamiltonian $`H(z,z^{})`$ of system (4).
The statement of the theorem is clear since we know that all solutions are algebraic-geometrical. However, we would like to present an independent direct proof. It can be found in the Appendix I.
Let us notice that $`L(\mathrm{}(z_j))=e^{\rho _j}`$. Using this identity we can rewrite the form $`\omega `$ in the following way:
$`\omega `$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}dz_id\rho _i={\displaystyle \underset{i=1}{\overset{n}{}}}dz_id\mathrm{ln}L(\mathrm{}(z_i))=`$
$`={\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{1}{L(\mathrm{}(z_i))}}dz_id\left({\displaystyle \underset{s=0}{\overset{n1}{}}}H_s\mathrm{}^s(z_i)\right)={\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \underset{s=0}{\overset{n1}{}}}{\displaystyle \frac{\mathrm{}^s(z_i)}{L(\mathrm{}(z_i))}}dz_idH_s=`$
$`={\displaystyle \underset{s=0}{\overset{n1}{}}}d\left({\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \stackrel{\mathrm{}(z_i)}{}}{\displaystyle \frac{E^sdE}{L(E)y(E)}}\right)dH_s+{\displaystyle \underset{s,k=0}{\overset{n1}{}}}\left({\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \stackrel{\mathrm{}(z_i)}{}}{\displaystyle \frac{E^{s+k}dE}{L(E)^2y(E)}}\right)dH_kdH_s=`$
$`={\displaystyle \underset{s=0}{\overset{n1}{}}}d\left({\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \stackrel{\mathrm{}(z_i)}{}}{\displaystyle \frac{E^sdE}{L(E)y(E)}}\right)dH_s,`$
where the function $`y(E)`$ is given by (1).
Thus we have proved the following statement.
###### Theorem 3.
The variables
$$\phi _s=\underset{i=1}{\overset{n}{}}\stackrel{\mathrm{}(z_i)}{}\frac{E^sdE}{L(E)y(E)},s=0,\mathrm{},n1$$
and $`H_s`$ defined by (7) are the action-angle type variables for system (4).
We would like however to rewrite the form $`\omega `$ once again in terms of the zeros of the polynomial $`L(\lambda |z,z^{})`$ which we shall denote by $`\widehat{\kappa }_j`$, $`j=1,\mathrm{},n1`$. In order to do this we introduce the new variables
$$\chi _j=\underset{i=1}{\overset{n}{}}\stackrel{\mathrm{}(z_i)}{}\frac{dE}{(E\widehat{\kappa }_j)y(E)},j=1,\mathrm{},n1.$$
Let us also introduce the variable
$$\chi =\underset{i=1}{\overset{n}{}}\stackrel{\mathrm{}(z_i)}{}\frac{dE}{y(E)}.$$
###### Theorem 4.
The above-defined form $`\omega `$ admits the following representation
$$\omega =d\chi d(\mathrm{ln}H)+\underset{j=1}{\overset{n1}{}}d\chi _jd\widehat{\kappa }_j.$$
(8)
###### Remark.
We want to emphasize the fact that the variables $`\{\chi ,\chi _j,j=1,\mathrm{},n1\}`$ are the degenerate curve analogs of the components of the Abel map. So our Hamiltonian structure fits in the general scheme proposed in and developed in .
The proof is a straightforward computation (see Appendix II).
## Appendix I
Let us consider the polynomial
$$L(\lambda |z,z^{})=\underset{j=1}{\overset{n}{}}z_j^{}\underset{ij}{}(\mathrm{}(z_i)\lambda )=\underset{k=0}{\overset{n1}{}}H_k(z,z^{})\lambda ^k.$$
The explicit formulae for the coefficients $`H_k`$ are
$$H_k=\underset{|J|=nk1}{}(1)^k\underset{jJ}{}\mathrm{}(z_j)\left(\underset{kJ}{}z_k^{}\right),$$
where summation is taken over all subsets $`J\{1,\mathrm{},n\}`$ of cardinality $`nk1`$.
We are going to show that the functions $`H_k`$ are time-independent, i. e. $`dH_k/dx=0`$. Indeed,
$$\begin{array}{c}\frac{d(1)^kH_k(z,z^{})}{dx}=\underset{J}{}\underset{jJ}{}\mathrm{}(z_i)\underset{kJ}{}z_k^{\prime \prime }+\underset{J}{}\left(\underset{sJ}{}\frac{\mathrm{}^{}(z_s)}{\mathrm{}(z_s)}z_s^{}\right)\underset{jJ}{}\mathrm{}(z_i)\underset{kJ}{}z_k^{\prime \prime }=\hfill \\ \hfill =\underset{J}{}\underset{jJ}{}\mathrm{}(z_j)\left(\underset{kJ}{}\underset{ik}{}\frac{\mathrm{}^{}(z_k)+\mathrm{}^{}(z_i)}{\mathrm{}(z_k)\mathrm{}(z_i)}z_i^{}z_k^{}\right)+\underset{J}{}\underset{jJ}{}\mathrm{}(z_j)\left(\underset{kJ}{}\underset{sJ}{}\frac{\mathrm{}^{}(z_s)}{\mathrm{}(z_s)}z_k^{}z_s^{}\right)=\\ \hfill =\underset{J}{}\underset{jJ}{}\mathrm{}(z_j)\left(\underset{kJ}{}\underset{sJ}{}\left[\frac{\mathrm{}^{}(z_s)}{\mathrm{}(z_s)}+\frac{\mathrm{}^{}(z_k)+\mathrm{}^{}(z_s)}{\mathrm{}(z_k)\mathrm{}(z_s)}\right]z_k^{}z_s^{}\right)=\underset{J,kJ,sJ}{}\alpha (J,k,s),\end{array}$$
where
$$\alpha (J,k,s)=\underset{jJ}{}\mathrm{}(z_j)\left[\frac{\mathrm{}^{}(z_s)}{\mathrm{}(z_s)}+\frac{\mathrm{}^{}(z_k)+\mathrm{}^{}(z_s)}{\mathrm{}(z_k)\mathrm{}(z_s)}\right]z_k^{}z_s^{}.$$
Let us consider the involution on the set of triples $`\{J,kJ,sJ\}`$ which maps $`\{J,k,s\}`$ into $`\{J^{},s,k\}`$, where $`J^{}=J\{k\}\{s\}`$.
Now note that
$$\begin{array}{c}\alpha (J,k,s)+\alpha (J^{},s,k)=\underset{jJJ^{}}{}\mathrm{}(z_j)z_k^{}z_s^{}[\mathrm{}^{}(z_s)+\mathrm{}^{}(z_k).+\hfill \\ \hfill +\mathrm{}(z_s)\frac{\mathrm{}^{}(z_k)+\mathrm{}^{}(z_s)}{\mathrm{}(z_k)\mathrm{}(z_s)}+\mathrm{}(z_k)\frac{\mathrm{}^{}(z_s)+\mathrm{}^{}(z_k)}{\mathrm{}(z_s)\mathrm{}(z_k)}]=0\end{array}$$
and therefore the whole sum $`_{J,kJ,sJ}\alpha (J,k,s)`$ vanishes.
## Appendix II
Consider the 2-form $`\omega =_{s=0}^{n1}d\phi _sdH_s`$. Recall that $`H_s=(1)^sH\sigma _{ns1}(\widehat{\kappa })`$, $`s=0,\mathrm{},n1`$, where $`\sigma _{ns1}(\widehat{\kappa })`$ denotes the coefficient of $`\lambda ^s`$ in the polynomial $`_{i=1}^{n1}(\lambda +\widehat{\kappa }_i)`$. By $`\sigma _{ns2}^j(\widehat{\kappa })`$ we denote the coefficient of $`\lambda ^s`$ in the polynomial $`_{ij}(\lambda +\widehat{\kappa }_i)`$. Then
$$\begin{array}{c}\omega =\underset{s=0}{\overset{n1}{}}d\phi _sdH_s=\underset{s=0}{\overset{n2}{}}d\phi _s(1)^sd(H\sigma _{ns1}(\widehat{\kappa }))+(1)^{n1}d\phi _{n1}dH=\hfill \\ \hfill =\underset{s=0}{\overset{n1}{}}(1)^s\sigma _{ns1}(\widehat{\kappa })d\phi _sdH+\underset{j=1}{\overset{n1}{}}\underset{s=0}{\overset{n2}{}}(1)^sH\sigma _{ns2}^j(\widehat{\kappa })d\phi _sd\widehat{\kappa }_j.\end{array}$$
(9)
Now let us notice that
$$\begin{array}{c}\underset{s=0}{\overset{n2}{}}(1)^sH\sigma _{ns2}^j(\widehat{\kappa })d\phi _s=\underset{s=0}{\overset{n2}{}}\underset{l=1}{\overset{n}{}}(1)^sH\sigma _{ns2}^j(\widehat{\kappa })d\stackrel{\mathrm{}(z_l)}{}\frac{E^sdE}{L(E)y(E)}=\hfill \\ \hfill =\underset{l=1}{\overset{n}{}}d\stackrel{\mathrm{}(z_l)}{}\frac{\underset{s=0}{\overset{n2}{}}(1)^sH\sigma _{ns2}^j(\widehat{\kappa })E^s}{L(E)y(E)}𝑑E\underset{l=1}{\overset{n}{}}\stackrel{\mathrm{}(z_l)}{}d\left(\frac{\underset{s=0}{\overset{n2}{}}(1)^sH\sigma _{ns2}^j(\widehat{\kappa })E^s}{L(E)y(E)}\right)𝑑E=\\ \hfill =\underset{l=1}{\overset{n}{}}d\stackrel{\mathrm{}(z_l)}{}\frac{dE}{(E\widehat{\kappa }_j)y(E)}+\underset{l=1}{\overset{n}{}}\stackrel{\mathrm{}(z_l)}{}\frac{dE}{(E\widehat{\kappa }_j)^2y(E)}𝑑\widehat{\kappa }_j.\end{array}$$
In the same way one can show that
$$\underset{s=0}{\overset{n1}{}}(1)^s\sigma _{ns1}(\widehat{\kappa })d\phi _s=\underset{l=1}{\overset{n}{}}d\stackrel{\mathrm{}(z_l)}{}\frac{dE}{Hy(E)}+\underset{l=1}{\overset{n}{}}\stackrel{\mathrm{}(z_l)}{}\frac{dE}{H^2y(E)}𝑑H.$$
Plugging these two formulae into (9) we obtain
$$\begin{array}{c}d\omega =\underset{j=1}{\overset{n1}{}}d\left(\underset{l=1}{\overset{n}{}}^{\mathrm{}(z_l)}\frac{dE}{(E\widehat{\kappa }_j)y(E)}\right)d\widehat{\kappa }_j+d\left(\underset{l=1}{\overset{n}{}}d^{\mathrm{}(z_l)}\frac{dE}{Hy(E)}\right)dH=\hfill \\ \hfill =\underset{j=1}{\overset{n1}{}}d\chi _jd\widehat{\kappa }_j+d\chi d(\mathrm{ln}H).\end{array}$$
### Acknowledgements
The authors are grateful to Professor I. M. Krichever for constant attention to this work.
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# Mixing Metals in the Early Universe
## 1 Introduction
Primordial nucleosynthesis enriched gas in the universe with the light elements He, D, and Li. It is only when the first galaxies and their stars appeared that heavier elements could be synthesized and, in some cases, ejected into the intergalactic medium. In currently popular models of galaxy formation based on hierarchical clustering the first galaxies to form were low mass systems with such shallow potential wells that a few supernovae could deposit sufficient kinetic energy to expel the entire interstellar medium of the galaxy (Ferrara 1998). In this way, these initial episodes of star formation in the universe (sometimes referred to as Population III) may have enriched the intergalactic medium (IGM) to an average metallicity $`Z_{\mathrm{IGM}}10^4Z_{}`$ (Miralda-Escudé & Rees 1997, Nath & Trentham 1997, Gnedin & Ostriker 1997, Ciardi et al. 2000b), comparable to that of the most metal-poor stars in the halo of our Galaxy (Ryan, Norris, & Beers 1996).
Most of the metals in the universe were presumably produced in larger collapsed structures at redshifts $`z\stackrel{<}{}10`$. In general such galaxies were better able to retain the products of stellar nucleosynthesis and could therefore achieve the abundance levels observed today. It is also likely that some metal-enriched gas escaped into the IGM, but it is far from clear by which process and to what degree the metals became distributed over large volumes, far from their production sites. These are the questions we consider in the present paper. Recent numerical hydrodynamic simulations of large-scale structure formation (Hellsten et al. 1997; Rauch et al. 1997; Zhang et al. 1998) have shown that baryons in the universe are distributed in a network structure (the “cosmic web”) with galaxies located at the high overdensity peaks. This implies that most of the volume is occupied by voids filled with gas at, or below, the mean cosmic density; such gas can be viewed as the long-searched-for true intergalactic medium. Since the voids are very large, with typical dimensions of several Mpc, their pollution by heavy elements produced in SN explosions is far from a trivial problem. In any case there must have been a protracted epoch when the metal content of the IGM was highly patchy.
Observationally, our best view of the IGM is still provided by the Ly$`\alpha `$ forest in the spectra of distant QSOs (Sargent et al. 1980) and the detection of metals in the forest ranks as one of the most significant discoveries made possible by the Keck telescopes (Cowie et al. 1995; Tytler et al. 1995). From an analysis of these data Hellsten et al. (1997) and Rauch, Haehnelt, & Steinmetz (1997) concluded that the measured column density ratios $`N`$(C IV)/$`N`$(H I) imply that typically \[C/H\] $`2.5`$ at $`z3`$, with a one order of magnitude dispersion in the metallicity of different clouds about this mean value.<sup>1</sup><sup>1</sup>1In the usual notation, \[C/H\] = log (C/H) $``$ log (C/H). At lower redshifts ($`z=0.30.8`$) Barlow & Tytler (1998) using HST/FOS spectra, with reasonable assumptions concerning the ionization correction and the line clustering properties, conclude that metallicities are as high as \[C/H\]$`\stackrel{>}{}1.3`$, roughly an order of magnitude larger than the value at $`z=2.5`$. These measurements, however, still refer to overdense regions of the universe, traced by Ly$`\alpha `$ clouds with column densities in excess of log $`N`$(H I) = 14.5 . The situation is far less clear-cut when we turn to the voids—or log $`N`$(H I)$`<14.0`$—where the observations are very challenging even with a 10-m telescope. Two studies have addressed this problem, with conflicting results. Lu et al. (1998) applied a stacking technique to nearly 300 C IV regions in QSO spectra but still found no composite signal. They interpreted this non-detection as evidence for a highly non-uniform degree of metal enrichment at $`z3`$, with the voids having \[C/H\] $`\stackrel{<}{}3.5`$ . On the other hand, Cowie & Songaila (1998) used a pixel-to-pixel optical depth technique to conclude that the average C IV/H I ratio remains essentially constant over the full range of neutral hydrogen column densities tested, down to log $`N`$(H I) $`13.5`$ . If this is indeed the case, the ejection and transport of metals away from galaxies must have been much more efficient than envisaged. Most recently, Ellison et al. (1999, 2000) have re-examined the problem and found that both approaches suffer from limitations which had not been properly taken into account in previous analyses. In their view, whether there are metals in the voids remains an open question.
On the theoretical side too, only a limited amount of work has been done on this topic. The common assumption that supernova driven winds may be able to distribute metals over large distances has been shown to be too simplistic and does not stand up to close quantitative scrutiny (MacLow & Ferrara 1999; Murakami & Babul 1999; D’Ercole & Brighenti 1999). These studies have come to the conclusion that efficient blowout is likely to be inhibited by the galaxy ISM and, at least at high redshift where densities are higher, by the pressure of the surrounding intergalactic gas (Babul & Rees 1992; Ciardi & Ferrara 1997). Gnedin & Ostriker (1998) numerically simulated the IGM enrichment in a CDM+$`\mathrm{\Lambda }`$ cosmological model. They found that the mean metallicity of the universe at $`z=4`$ is about 1/200 $`Z_{}`$ but the variations around this value are large and dependent on the overdensity. At this epoch the enrichment was incomplete, with some regions of the universe still of pristine composition. Their simulations, however, end at $`z=4`$ so that a direct connection with the nearby universe is difficult. In addition, quantitative conclusions on the transport mechanism are affected by uncertainties arising from insufficient numerical resolution so that, for example, they are unable to resolve the snowplough phase of explosions. To alleviate the problem of inefficient blowout Gnedin (1998) developed the idea that metals are predominantly transported following merger events, a mechanism which would also predict a highly inhomogeneous distribution of heavy elements in the IGM at $`z=4`$. However, the merger history has not yet been followed to more recent epochs.
In this paper we make explicit predictions for the IGM metallicity evolution by suggesting a two-step mechanism for the transport and mixing of heavy elements. In our scenario, SN explosions first eject metals over a relatively small region in the surroundings of the host galaxy; subsequently, a yet unknown diffusion process transports and mixes these elements over much larger scales. The plan of the paper is as follows. In §2 we consider the conditions for gas ejection from galaxies, paying particular attention to the treatment of blowout, which we show to be quite a rare occurrence in massive galaxies. In §3 we explicitly derive the metallicity evolution of the IGM for a particular (Cold Dark Matter) cosmological model. The results are discussed in §4 and a brief summary (§5) concludes the paper.
## 2 Metal Ejection by Galaxies
It is likely that the first episodes of star formation in the universe had a dramatic impact on the small mass galaxies in which they occurred. The combined effects of supernova explosions in relatively small volumes lead to the formation of superbubbles (SBs) which can result in partial (blowout) or even complete (blowaway) removal of the interstellar medium (Kovalenko & Shchekinov 1985; MacLow & McCray 1988; Ciardi & Ferrara 1997; Ferrara 1998; MacLow & Ferrara 1999). The resulting large scale outflows (galactic superwinds) have been observed directly in local starbursts (González Delgado et al. 1998; Heckman et al. 1998) and in Lyman break galaxies at $`z3`$ (Lowenthal et al. 1997; Franx et al. 1997; Pettini et al. 1998; Pettini et al. 2000). For the purpose of this paper, we are interested in establishing the maximum galactic mass which allows blowout to occur. MacLow & Ferrara (1999) have shown that when blowout does take place the escape efficiency of the metals produced by the supernova progenitors is close to unity.
### 2.1 Conditions for Blowout
The blowout condition can be derived by comparing two characteristic velocities: the blowout velocity, $`v_b`$, and the escape velocity of the galaxy, $`v_e`$. To calculate these two velocities we start by defining a protogalaxy as a two component system consisting of a dark matter halo and a gaseous disk. We assume a modified isothermal halo density profile $`\rho _h(r)=\rho _c/[1+(r/r_a)^2]`$ extending out to a radius $`r_{200}r_h=[3M_h/4\pi (200\rho _{crit})]^{1/3}`$, defined as the radius within which the mean dark matter density is 200 times the critical density $`\rho _{crit}=3H_0^2(1+z)^3/8\pi G`$ at the redshift $`z`$ when the halo is identified. $`M_h`$ is the halo mass and $`H_0=100h`$ km s<sup>-1</sup> Mpc<sup>-1</sup> is the present-day Hubble constant; throughout this paper we adopt $`\mathrm{\Omega }=1`$.
For such a halo the escape velocity can be written as
$$v_e^24\pi pG\rho _cr_a^2=\frac{2pGM_h}{r_h},$$
(1)
with $`p=1.65`$ (MacLow & Ferrara 1999). Note that $`v_e=\sqrt{2p}v_c`$, where $`v_c`$ is the circular velocity of the halo.
Due to the dissipative nature of the gas, the baryons which initially should be distributed approximately as the dark matter, will lose pressure and collapse in the gravitational field of the dark matter. The density ratio between these two components is assumed here to be equal to its cosmological value, and therefore the initial density of the gas in the protogalaxy is $`\rho _g=\mathrm{\Omega }_b\rho _h`$. If the halo is rotating, the gas will collapse in a centrifugally supported disk. The radius of the disk can be estimated by imposing that the specific angular momentum of the disk, $`j_d`$, is equal to that of the halo, $`j_h`$ (Mo et al. 1998, Weil et al. 1998)
$$j_h=\sqrt{2}\lambda v_cr_h=2v_c\mathrm{}_d=j_d$$
(2)
where $`\mathrm{}_d`$ is the disk scale length and $`\lambda `$ the standard halo spin parameter. As shown by numerical simulations (Barnes & Efstathiou 1987; Steinmetz & Bartelmann 1995) $`\lambda `$ depends very weakly on $`M_h`$ and on the density fluctuation spectrum; its distribution is approximately log-normal and peaks around $`\lambda =0.04`$. From eq. 2 we then obtain $`\mathrm{}_d=(\lambda /\sqrt{2})r_h`$. For an exponential disk, implicitly assumed in deriving eq. 2 above, the optical radius (i.e. the radius encompassing 83% of the total integrated light) is $`3.2\times \mathrm{}_d`$. Since galaxies typically extend $`2`$ times their optical radius in H I (Salpeter & Hoffman 1996), we adopt for the radius of the gaseous disk
$$r_d=4.5\lambda r_h.$$
(3)
We now evaluate the scale height of the gas in the disk. This is roughly given by
$$H\frac{c_s^2r_d^2}{GM_h}70\lambda ^2\left(\frac{c_s}{v_e}\right)^2r_h$$
(4)
where we have used the expression for $`r_d`$ above; $`c_s`$ is the effective gas sound speed which also includes a possible turbulent contribution. Note that
$$\frac{H}{r_d}=15.3\lambda \left(\frac{c_s}{v_e}\right)^2$$
(5)
If all baryons bound to the dark matter halo were able to collapse, the mean gas density in the disk would be
$$\rho _d\left(\frac{r_h}{r_d}\right)^3\left(\frac{r_d}{H}\right)\rho _g.$$
(6)
The previous assumption, although uncertain, gives approximately the correct density when scaled to Milky Way parameters; in addition, $`\rho _d`$ enters in the expression for $`v_b`$ below to the 1/3 power. Thus, unless the fraction of collapsed baryons is unreasonably low, a value lower than unity does not lead to qualitatively different conclusions. The explicit expression for the blowout velocity, $`v_b`$, has been obtained by Ferrara & Tolstoy (2000):
$$v_b=2.7\left(\frac{L}{\rho _dH^2}\right)^{1/3},$$
(7)
where $`L`$ is the mechanical luminosity of the parent SB. Note that $`v_b`$ by definition (Ferrara & Tolstoy 2000) is estimated at a height equal to 3H above the galactic plane. Obviously, if the blowout velocity has already decreased below the local sound speed at that point, the blowout is inhibited. By using eqs. 3, 6 and 4 it is easy to show that $`H^2\rho _d\mathrm{\Omega }_bc_s^2`$, that is the product $`H^2\rho _d`$ is independent of mass and redshift; therefore $`v_b`$ depends on these quantities only through $`L`$. Blowout occurs if $`v_b>v_e`$, or
$$2.7\left(\frac{L}{H^2\rho _d}\right)^{1/3}\left(\frac{2pGM_h}{r_h}\right)^{1/2}.$$
(8)
This implies that the mechanical luminosity for a SB to blowout must be larger than the critical value
$$L_c=(0.27pG)^{3/2}\rho _dH^2\left(\frac{M_h}{r_h}\right)^{3/2}.$$
(9)
The total mechanical luminosity of the corresponding SB in a primordial galaxy that virialized at redshift $`z`$ can be written as (Ferrara 1998)
$$L_t(z)=ϵ_0\nu \dot{M}_{}=\frac{ϵ_0\nu \mathrm{\Omega }_bf_b}{\tau _{}t_{ff}(z)}M_h,$$
(10)
where $`ϵ_0=10^{51}`$ ergs is the energy of a SN explosion and $`\dot{M}_{}`$ is the star formation rate; we assume a Salpeter IMF, for which one SN is produced for each 56 $`M_{}=\nu ^1`$ of stars formed. The baryon density parameter is $`\mathrm{\Omega }_b=0.05\mathrm{\Omega }_{b,5}`$, of which a fraction $`f_b0.08f_{b,8}`$ (Ciardi et al. 2000a) is able to cool and become available to form stars. The free-fall time is $`t_{ff}=(4\pi G\rho _h)^{1/2}`$; $`\tau _{}^1=3\%`$ is the star formation efficiency, which is fixed by matching to the cosmic star formation history as described in §4.2 below.
For Population III objects, which as we have argued are likely to have been of low mass, we make the simplifying assumption that star formation is confined to a relatively small region, and that all the exploding SNae associated with the initial star formation episode result in the formation of a single superbubble. In this case a comparison between eq. 9 and eq. 10 shows that $`L_t(z)/L_c`$ does not depend on mass and redshift, and the condition $`L_t(z)>L_c`$ is satisfied when
$$\frac{ϵ_0\nu f_b}{\tau }>0.1c_s^2,$$
(11)
which for typical parameters is equivalent to $`4.4\times 10^{12}>0.1c_s^2`$, and is therefore valid for galaxies with $`c_s70`$ km/s. Thus low mass galaxies with coherent star formation (i.e. with all SNae driving a single SB) are very likely blowing out.
For more massive and larger galaxies we have to consider the more realistic situation where the SNae are more widely distributed within the disk, and occur in OB associations with different values of $`L_{OB}`$ (or, equivalently, with a different number $`N`$ of SNae). In nearby galaxies it is found that the luminosity function of OB associations is well approximated by a power-law
$$\varphi (N)=\frac{d𝒩_{OB}}{dN}=AN^\beta $$
(12)
with $`\beta 2`$ (McKee & Williams 1997; Oey & Clarke 1998). Here $`𝒩_{OB}`$ is the number of associations containing $`N`$ OB stars; normalization of $`\varphi (N)`$ to unity requires $`A=1`$. Thus the probability for a cluster of OB stars to host $`N`$ SNae is $`N^2`$, where $`N=L_{OB}t_{OB}/ϵ_0`$, and $`t_{OB}=40`$ Myr is the time at which the lowest mass ($`8M_{}`$) SN progenitors expire. The total mechanical luminosity, which must be equal to $`L_t(z)`$ in eq. 10, is then found to be
$$L_t(z)=\underset{N_m}{\overset{N_M}{}}L_{OB}(N)\varphi 𝑑N,$$
(13)
where $`N_m=1`$ ($`N_M`$) is the minimum (maximum) possible number of SNae in a cluster. This gives
$$L_t(z)=\mathrm{const}.\frac{ϵ_0}{t_{OB}}\mathrm{ln}\frac{N_M}{N_m}.$$
(14)
The contribution to the total luminosity from clusters powerful enough to lead to blowout is
$$L_B(z,>L_c)=\mathrm{const}.\frac{ϵ_0}{t_{OB}}\mathrm{ln}\frac{N_M}{N_c},$$
(15)
where $`N_c`$ is the number of SNae in a cluster with mechanical luminosity equal to $`L_c`$, i.e.
$$N_c=\frac{L_ct_{OB}}{ϵ_0}.$$
(16)
Thus, the fraction of the mechanical energy which can be blown out is
$$\delta _B=\frac{\mathrm{ln}(N_M/N_c)}{\mathrm{ln}(N_M/N_m)}<1.$$
(17)
Clearly, $`N_M`$ (and therefore $`\delta _B`$) is an intrinsically stochastic number. To determine its dependence on the total number of supernovae $`N_t=L_t(z)t_{OB}/ϵ_0`$ produced by a galaxy during the lifetime of an OB association, we have used a Monte Carlo procedure applied to the distribution function in eq. 12. The results for $`N_M`$ and $`\delta _B`$ as a function of $`N_t`$ are shown in Fig. 1; we recall that the star formation rate $`\dot{M}_{}`$ in the galaxy is related (through eq. 10) to $`N_t`$ by $`\dot{M}_{}5\times 10^6N_tM_{}`$ yr<sup>-1</sup>. As can be seen from Fig. 1, for low values of $`N_t`$ the quantity $`N_M`$ is larger than $`N_c`$, implying that in every galaxy at least some SBs are able to blow out. However, near $`N_t=10^4`$ $`N_M`$ flattens and eventually becomes equal to $`N_c`$ at $`N_t\mathrm{45\hspace{0.17em}000}`$. Above this limit (corresponding to a galaxy with $`\dot{M}_{}0.35M_{}`$ yr<sup>-1</sup> or $`M_h10^{12}(1+z)^{3/2}M_{}`$) blowout is inhibited. The fraction $`\delta _B`$ can be seen from Fig. 1 to be a decreasing function of $`N_t`$; an approximate analytical form is
$`\delta _B(N_t)=1\mathrm{for}N_t<100`$ (18)
$`\delta _B(N_t)=a+b\mathrm{ln}(N_t^1)\mathrm{for}N_t>100,`$
with $`a=1.76`$, $`b=0.165`$. Clearly, in small galaxies even the smallest associations are capable of producing blowout, so that the issue of coherence discussed above is irrelevant.
Thus, if we assume the most of the baryonic matter has collapsed into the disk, the main metal polluters of the IGM at the present epoch are galaxies with visible mass lower than $`M_d^{up}=\mathrm{\Omega }_bM_h\stackrel{<}{}5\times 10^{10}M_{}`$. This critical mass decreases with increasing redshift as $`(1+z)^{3/2}`$, indicating that the initial metal enrichment of the universe must have been produced by even smaller galaxies. For example, at $`z=20`$, the expected value of $`M_d^{up}`$ is only $`5\times 10^8M_{}`$.
### 2.2 Confinement of Blowout-driven Outflows
Blowout will drive an outflow which will eventually be confined by the IGM pressure. What is the characteristic length, $`R_e`$, at which such pressure equilibrium is achieved? By requiring that the outflow ram pressure, $`\rho _wv_w^2`$, is equal to the IGM pressure, $`p_i`$ we obtain
$$R_e=\left(\frac{\dot{M}_wv_w}{4\pi p_i}\right)^{1/2};$$
(19)
where we have used mass conservation and assumed that the flow is approximately spherical. Once recombination is complete at redshift $`z200`$, and before the IGM is reheated by the energy input from the first galaxies, the IGM pressure will evolve almost adiabatically, $`p_i=p_{}[(1+z)/200)]^5`$, where $`p_{}/k_B1500h^2`$ cm<sup>-3</sup> K is the value of $`p_i`$ at $`1+z=200`$. However, almost unavoidably, even the earliest SN-driven bubbles will expand in an IGM which has been pre-ionized by the same massive stars which later exploded as SNae. Thus, it seems more appropriate to calculate the IGM pressure in the surroundings of the galaxy as $`p_i=(R/\mu )\rho _i(z)T`$, with $`T2\times 10^4`$ K as a result of photoionization heating. This is a reasonable hypothesis since it can be shown Ciardi et al. 2000a) that the ionization spheres are much larger than the metal-enriched bubbles considered here. The mass loss rate $`\dot{M}_w`$ is typically a fraction $`\xi 10`$% of the star formation rate $`\dot{M}_{}=L_e/ϵ_0\nu `$ (Ferrara & Tolstoy 2000) where $`L_e=\delta _BL_t`$ is the effective mechanical luminosity, that is the fraction available for blowout; we can also assume that $`v_wv_e`$. With these hypotheses, and using the relations derived above for $`v_e`$, one obtains
$$R_e(M_h,z)=0.9\left(\frac{\mathrm{\Omega }_{b,5}}{\lambda _4}\frac{v_e^4\delta _B}{p_i}\right)^{1/2}\mathrm{cm},$$
(20)
where $`\lambda _4=\lambda /0.04`$. The behavior of $`R_e`$ with $`M_h`$ is shown in Fig. 2 for different redshifts. At lower redshifts, the bubbles are larger because the pressure of the confining IGM is lower; also evident from the figure is the increase of the critical mass for blowout with decreasing $`z`$.
The major conclusion that can be drawn from Fig. 2 is that supernova-driven outflows are quite inadequate for dispersing heavy elements far from their production sites and cannot account, by themselves, for the ubiquitous presence of metals in the IGM at $`z3`$. This can be readily realized when we consider that in standard Cold Dark Matter models, taken here as representative of a larger class of hierarchical models of structure formation, the typical (physical) separation between objects with mass $`M_h10^{10}M_{}`$ varies between $`20.2`$ Mpc in the interval $`0<z<8`$ (Ciardi et al. 2000a). This scale is much larger than the values of $`R_e`$ in Fig. 2, so that we would expect regions outside the metal enriched bubbles to retain their primordial composition (this argument is reconsidered more rigorously in §4 below). Clearly some other mechanism must be operating to remove the metals from the surroundings of galaxies and mix them with the more generally distributed IGM.
## 3 Predictions for CDM Models
To make further progress, we need to fix a specific cosmological model. As an example, we consider the so-called Standard Cold Dark Matter (SCDM), with $`\mathrm{\Omega }_M=1,\mathrm{\Omega }_\mathrm{\Lambda }=0`$, and $`h=0.5`$; the power spectrum $`|\delta _k|^2`$ is taken from Efstathiou et al. (1992), normalized to the present-day abundance of rich clusters ($`\sigma _8=0.6`$; Eke, Cole, & Frenk 1996). To calculate the number density of dark matter halos as a function of redshift we use the Press & Schechter (1974, hereafter PS) formalism; this technique is widely used in semi-analytical models of galaxy formation and gravitational lensing (White & Frenk 1991, Kauffman 1995, Ciardi & Ferrara 1997, Baugh et al. 1998, Guiderdoni et al. 1998, Marri & Ferrara 1998) and it has been shown to be in surprisingly good agreement with the results from N-body numerical simulations. Given a power spectrum $`|\delta _k|^2`$, one can write the Gaussian variance of the fluctuations on the mass scale $`M`$:
$$\sigma _M^2=\frac{d^3k}{(2\pi )^3}W^2(k,R)|\delta _k|^2,$$
(21)
where $`M=(4/3)\pi \rho R^3`$, $`\rho `$ is the matter density, and
$$W=\frac{3}{(kR)^3}\left[\mathrm{sin}(kR)(kR)\mathrm{cos}(kR)\right];$$
(22)
is a top-hat filter function. From the results of the nonlinear theory of gravitational collapse, stating that a spherical perturbation with overdensity $`\delta _c=\delta \rho /\rho >1.69`$ with respect to the background matter collapses to form a bound object, through the PS formalism we can derive the normalized fraction of collapsed objects per unit mass at a given redshift:
$$f(M,z)=\sqrt{\frac{2}{\pi }}\frac{\delta _c(1+z)}{\sigma _M^2}e^{\delta _c^2(1+z)^2/2\sigma _M^2}\left(\frac{d\sigma _M}{dM}\right).$$
(23)
Then the comoving number density of dark matter halos per unit mass is
$$n_h(M_h,z)=\frac{\mathrm{\Omega }_M\rho _{crit}}{M}f(M_h,z).$$
(24)
We can now ask what is the fraction of the IGM polluted by metals at different redshifts. To this end we calculate an IGM porosity parameter, $`Q(z)`$, defined by
$$dQ(z)=M_h|\frac{dn_h}{dz}|R^3(M_h,z)dz,$$
(25)
in two different cases. First, we consider the case in which metals are only injected in the IGM by superbubbles and therefore $`RR_e`$ in the previous equation. The resulting porosity evolution (complete overlap of the bubbles occurs for $`Q=0.16`$, Smith 1976) is shown in Fig. 3. Again we see that blowout by itself would lead only to a negligible dispersal of metals ($`Q<10^4`$), with most of the IGM maintaining its primordial composition to the present day.
If we are to explain the relatively ubiquitous presence of metals in the ‘true’ IGM, as deduced from at least some observations, we are then forced to assume that some additional physical mechanism, capable of efficiently transporting metals away from their production sites, must be at work. The nature of such a process can only be matter of speculation at present, because neither the numerical simulations nor the observations have yet reached the required levels of sophistication or sensitivity to address this question properly. Among the options which are worth exploring are galaxy collisions, diffusive processes and the peculiar motions of galaxies. We consider the relative importance of these different possibilities in future work.
Here we take a strictly phenomenological approach which nevertheless has considerable predicting power. We introduce a diffusive radius, $`R_d`$, defined as the mean interdistance between the galaxies responsible for providing the predominant contribution to the metal enrichment of the IGM (see Fig. 4), i.e. $`R_d=0.2`$ Mpc (comoving) for $`M_h>2\times 10^8M_{}`$. As the largest Doppler parameters measured in Ly$`\alpha `$ clouds are of order $`50`$ km s<sup>-1</sup>, the time required for a pollution front to travel such a distance is shorter than the Hubble time only at redshifts $`z<1`$. Thus, our assumption is equivalent to fixing at $`z\stackrel{<}{}1`$ the epoch at which metals become homogeneously distributed. As we will see shortly, this simple hypothesis leads to a number of consequences which are in accord with available data. If future studies of possible mechanisms for the diffusion of metals are able to determine directly the value of $`R_d`$ and its evolution with time, it will be relatively easy to include their results in the general framework of this paper and explore any differences with the conclusions presented here.
As can be seen by Fig. 3, when the condition $`R=max(R_e,R_d)`$ is introduced, mixing is improved dramatically—as expected, and metal-enriched bubbles indeed overlap at $`z1`$. At later epochs essentially all of the IGM is polluted with heavy elements produced in galaxies and subsequently redistributed by the combined effects of blowout and diffusion.
### 3.1 Metallicity of Polluted Regions
We can now calculate the average metallicity of the gas inside the diffusive spheres. In order to do so we need to know $`\mu _Z`$, the mass of metals produced by the typical supernova. Nucleosynthesis calculations in Type II SNae by Tsujimoto et al. (1995) predict $`\mu _Z23M_{}`$ for a Salpeter IMF, depending on the upper mass limit above which a black hole is formed. Here we adopt $`\mu _Z=2.58M_{}`$, also consistent with a matching to the cosmic star formation history (§4.2). By design, mixing is efficient and the metal distribution is therefore homogeneous within the radius $`R_d`$. The total number of SNae per galaxy is
$$𝒩_{SN}=\frac{\nu \mathrm{\Omega }_bf_b}{\tau _{}}M_h.$$
(26)
Then the mass of metals ejected by a galaxy inside a halo of mass $`M_h`$ is
$$M_e=\frac{\mu _Z\nu \mathrm{\Omega }_bf_b}{\tau _{}}M_h\delta _B,$$
(27)
with $`\delta _B`$ given by eq. 18. The metal density in a diffusive sphere surrounding a given halo is $`\rho _Z=3M_e/4\pi R_d^3`$, and its average metallicity is $`\rho _Z/\rho _i(z)`$, where $`\rho _i(z)`$ is the mean density of IGM at redshift $`z`$. This assumption restricts our analysis to the case of the ‘true’ IGM, characterized by overdensities close to unity. Extrapolation of the results to hydrogen column densities $`\mathrm{log}N_{HI}\stackrel{>}{}14`$ is only very qualitative and deserves more study. The dependence of metallicity in the diffusive spheres on halo mass and redshift is plotted in Fig. 4. At any given redshift, $`Z`$ increases as a function of mass (because of the increasing metal production) up to an abrupt cutoff when blowout becomes inhibited, as described by the behavior of $`\delta _B`$ in Fig. 1. As the spheres grow with time, the metals are distributed over larger volumes and this is reflected by the shift to lower metallicities with decreasing $`z`$ in Fig. 4. Note that for the same reason pockets of very high metal content ($`Z1Z_{}`$) are expected at large $`z`$, although the size of these regions is very small.
### 3.2 IGM Metallicity Evolution
Once the cosmological model has been fixed, we need to determine the values of the the star formation efficiency, $`\tau _{}`$, in order to calculate the metallicity of the IGM as a function of time. We choose this value by comparison with observations, as follows. We calculate the evolution of $`\mathrm{\Omega }_{}`$, the density parameter of stars formed in the universe (see Fig. 5), with the requirement that at $`z=0`$ it matches the estimate by Fukugita, Hogan & Peebles (1998). These authors conclude that $`\mathrm{\Omega }_{}(0)`$ in spheroids, disks and irregulars is $`0.0049h_{50}^2`$. A similar value is found by integrating current estimates of the star formation rate density in the universe as a function of redshift, as deduced from deep imaging surveys (Pettini 1999). With this normalization (giving $`\tau _{}^13\%`$), we can then derive the corresponding evolution of the metals produced by stars and returned to the ISM $`\mathrm{\Omega }_Z`$ (Fig. 5), which is directly proportional $`\mathrm{\Omega }_{}`$.
However, not all the metals can escape from the galaxy where they have been produced. The cosmic ejection fraction, $`f_{ej}(z)`$, (i.e. the ejection fraction averaged over the entire population of halos) is very close to unity at high redshift where predominantly small galaxies are present, but it steadily decreases to about 50% at $`z=0`$, as the number of more massive galaxies able to retain their metals increases. For this reason, the curve describing the density parameter of ejected metals $`\mathrm{\Omega }_Z^{ej}=f_{ej}(z)\mathrm{\Omega }_Z(z)`$ in Figure 5 progressively deviates from that for $`\mathrm{\Omega }_Z(z)`$ with decreasing redshift. Stated differently, today about 50% of the metals produced should reside in the IGM. If the metals were homogeneously mixed with the baryons in the universe at any redshift the average IGM metallicity (top curve in Fig. 5) would be $`Z=\mathrm{\Omega }_Z^{ej}/\mathrm{\Omega }_b1/25Z_{}`$ at $`z=3`$ and $`Z0.1Z_{}`$ at $`z=0`$. Note that the average metallicity of today’s galaxies would be higher by a factor $`\mathrm{\Omega }_b/(\mathrm{\Omega }_{}+\mathrm{\Omega }_g)`$ (the baryon density divided by the sum of star and gas density in galaxies); we therefore expect that for luminous matter today $`ZZ_{}`$, in accord with observational estimates (Edmunds & Phillips 1997).
Our main results are displayed in Figure 6 and 7, which show in the spread in metallicity as a function of redshift, compared with the average $`Z`$ of the IGM from Figure 5. The two figures correspond to different metallicities of the outflowing hot gas in the superbubbles, assumed to be $`Z_{}`$ in Fig. 6 and $`8Z_{}`$ in Fig 7. In the figures the metallicity distribution at each redshift considered is shown in the vertical direction with the density of symbols approximately proportional to the amplitude of the distribution. The distributions are relatively flat, but at redshifts $`z\stackrel{<}{}5`$ they also show a double-horned profile with points accumulating at the highest and lowest metallicities. We illustrate this effect in Figure 8, where we have reproduced the metallicity histograms at $`z=0`$ and $`z=1`$ from Figure 6. The more massive galaxies are responsible for the high metallicity peaks in these distributions, caused by the cutoff in the ejection fraction described by the parameter $`\delta _B`$ (Figure 1). The low metallicity maximum is due to the increasing number of low mass halos.
There are several interesting features of Figures 6 and 7 which we now discuss briefly. First, at $`z>6`$ the average IGM metallicity is not within the predicted ranges of metallicities, being lower than the lowest values in the distributions. What we are seeing is a highly inhomogeneous distribution of metals which are still clumped in small regions around the galaxies which produced them. As mixing proceeds, the average metallicity increasingly becomes a better description of the true mean metallicity, and approaches the mean of the distribution. At $`z1`$ the postulated diffusion leads to all the volume in the universe being polluted with heavy elements to some extent (Figure 3). The merging of the metal-enriched spheres produced by different galaxies has the effect of erasing regions of low metallicity: the minimum metallicity corresponds to that in spheres marginally overlapping at a given redshift. This effect is reflected in the growth of the lower boundary of the distributions from $`(1+z)\stackrel{<}{}5`$ in Figures 6 and 7, and results in the present-day IGM metallicity being confined within the narrow range $`0.1\pm 0.03Z_{}`$(1 $`\sigma `$).
The decreasing spread in metallicity with time is a direct result of the growth of $`Z`$. At high $`z`$, where the overall level of enrichment is low, even relatively metal-poor spheres have a chance to stand out. However, as the average IGM metallicity increases, only a few massive objects are able to produce diffusive spheres with $`Z`$ sufficiently high to be recognized as metallicity peaks (and still be able to eject their metals). The situation is similar to a “flooding effect”, where as the water level rises fewer and fewer mountain peaks can be seen.
At redshifts $`z>1`$, when metal enrichment is highly inhomogeneous, a considerable fraction of the IGM is still of essentially primordial composition. Thus, at $`z3`$ for example, we expect that only some low column density Ly$`\alpha `$ forest clouds will show associated metal lines, while the majority will not; the ratio between Ly$`\alpha `$ clouds with and without metals depends on the covering factor of the diffusive spheres and grows with time. This has implications for the interpretation of the results of searches for metals in the forest. Given the low optical depths when log $`N`$(H I) $`<14`$, such searches are normally conducted in a statistical way, by considering together the data from many absorption lines. In such cases one obtains some gross average over all the absorbers which, given the dilution with the unpolluted IGM, will be systematically lower than the values of $`Z`$ plotted in Figures 6 and 7. For this reason we also show in these figures the covering factor-weighted metallicity $`ZP(z)`$, where $`P(z)`$ describes the evolution of the covering factor down to $`z=1`$ .
In Figures 6 and 7 we have used different symbols to represent the contributions to the IGM metallicity from halos with virial temperatures above (hexagons) and below (triangles) $`10^4`$ K, to give a qualitative idea of the role of large and small objects in the enrichment process. It is seen that at high $`z`$ small objects are controlling the process, whereas at lower redshift enrichment is largely regulated by more massive galaxies (up to the ejection limit set by $`\delta _B`$, see Figure 3). If some inhibiting effect, such as photoheating by the UV background, affects preferentially low mass galaxies, the lower metallicity bounds at high redshift will move up accordingly, shifting closer to the line separating hexagons from triangles. The flat upper boundary of the distribution at high $`z`$ reflects the fact that the volumes involved are so small that the IGM baryon loading results in a negligible dilution of the metallicity of the ejecta. Since there is no firm measurement of this quantity at present (Heckman, private communication), we have considered two possibilities, $`Z_{}`$ (our standard case) and 8 $`Z_{}`$ in Figures 6 and 7 respectively. The IGM metallicity distribution is affected by this choice at high redshift, but the difference becomes much smaller at low $`z`$, where the IGM baryon loading regulates the dilution of the larger diffusive spheres.
Finally, we also show in the two figures the average metallicities for the case in which metal-enriched spheres with size below $`R_e(min)=1`$ kpc have been excluded. In our models these are the absolute lower limits to the average IGM metallicity at a given redshift. $`R_e(min)`$ has been calculated by imposing the condition that the collisional timescale between galaxies with a given impact parameter is shorter than the Hubble time at any redshift up to $`z=10`$. Obviously, a sphere with size greater than $`R_e(min)`$ can be strongly disturbed by tidal interactions following an encounter at greater impact parameter, and for this reason the curves obtained in this way are strictly lower limits. Nevertheless, as can be seen from Figure 6 and 7, imposing this condition reduces the metallicity by only about 40%.
In addition to the SCDM model, we have also considered a CDM model with a cosmological constant (CDM+$`\mathrm{\Lambda }`$) with $`\mathrm{\Omega }_M=0.4,\mathrm{\Omega }_\mathrm{\Lambda }=0.6,h=0.5`$). By normalizing the star formation rate following the same procedure as for the SCDM model we obtain a higher value of $`\tau _{}^1=10.5\%`$. Because of the normalization, the resulting metal distribution is qualitatively very similar to the one derived for the SCDM case. The only notable differences are a larger metallicity spread between $`0.5<z<2`$ and a somewhat higher mean value at $`z=0`$, $`Z=0.15\pm 0.03Z_{}`$.
## 4 Discussion
The basic conclusion of this paper is that metal ejection driven by SN events fails, by more than one order of magnitude, to distribute the products of stellar nucleosynthesis over volumes large enough to pollute the whole IGM to the typical metallicity of Ly$`\alpha `$ clouds, \[C/H\] $`2.5`$. We are therefore forced to conclude that some additional physical process must be at play, the nature of which remains to be determined. In our scheme transport of metals occurs in two sequential steps. First, SNae provide the initial kick necessary to eject heavy elements outside the potential well of a galaxy, but the ejecta are then confined by the IGM pressure to a relatively small bubble, of radius $`R_e`$. We then postulate that a second process is responsible for the diffusion of metals on a typical scale $`R_d`$, comparable to the mean distance between the galaxies which are the most efficient pollutants of the IGM.
Our scenario is qualitatively different from that proposed by Gnedin (1998), who attributed the mixing to more violent and rarer galaxy mergers. In the model proposed here SNae are of fundamental importance as they initiate the mixing process. However, galaxy collisions might well play a role in the subsequent phase during which metals are spread over larger scales, of order $`R_d`$. Other diffusive processes, such as thermal conduction and turbulent mixing layers at the interfaces between cosmological flows, as well as the peculiar motions of galaxies, may also be important in determining the spatial structure of the distribution of metals. A detailed study of such processes is beyond the scope of this paper and should, in any case, be based on cosmological simulations which are in progress. We will report of this extension of the work in a forthcoming paper.
In the present study we have stressed that the IGM enrichment proceeds in a very inhomogeneous manner, with pockets of metal rich gas gradually increasing both in number and in size until they overlap at $`z1`$. The average metallicity of the IGM increases with time, a trend confirmed by the results in Barlow & Tytler (1998), who found an order of magnitude increase in the metallicity of the Ly$`\alpha `$ forest between $`z=2.5`$ and $`z=0.5`$. Although our results are strictly applicable only to the ‘true’ IGM with overdensities close to unity, we nevertheless regard this as an encouraging performance of our model.
The metallicity spread is predicted to decrease with the progress of time. At $`z\stackrel{<}{}1`$, when the entire volume of the universe has been exposed to metal pollution, the spread in metallicity is less than one order of magnitude, and at the present epoch $`Z=0.1\pm 0.03Z_{}`$. Thus we predict that at $`z<1`$ essentially all absorbers should have associated C IV absorption, irrespectively of their column density. The effect should be very pronounced. Not only are the voids polluted by the overlap of metal-enriched diffusive spheres, as discussed above, but the decreasing intensity and hardness of the ionizing background lead to a further increase in the fraction of C which is triply ionized, so that the ratio $`N`$(C IV)/$`N`$(H I) increases for a fixed \[C/H\] (Rauch et al. 1997). It should be possible to test these predictions with forthcoming observations. UV-efficient echelle spectrographs, now available on the VLT and soon on the Gemini South telescope, will allow sensitive searches for C IV doublets at significantly lower redshifts than probed so far, while STIS on HST and FUSE will map the Ly$`\alpha `$ forest with the required spectral resolution at wavelengths below 3000 Å, which are inaccessible from the ground.
While our study has concentrated on the IGM, it also has important consequences for the metallicity of the intracluster medium. which is found to have a rather uniform value $`1/3Z_{}`$ with little, if any, evolution up to $`z0.3`$ (Fukazawa et al. 1998, Renzini 1999). This has been interpreted as evidence supporting the view that the enrichment process in clusters was already completed by that epoch. In general we expect two sources to contribute to the build up of the intracluster medium, infalling IGM and gas stripped from cluster galaxies. In our models these two components have significantly different composition. If the clusters formed at $`z1`$, we expect the IGM to have $`0.1Z_{}`$ (roughly constant from $`z1`$ to the present), while the galaxies’ ISM has approximately solar composition. The simplest estimate of the resulting metallicity is the geometric mean of these two values (as appropriate for hydrodynamical mixing problems, see Begelman & Fabian 1990), that is $`Z_{icm}=(0.1\times 1)^{1/2}Z_{}=0.32Z_{}`$, as observed. While this may well be a fortunate coincidence, it is also true that this conclusion lends qualitative support to the model and the assumptions made.
We make a final point concerning the role of intergalactic dust. The gas phase abundances we have derived do not take into account the possibility that some fraction of heavy elements may be locked up into dust grains. If this were the case, clearly the distributions of metallicities in Figure 6 and 7 would shift towards lower values of $`Z`$. In addition, at redshifts were the enrichment is still inhomogeneous, dust associated with regions of high metal concentration may reprocess UV/optical photons into IR radiation and give rise small scale anisotropies in the Cosmic Microwave Background which may be detectable (Ferrara et al. 1999). Having said this, however, we consider it highly speculative whether dust can survive at all in the hostile environments associated with outflowing superbubbles, where the temperatures are high and the gas has been shocked by SN explosions.
## 5 Summary
In this paper we have investigated the evolution of the metallicity of the intergalactic medium with particular emphasis on its spatial distribution. We have derived the conditions under which supernova-driven ejection of metals from galaxies can occur. A strong conclusion of our calculations is that if SNae were the only source of kinetic energy for the metals, a highly inhomogeneous distribution would result at any redshift. Under these conditions we would expect most of the volume of the universe to remain at near-primordial composition, with a metallicity $`Z10^4Z_{}`$, in contrast with the observational results discussed in the Introduction. Thus, an additional (but yet unknown) physical mechanism must be invoked to mix the metals on scales comparable to the mean distance between the galaxies which are the most efficient pollutants. From this simple hypothesis we have derived a number of testable predictions for the evolution of the IGM metallicity.
Quantitatively, we find that:
1. Metal ejection, or blowout, is inhibited in galaxies with total mass above $`M_h10^{12}(1+z)^{3/2}M_{}`$ due to the combined effects of their larger gravitational field and less coherent SN energy deposition.
2. The fraction of metals ejected over the star formation history of the universe is about 50% at $`z=0`$. We expect that at the present epoch approximately half of the metals are to be found in the IGM and the average metallicity of luminous matter to be approximately solar.
3. If the ejected metals were homogeneously mixed with the baryons in the universe, the average IGM metallicity would be $`Z=\mathrm{\Omega }_Z^{ej}/\mathrm{\Omega }_b1/25Z_{}`$ at $`z=3`$. However, due to the spatial inhomogeneity, $`Z`$ is actually lower than the mean of the distribution in the diffusive metal-enriched spheres.
4. Metals become homogeneously distributed in the IGM at $`z\stackrel{<}{}1`$, when the metal-enriched zones overlap, and the spread of the distribution is reduced. We calculate that at $`z=0`$ the IGM metallicity is in the range $`Z0.1\pm 0.03Z_{}`$.
5. The uniform metal abundance of intracluster gas from $`z0.3`$ to the present is naturally explained by a mixture of IGM infalling into the cluster and gas stripped from cluster members, with element abundances for both components as predicted by our models.
We should like to thank E. Corbelli, A. Meiksin and B. Nath for useful discussions. YS acknowledges support from the OAArcetri.
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# DARK MATTER CAUSTICS
## 1 Introduction
Before the onset of galaxy formation but after the time $`t_{eq}`$ of equality between matter and radiation, the velocity dispersion of the cold dark matter candidates is very small, of order $`\delta v_a(t)310^{17}\left(\frac{10^5eV}{m_a}\right)\left(\frac{t_0}{t}\right)^{2/3}`$ for axions and $`\delta v_W(t)10^{11}\left(\frac{GeV}{m_W}\right)^{1/2}\left(\frac{t_0}{t}\right)^{2/3}`$ for WIMPs, where $`t_0`$ is the present age of the universe and $`m_a`$ and $`m_W`$ are respectively the masses of the axion and the WIMP. In the context of galaxy formation, such small velocity dispersions are entirely negligible. Massive neutrinos, on the other hand, have primordial velocity dispersion $`\delta v_\nu (t)5.310^4\left(\frac{eV}{m_\nu }\right)\left(\frac{t_0}{t}\right)^{2/3}`$ which is comparable to the virial velocity in galaxies and therefore non-negligible in the context of galaxy formation $`^\mathrm{?}`$. This is the reason why massive neutrinos are called ‘hot dark matter’.
Collisionless dark matter particles lie on a thin 3-dimensional (3D) sheet in 6D phase-space. The thickness of this sheet is the primordial velocity dispersion $`\delta v`$. If each of the aforementioned species of collisionless particles is present, the phase-space sheet has three layers, a very thin layer of axions, a medium layer of WIMPs and a thick layer of neutrinos. The phase-space sheet is located on the 3D hypersurface of points $`(\stackrel{}{r},\stackrel{}{v}):\stackrel{}{v}=H(t)\stackrel{}{r}+\mathrm{\Delta }\stackrel{}{v}(\stackrel{}{r},t)`$ where $`H(t)=\frac{2}{3t}`$ is the Hubble expansion rate and $`\mathrm{\Delta }\stackrel{}{v}(\stackrel{}{r},t)`$ is the peculiar velocity field. Fig. 1 shows a 2D section of
6D phase-space along the $`(z,\dot{z})`$ plane. The wiggly line is the intersection of the 3D sheet on which the particles lie in phase-space with the plane of the figure. The thickness of the line is the velocity dispersion $`\delta v`$, whereas the amplitude of the wiggles in the line is the peculiar velocity $`\mathrm{\Delta }v`$. If there were no peculiar velocities, the line would be straight since $`\dot{z}=H(t)z`$ in that case.
The peculiar velocities are associated with density perturbations and grow by gravitational instability as $`\mathrm{\Delta }vt^{2/3}`$. On the other hand the primordial velocity dispersion decreases on average as $`\delta vt^{2/3}`$, consistently with Liouville’s theorem. When an overdensity enters the non-linear regime, the particles in its vicinity fall back onto it. This implies that the phase-space sheet ‘winds up’ there in clockwise fashion. One such overdensity is shown in Fig. 1. In the linear regime, there is only one value of velocity, i.e. one single flow, at a typical location in physical space, because the phase-space sheet covers physical space only once. On the other hand, inside an overdensity in the non-linear regime, the phase-space sheet covers physical space multiple times implying that there are several (but always an odd number of) flows at such locations.
At the boundary surface between two regions one of which has $`n`$ flows and the other $`n+2`$ flows, the physical space density is very large because the phase-space sheet has a fold there. At the fold, the phase-space sheet is tangent to velocity space and hence, in the limit of zero velocity dispersion $`(\delta v=0)`$, the density diverges since it is the integral of the phase-space density over velocity space. The structure associated with such a phase-space fold is called a ’caustic’. It is a surface in physical space. It is easy to show that, in the limit of zero velocity dispersion, the density diverges as $`d\frac{1}{\sqrt{\sigma }}`$ when the caustic is approached from the side with $`n+2`$ flows, where $`\sigma `$ is the distance to the caustic. Velocity dispersion cuts off the divergence.
As mentioned above, the process of galactic halo formation involves the local winding up of the phase-space sheet of collisionless dark matter particles. If the galactic center is approached from an arbitrary direction at a given time, the local number of flows increases. First, there is one flow, then three flows, then five, seven… The number of flows at our location in the Milky Way galaxy today has been estimated $`^\mathrm{?}`$ to be of order 100. The boundary between the region with one (three, five, …) and the region with three (five, seven, …) flows is the location of a caustic which is topologically a sphere surrounding the galaxy. When these caustic spheres are approached from the inside the density diverges as $`d\frac{1}{\sqrt{\sigma }}`$ in the zero velocity dispersion limit. These spheres are the outer caustics in the phase-space structure of galactic halos. In addition there are inner caustics.
It is a little more difficult to see why there must be inner caustics, and to derive their structure. See refs. for details. The inner caustics are rings. They are located near where the particles with the most angular momentum in a given in and out flow are at their distance of closest approach to the galactic center. A ring is a closed tube whose cross-section is a $`D_4`$ catastophe $`^\mathrm{?}`$. The cross-section is shown in Fig. 2 in the limit of axial and reflection symmetry, and where the transverse dimensions, $`p`$ and $`q`$, are much smaller than the ring radius $`a`$.
In the absence of any symmetry, the cross-section of the tube does not have the exact shape shown in Fig. 2 but it still has that shape qualitatively, i.e. it is still a closed line with three cusps one of which points away from the galactic center.
The existence of caustic rings of dark matter follows from only two assumptions:
1. the existence of collisionless dark matter
2. that the velocity dispersion of the infalling dark matter is much less, by a factor ten say, than the rotation velocity of the galaxy.
Only the second assumption requires elaboration. Velocity dispersion has the effect of smoothing out caustics. The question is when is the velocity dispersion so large as to smooth caustic rings over distance scales of order the ring radius $`a`$, thus making the notion of caustic ring meaningless. In ref. this critical velocity dispersion was estimated to be 30 km/s = $`10^4`$ for the caustic rings in our own galaxy, whose rotation velocity is 220 km/s. $`10^4`$ is much less than the primordial velocity dispersion $`\delta v`$ of the cold dark matter candidates. However the velocity dispersion $`\mathrm{\Delta }v`$ associated with density perturbations also smoothes caustics in coarse grained observations. So the question is whether the velocity dispersion $`\mathrm{\Delta }v`$ of cold dark matter particles associated with density perturbations falling onto our galaxy is less than 30 km/s. The answer is yes with high probability since the infalling dark matter particles are not associated with any observed inhomogeneities. 30 km/s is of order the velocity dispersion of the Magellanic Clouds. For the velocity dispersion of the dark matter particles presently falling onto our galaxy to be as large 30 km/s, these particles would have to be part of clumps whose mass/size ratio is of order that of the Magellanic Clouds. But if that were the case, why did these clumps fail to become luminous?
One might ask whether caustic rings can be seen in N-body simulations of galaxy formation. The generic surface caustics associated with simple folds of the phase-space sheet have been seen $`^\mathrm{?}`$. However, caustic rings would require far greater resolution than presently available, at least in a 3D simulation of our own halo. Indeed, the largest ring in our galaxy has been estimated $`^\mathrm{?}`$ to have radius of order 40 kpc. It is part of an in and out flow that extends to the Galaxy’s current turnaround radius, of order 2 Mpc. To resolve this first ring, the spatial resolution would have to be considerably smaller than 10 kpc. Hence a minimum of $`2\frac{1}{(10\mathrm{k}\mathrm{p}\mathrm{c})^3}\frac{4\pi }{3}(2\mathrm{M}\mathrm{p}\mathrm{c})^3710^7`$ particles would be required to see the caustic ring in a simulation of this one flow. However, the number of flows at 40 kpc in our halo $`^\mathrm{?}`$ is of order 10. So it appears that $`10^9`$ particles is a strict minimum in a 3D simulation of our halo. Even so, this addresses only the kinematic requirement of resolving the halo in phase-space, assuming moreover that the particles are approximately uniformly distributed on the phase-space sheets. There is a further dynamical requirement that 2-body collisions do not artificially ’fuzz up’ the phase-space sheets. Indeed 2-body collisions are entirely negligible in the flow of cold dark matter particles such as axions or WIMPs. On the other hand, 2-body collisions are present, and hence the velocity dispersion is artificially increased, in the simulations. This may occur to such an extent that the caustics are washed away even if $`10^9`$ particles are used.
In the self-similar infall model $`^{\mathrm{?},\mathrm{?}}`$ of galactic halo formation the caustic ring radii $`a_n`$ are predicted $`^\mathrm{?}`$:
$$\{a_n:n=1,2,\mathrm{}\}(39,19.5,13,10,8,\mathrm{})\mathrm{kpc}\left(\frac{j_{\mathrm{max}}}{0.25}\right)\left(\frac{0.7}{h}\right)\left(\frac{v_{rot}}{220\frac{km}{s}}\right)$$
(1)
where $`h`$ is the present Hubble rate in units of 100 km/s.Mpc, $`v_{rot}`$ is the rotation velocity of the galaxy and $`j_{\mathrm{max}}`$ is the maximum of its dimensionless angular momentum distribution as defined in ref. . In Eq. (1) we assume that the parameter $`^{\mathrm{?},\mathrm{?}}`$ $`ϵ=0.3`$.
Eq. (1) predicts the caustic ring radii of a galaxy in terms of its first ring radius $`a_1`$. If the caustic rings lie close to the galactic plane they cause bumps in the rotation curve, at the caustic ring radii. As a possible example of this effect, consider $`^\mathrm{?}`$ the rotation curve of NGC3198, one of the best measured. It has three faint bumps at radii: 28, 13.5 and 9 kpc, assuming $`h=0.75`$. The ratios happen to be consistent with Eq. (1) assuming the bumps are caused by the first three ($`n=1,2,3`$) ring caustics of NGC3198. Moreover, since $`v_{\mathrm{rot}}=150`$ km/s, $`j_{\mathrm{max}}`$ is determined to equal 0.28. The uncertainty in $`h`$ is a systematic effect that can be corrected for when determining $`j_{\mathrm{max}}`$ because the bump radii scale like $`1/h^{}`$ where $`h^{}`$ is the Hubble rate assumed by the observer in constructing the rotation curve, and the caustic ring radii scale as $`1/h`$. Rises in the inner rotation curve of the Milky Way were also interpreted $`^\mathrm{?}`$ as due to caustics $`n=6,7,8,9,10,11,12`$ and 13. This determined the value of $`j_{\mathrm{max}}`$ of our own galaxy to be 0.263. The first five caustic ring radii in our galaxy are then predicted to be: 41, 20, 13.3, 10, 8 kpc.
## 2 Evidence for universal structure in galactic halos
Motivated by these findings, we analyzed $`^\mathrm{?}`$ a set of 32 extended well-measured galactic rotation curves which had been previously selected $`^\mathrm{?}`$ under the criteria that each is an accurate tracer of the galactic radial force law, and that it extends far beyond the edge of the luminous disk.
According to the self-similar caustic ring model, each galaxy has its own value of $`j_{\mathrm{max}}`$. Over the set of 32 galaxies selected in ref. , $`j_{\mathrm{max}}`$ has some unknown distribution. However, the fact that the values of $`j_{\mathrm{max}}`$ of NGC3198 and of the Milky Way happen to be close to one another, within 7%, suggests that the $`j_{\mathrm{max}}`$ distribution may be peaked near a value of 0.27 . Our strategy is to rescale each rotation curve according to
$$r\stackrel{~}{r}=r(\frac{220\mathrm{km}/\mathrm{s}}{v_{rot}})$$
(2)
and to add them in some way. Since Eq. (1) predicts the $`n`$th caustic radius $`a_n`$ to be distributed like $`j_{\mathrm{max}}`$ for all n, and it fixes the ratios $`a_n/a_11/n`$, the sum of rotation curves should show the $`j_{\mathrm{max}}`$ distribution, once for $`n=1`$, then at about half the $`n=1`$ radii for $`n=2`$, then at about 1/3 the $`n=1`$ radii for $`n=3`$,and so on. If the $`j_{\mathrm{max}}`$ distribution is broad, the sum of rotation curves is unlikely to show any feature. However, if it is peaked, then the sum should show a peak for $`n=1`$ at some radius, then again at 1/2 that radius for $`n=2`$, at 1/3 the radius for $`n=3`$, and so on. If the $`j_{\mathrm{max}}`$ distribution is peaked at 0.263 (the value for the Milky Way) the peaks in the sum of rotation curves should appear at 41 kpc, 20 kpc, 13.3 kpc …
The procedure followed to add the 32 rotation curves is described in detail in ref. . Briefly, we proceeded as follows. For each rotation curve, all data points with rescaled radii $`\stackrel{~}{r}<10\mathrm{kpc}`$ were deleted to remove the effect of the luminous disk. The remaining points were then fitted to a line. The rotation velocity $`v_{rot}`$ used to rescale the radii in Eq. (2) is the average of that line. The rms deviation $`\sqrt{\delta v{}_{}{}^{2}}`$ from the linear fit was determined for each galaxy. This was taken to be the error on the residuals $`\delta v_i`$, i.e. the differences between the measured velocities in a rotation curve and the linear fit. Finally the sample of 32 galaxies was averaged in $`2\mathrm{kpc}`$ radial bins:
$$b_i\frac{1}{N_i}\underset{j=1}{\overset{N_i}{}}\delta \stackrel{~}{v}_j,$$
(3)
where $`N_i`$ is the number of data points in the bin. The assigned error on each $`b_i`$ is then simply $`1/\sqrt{N_i}`$. Fig. 3 shows the result.
There are two features evident at roughly $`20`$ and $`40\mathrm{kpc}`$. A fit to two Gaussians plus a constant indicates features at $`19.4\pm 0.7\mathrm{kpc}`$ and $`41.3\pm 0.8\mathrm{kpc}`$, with overall significance of $`2.4\sigma `$ and $`2.6\sigma `$, respectively. Fig. 3 shows the fitted curve. When the same fit is applied to the same data in 1 kpc bins, the significance of the two peaks is 2.6 and 3.0 $`\sigma `$ respectively. The locations of the features agrees with the predictions of the self-similar caustic ring model with the $`j_{\mathrm{max}}`$ distribution peaked at 0.27. The use of Gaussians to fit the peaks in the combined rotation curve was an arbitrary choice in the absence of information on the $`j_{\mathrm{max}}`$ distribution.
The existence of velocity peaks and caustic rings in the cold dark matter distribution is relevant to axion $`^\mathrm{?}`$ and WIMP searches $`^\mathrm{?}`$. Caustics may also be investigated using gravitational lensing techniques $`^\mathrm{?}`$.
## Acknowledgements:
This work was supported in part by the US Department of Energy under grant No. DEFG05-86ER40272.
## References
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# CPHT S 040.0400 Reinterpretation of Effective Chiral Lagrangian
(April 2000)
Effective Tree Chiral Lagrangian is interpreted as a power series expansion of the kinematical variables. In the presence of the strong interaction this expansion is valid below the unitarity cut, hence in the unphysical region. Consequences of this reinterpretation of the Chiral Lagrangian are analysed for the relation between $`K\pi `$ and $`K2\pi `$ transitions.
There has been recent interest how to handle the problem of the strong final state pion pion interaction in $`K_s2\pi `$ decay. In particular how to determine the off-shell $`K\pi `$ transition using the input of the measured $`K2\pi `$ rate. This last reaction is quite difficult to calculate by the technique of the lattice gauge theory because of the strong $`\pi \pi `$ final state interaction and of other reasons. The $`K\pi `$ transition should be easier to calculate using the lattice technique and could provide the answer to the origine of the $`\mathrm{\Delta }I=1/2`$ problem .
In a related problem which is now of central interest, is to understand the experimental ratio $`ϵ^{}/ϵ`$ of the CP violation problem. Eventually, using the lattice theory, one should be able to calculate numerically this CP violation effect by reducing them to the $`K\pi `$ problem. We note that there is a recent calculation of this ratio taking into account of the strong final state interaction using the technique discussed previously .
A possible solution for the first problem was given a long time ago . It was based on a reinterpretation of the Effective Chiral Lagrangian (CL) at the tree level which has not been clear to all readers. This has led to some questions raised by a number of authors and more recently by Buras et al.. We wish to clarify in this note some points raised by these authors. It is shown here that interpreting the tree CL of the process $`K2\pi `$ amplitude as a power series expansion in momentum together with requiring the zero of the matrix element as demanded by the Cabibbo and Gell-Mann theorem , assure that the final result is model independent. Our line of approach has been used in a series of articles .
In this note the problem of the final state interaction is reexamined. The main idea is that, just the same as in our previous work, the tree CL is an explicit manifestation of the current algebra soft pion theorems. It was first invented to avoid some complicated and cumbersome manipulations of the current algebra technique . It is crucial to note that the current algebra soft pion theorems are still valid in the presence of the hadronic initial or final state interactions. The current algebra relation between the form factors of $`K\pi \pi e\nu `$ and $`K\pi e\nu `$ is such an example . It is still valid in the presence of the strong interaction between $`K\pi `$ and $`\pi \pi `$
As was previously suggested, the tree CL, should bear this important property i.e it should be valid even in the presence of the strong interactions among the hadrons involved. At first sight this cannot be done because CL only gives a power series expansion in the invariant variables $`s,t,u\mathrm{}`$ of the matrix element which must be real in the physical region while the strong final state interaction should make them complex. There is, however, a region where the power series expansion of the matrix element is valid, namely outside the cut in the unphysical region. We shall make use of the CL to give relations among different processes just the same as the current algebra technique. The tree CL can therefore be considered as low energy theorems in the unphysical region for our purpose. It is not difficult to analytically continue these low energy theorems to the physical region by using the technique of the integral equations of the Omnes-Muskhelishvili type (OM) , or the inverse amplitude method etc.
The effective Lagrangian for the $`\mathrm{\Delta }I=1/2`$ $`K2\pi `$ is given by :
$$M(K_S(k)\pi ^+(p)+\pi ^{}(q))=\frac{1}{\sqrt{2}}Cf_\pi (2k^2p^2q^2)$$
(1)
and
$$M(K_L\pi ^0)=iC\sqrt{2}f_\pi ^2q(\pi ).k(K)$$
(2)
where $`f_\pi =93MeV`$ and is the pion decay constant. Notice the constant C are common to both equations which the manifestation of the CA soft pion theorem relating the off-shell $`K\pi `$ to $`K2\pi `$ amplitude.
Let us consider the Eq. (1) when both pions are on their mass shell and the Kaon off its mass shell. The usual interpretation of this equation is simply the result of the effective tree Lagrangian or the result with the strong final state $`\pi \pi `$ interaction switched off. To take into account of the final state $`\pi \pi `$ interaction, Chiral Perturbation Theory (ChPT) could be used. Because of the presence of the undetermined counterterms and of the violation of the unitarity inherited in the perturbative schemes, the phase theorem is no longer satisfied in this scheme and hence this approach is not useful for our purpose.
The crucial point is to reinterprete Eq. (1) as the result of the first two terms of a power series expansion in $`s=k^2`$ variable of the matrix elements $`K2\pi `$ with the presence $`\pi \pi `$, $`4\pi `$, $`K\pi `$ … interactions. This expansion is only valid in the unphysical region. Let us denote this matrix element with the two pions on their mass shell as $`A(s)`$. It is assumed that $`A(s)`$ is analytic in the cut plane with a cut from $`4m_\pi ^2`$ to $`\mathrm{}`$. In reality, $`A(s)`$ is a product of two functions, the self energy operators of the Kaon and the $`K2\pi `$ vertex. (We shall neglect in the following the Kaon self energy operators due to its higher threshod $`K2\pi `$…).
Let us first discuss the solution of this problem from a more general viewpoint. In our non-perturbative approach, the effective tree CL represents low energy theorems with strong (final state) interaction taken into account. Below the cut their contribution can be represented by a polynomial in $`s`$ of degree $`n`$ and with real coefficients. Without loss of generality this polynomial can also be rewritten as as a polynomial in $`(ss_0)`$ variable where $`s_0`$ is in the unphysical region which will be taken on the real $`s`$ axis below the branch point $`4m_\pi ^2`$.
Assuming that $`A(s)`$ is polynomially bounded and that $`A(s)s^{(n+1)}0`$ as $`s\mathrm{}`$, $`n>0`$, we can write the following dispersion relation
$$A(s,s_0)=a_0+a_1(ss_0)+\mathrm{}a_n(ss_0)^n+\frac{(ss_0)^{n+1}}{\pi }_{4m_\pi ^2}^{\mathrm{}}\frac{ImA(z)dz}{(zs_0)^{n+1}(zsiϵ)}$$
(3)
Around $`s=s_0`$, the dispersion integral is of the order $`(ss_0)^{n+1}`$ and can be neglected, hence the low energy theorem is recovered. Needless to say, $`a_n`$ are, apart from a factorial factor $`n!`$, the derivatives of $`A(s)`$ evaluated at $`s=s_0`$.
The mathematical problem is now clear: Find the solution of the integral equation of the OM type for $`A(s,s_0)`$ with its imaginary part given by the elastic unitarity:
$$ImA(s)=A(s,s_0)e^{i\delta (s)}sin\delta (s)$$
(4)
with the boundary conditions around $`s_0`$ given by Eq. (3) and where $`\delta `$ is the S-wave $`I=0`$ $`\pi \pi `$ phase shifts.
To solve this integral equation Eq. (3) which is of OM type , let us define the function $`\mathrm{\Omega }(s,s_0)`$ normalized to unity for convenience at $`s=s_0`$:
$$\mathrm{\Omega }(s,s_0)=\mathrm{exp}(\frac{ss_0}{\pi }_{4m_\pi ^2}^{\mathrm{}}\frac{\delta (z)dz}{(zs_0)(zsiϵ)})$$
(5)
The solution for our integral equation is:
$$A(s,s_0)=P_n(s)\mathrm{\Omega }(s,s_0)$$
(6)
where $`P_n(s)`$ is a polynomial in $`s`$ of order n with real coefficients. They can be determined by expanding the function $`\mathrm{\Omega }(s,s_0)`$ in a power series in $`(ss_0)`$, and compare Eq. (6) with Eq. (3). The expansion in Taylor’s series is possible because $`\mathrm{\Omega }(s,s_0)`$ is an analytic function with a cut from $`4m_\pi ^2`$ to $`\mathrm{}`$.
In the special case where only two terms in the series are known such as the case of the CL given by Eq. (1), the solution of our integral equation is given by:
$$A(s)=\{a_0+(ss_0)(a_1a_0\mathrm{\Omega }^{^{}}(s_0,s_0))\}\mathrm{\Omega }(s,s_0)$$
(7)
where $`\mathrm{\Omega }^{^{}}`$ denotes the first derivative of $`\mathrm{\Omega }(s,s_0)`$ evaluated at $`s_0`$ and is given by:
$$\mathrm{\Omega }{}_{}{}^{}(s_0,s_0)=\frac{1}{\pi }_{4m_\pi ^2}^{\mathrm{}}\frac{\delta (z)dz}{(zs_0)^2}$$
(8)
The presence of the term $`\mathrm{\Omega }^{^{}}(s_0,s_0)`$ is to ensure the boundary condition for $`A(s,s_0)`$ is satisfied.
It is straightforward to generalise the solution of Eq. (3) for other values of $`n`$. For example when $`n=2`$, the solution for the integral equation Eq.( 3) is obtained by adding to the curly bracket on the righthand side of Eq. (7) a term:
$$(ss_0)^2(a_2a_1\mathrm{\Omega }^{^{}}(s_0,s_0)a_0\frac{\mathrm{\Omega }^{^{\prime \prime }}(s_0,s_0)}{2}+a_0\mathrm{\Omega }^{}_{}{}^{}2(s_0,s_0))$$
(9)
One is tempted to write a simpler solution than that given by Eq. (7) by construcing for example:
$$A(s,s_0)=(a_0+a_1(ss_0))\mathrm{\Omega }_2(s,s_0)$$
(10)
where
$$\mathrm{\Omega }_2(s,s_0)=\mathrm{exp}(\frac{(ss_0)^2}{\pi }_{4m_\pi ^2}^{\mathrm{}}\frac{\delta (z)dz}{(zs_0)^2(zsiϵ)})$$
(11)
which satisfies the boundary conditions, but violates the condition on the polynomially boundedness. This is so because by partial fraction, one can show $`\mathrm{\Omega }_2(s,s_0)=\mathrm{\Omega }(s,s_0)/\mathrm{exp}((ss_0)\mathrm{\Omega }{}_{}{}^{}(s_0,s_0))`$ which has an exponential behavior. This result is totally expected because the dispersion relation for $`\mathrm{log}\mathrm{\Omega }(s,s_0)`$ obeys at most a once subtracted dispersion relation due to the polynomial boundedness.
As usual the solution given by Eq. (7) has the polynomial ambiguity because we can multiply the RHS of Eq. (7) by a polynomial factor $`1+_{n=2}^Nc_n(ss_0)^n`$ with $`N2`$. We assume, in the following, there is no such ambiguity, or the zeros introduced by such an ambiguity is sufficiently far away from the physical region of interest.
For the $`K\pi \pi `$ problem, invoking the Cabibbo and Gell-Mann theorem which requires the matrix element to vanish in the SU(3) limit, one has $`s_0=m_\pi ^2`$.
We now return to our problem. Let us rewrite Eq. (1) with the two pion on their mass shell as:
$$M(K_S(k)\pi ^+(p)+\pi ^{}(q))=\sqrt{2}Cf_\pi (sm_\pi ^2)+\mathrm{}$$
(12)
This equation should be considered as the power series around $`s=m_\pi ^2`$ keeping only the first 2 terms, namely the matrix element vanishes at $`s=m_\pi ^2`$ and its derivative at this point is known. The solution for the corresponding Omnes-Muskhelishvilli equation is therefore:
$$M(s)=\sqrt{2}Cf_\pi (sm_\pi ^2)\mathrm{\Omega }(s,m_\pi ^2)$$
(13)
as can be seen using the result of Eq. (7) and letting $`a_0=0`$ and $`s_0=m_\pi ^2`$ as required by Eq. (12). The condition on the position of the zero of the matrix element at $`s=m_\pi ^2`$ is a direct consequence of the Cabibbo-Gell-Mann theorem on the SU(3) symmetry of the problem . Eq. (13) were derived earlier without giving explicitly a justification .
The physical value of the matrix element is obtained by setting $`s=m_K^2`$.
Using the experimental rate for $`K_s2\pi `$ and the S-wave, $`I=0`$ phase shifts as given by the unitarized one loop ChPT which fit to the experimental data , one obtains $`C=\mathrm{0.90.10}^{11}MeV^2`$. This is the result of the reference .
Our strategy to study the $`K_s2\pi `$ and $`K_L3\pi `$ is therefore to calculate first the $`K\pi `$ transition by lattice gauge theory or by some approximate schemes and compare them with the value $`C`$ given above.
For other problems, such as the $`\eta \pi ^+\pi ^{}\pi ^0`$, because there is no Cabibbo and Gell-Mann theorem when the $`\pi ^0`$ and the $`\pi ^{}`$ are soft, how can $`s_0`$ be determined ? An approximate answer to this problem is to realise that the effective CL is a realisation of the Current Algebra soft pion theorems . By taking $`\pi ^0`$ soft, using current algebra, this problem is reduced to the matrix element $`<\eta v\pi ^+(q_1)\pi ^{}(q_2)>`$ where $`v`$ is a pseudo-scalar operator. This matrix element is similar to the $`K2\pi `$ problem treated here. Taking one of the remaining two pions soft, the other pion on its mass shell, this matrix element is related to the $`\eta \pi `$ mixing problem. Going through this process for the $`\eta 3\pi `$ problem, one has $`s=(q_1+q_2)^2m_\pi ^2=s_0`$ when $`q_10`$. This point sets the scale to this problem in terms of the $`\eta \pi `$ mixing . The corresponding dispersion relation for this problem is similar to Eq. (3), except only a once subtraction at $`s_0`$ is needed, with $`a_0`$ given by the value of the $`\eta \pi `$ mixing.
This value of $`s_0`$ is different from that given in the reference , namely $`s_0=0`$, the chiral $`SU(3)`$ limit. The final results is insensitive to these two choices of $`s_0`$. They differ from each other by only a few percents.
We have presented here a reasonable method to determine $`s_0`$. This method is inspired by the current algebra technique .
This discussion of the $`\eta \pi ^+\pi ^{}\pi ^0`$ problem is also of interest for calculating the matrix element of $`(8_L,8_R)`$ operators of the $`\mathrm{\Delta }I=3/2`$ $`K2\pi `$ problem.
The following discussion on the subtraction point $`s_0`$, based on the theory of analytic function is of some interest. Because $`A(s,s_0)`$ is an analytic function in the cut s-plane, its Taylor’s series converges inside a circle with a center at $`s_0`$ and of a radius $`4m_\pi ^2s_0`$. The choice of the number of terms in the series to achieve a given accuracy depends on the physical situation and also on the choice of the point $`s_0`$.
For example for $`s_0`$ near to the branch point $`s=4m_\pi ^2`$, the radius of the convergence of the series is very small and hence more terms are needed in the power series if the series converges at all. For $`s_0=4m_\pi ^2`$, the radius of the convergence is zero which is expected because all the derivatives of $`A(s)`$ evaluated at this point become infinite due to the square root threshold singularity and hence the series diverges.
If the expansion point $`s_0`$ was taken far away from the origin and on the negative s axis, the radius of the convergence of the series would be, in principle, larger but we would have to take more terms in the series in order to explore the boundary conditions near the origin which is the chiral limit of the matrix element.
As it is shown above, some physical input must be made to restrict the determination of $`s_0`$.
To see the sensitivity of our solution on the expansion point $`s_0`$ we pretend to ignore the Cabibbo-Gell Mann theorem and study the solution of Eq.(7) as a function of $`s_0`$. Let us take $`s_0=0`$ and $`s_0=2m_\pi ^2`$. The former yields the zero of the $`K2\pi `$ amplitude at $`s0.96m_\pi ^2`$ which is prettty near to the Cabbibo and Gell-Mann point, and the latter at $`s0.92m_\pi ^2`$ which is a larger violation of this theorem. The only point where there is no violation of the Cabibbo and Gell-Mann theorem is $`s_0=m_\pi ^2`$ which is totally expected.
Normalising the factor $`\sqrt{2}Cf_\pi `$ to be unity, for $`s_0=0,m_\pi ^2`$, and $`2m_\pi ^2`$ we have, respectively, the absolute value of the physical matrix element ($`s=m_K^2`$) to be 19, 18.2 and 15.4 which shows some sensitivity on the choice of $`s_0`$. Fortunately for our problem, $`s_0=m_\pi ^2`$ as required by the Cabibbo and Gell-Mann theorem.
As $`s_0`$ approaches the branch point, the violation of this theorem becomes larger. For example at $`s_0=3m_\pi ^2`$, the zero of the matrix element is at $`0.5m_\pi ^2`$ which is a large violation. This is due to the square root threshold singularity of the matrix element, resulted by the threshold behavior of the S-wave $`\pi \pi `$ phase shift, $`\delta 0`$ as $`\sqrt{s4m_\pi ^2}`$.
In Fig. (1) the function $`\mathrm{\Omega }(s,0)`$ derived in is plotted as a function of s. For $`s<4m_\pi ^2`$ this function is real. For $`s>4m_\pi ^2`$, $`\mathrm{\Omega }(s,0)`$ is complex, only its real part is plotted. As one can see at $`s=4m_\pi ^2`$ there is a cusp associated with what was known as the square root (threshold) singularity, but is now misnamed as the ”chiral logarithm” singularity . This singularity is due to the threshold behaviour of the phase shift as discussed above. On the same graph, the imaginary part of $`\mathrm{\Omega }(s,0)`$ is plotted as a function of s.
As can be seen in Fig. (1), because of the square root singularity, a power series expansion for $`\mathrm{\Omega }(s)`$ near to the branch point requires many terms to give the correct energy dependence; the series converges inside a small circle with the center at $`s_0`$ and with a radius $`4m_\pi ^2s_0`$.
Our approach to this problem is heavily based on the reinterpretation of the tree CL as a power series expansion below the $`2\pi `$ threshold even in the presence of the strong $`\pi \pi `$ (final state) interaction and on the use of the current algebra low energy theorems. It is quite different from the spirit of the reference where ChPT is used and hence the assumption on the derivative on $`s_0=m_\pi ^2`$ has to be made. Our reinterpretation of the tree CL leads naturally to this condition. Some of the points discussed in their article are clarified in this article.
This author would like to thank Luis Oliver for pointing out the existence of the reference .
Figure Captions
Fig. 1 :The real part of the function $`\mathrm{\Omega }(s,0)`$ as a function of $`s`$ (in the unit $`m_\pi ^2=1`$) is shown by the solid line; the imaginary part of $`\mathrm{\Omega }(s,0)`$ is shown by the dotted line.
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# Untitled Document
Ref. SISSA 42/2000/FM
$`A(SL_q(2))`$ at roots of unity is a free module over $`A(SL(2))`$
Ludwik Da̧browski, Cesare Reina and Alessandro Zampa
SISSA, via Beirut 2-4, 34014 Trieste (Italy)
dabrow@sissa.it, reina@sissa.it, zampa@fm.sissa.it
Abstract. It is shown that when $`q`$ is a primitive root of unity of order not equal to $`2`$ mod $`4`$, $`A(SL_q(2))`$ is a free module of finite rank over the coordinate ring of the classical group $`SL(2)`$. An explicit set of generators is provided.
MSC: 17B37, 16W35, 81R50
Key words: quantum group
Quantum groups at roots of unity have particularly rich and interesting structure. In this letter we adopt the ‘quantum functions’ point of view, dual to the quantum universal enveloping algebras (cf. ). It is known that the quantum Hopf algebra $`A(SL_q(2))`$ at primitive roots of unity $`q^l=1,l`$ odd, or $`q^{2l}=1,l`$ even, can be viewed as a module over the algebra $`A(SL(2))`$, eg. by using the Frobenius mappings . This module is projective and finitely generated and thus it corresponds to a locally free coherent sheaf $``$ over the affine group scheme $`SL(2)=:SL(2,C)`$. It has been also established that $`SL_q(2)`$ at primitive odd roots of unity forms a principal fibre bundle (faithfully flat Hopf-Galois extension) over $`SL(2)`$ . (See also , , for the case of cubic root of unity, that was conjectured in to be important in understanding the (quantum) symmetry of the Standard Model of fundamental interactions).
It may have been guessed for a long time that $`A(SL_q(2))`$ at roots of unity is actually free over $`A(SL(2))`$ but, as far as we know, no proof has been given. Here we show that it is actually so. Our reasoning underlines the difficulties to tackle the more general problem of other “simple” quantum groups at roots of unity.
Recall that for $`Cq0`$, $`A(SL_q(2))`$ is the free algebra $`C<a,b,c,d>`$ modulo the ideal generated by the commutation relations
$$ab=qbaac=qca,bd=qdb,bc=cb,cd=qdc,adda=(qq^1)bc$$
and by the $`q`$-determinant relation $`adqbc=1.`$ When $`q`$ is a primitive $`l`$-th root of unity, for odd $`l`$, the Hopf subalgebra generated by
$$\alpha =a^l,\beta =b^l,\gamma =c^l,\delta =d^l$$
is isomorphic to the Hopf algebra
$$A(SL(2))=C[\alpha ,\beta ,\gamma ,\delta ]/<\alpha \delta \beta \gamma 1>$$
with the restricted coproduct, counit and coinverse. Moreover, $`A(SL_q(2))`$ is a finitely generated module over $`A(SL(2))`$. The same happens when $`l`$ is even and $`q`$ is a primitive $`2l`$-th root of unity, but this time $`A(SL_q(2))`$ has different left and right module structures, due to the appearance of some $`\pm `$ signs in the commutation relations.
Notice that, using the commutation relations to order monomials, $`A(SL_q(2))`$ can be generated over $`A(SL(2))`$ by the $`l^4`$ products
$$a^{r_a}b^{r_b}c^{r_c}d^{r_d},(0r_{}l1).$$
The determinant relation $`ad=1+qbc`$ reduces the generators above to a smaller set; all the monomials with $`r_ar_d>0`$ can be expressed as linear combinations over $`C`$ of the following set
$$a^{r_a}b^{r_b}c^{r_c},b^{r_b}c^{r_c}d^{r_d},$$
of $`2l^3`$ elements. Indeed
$$P_k(bc)=:a^kd^k=\underset{j=1}{\overset{k}{}}(1+q^{2j1}bc)=\underset{j=0}{\overset{k}{}}p_{k,j}(q)b^jc^j,$$
where
$$p_{k,j}(q)=q^{j^2}\underset{r=j+1}{\overset{k}{}}(q^{2r}1)\left(\underset{s=1}{\overset{kj}{}}(q^{2s}1)\right)^1.$$
Since $`A(SL_q(2))`$ is a finitely generated projective module over $`A(SL(2))`$, it corresponds to a locally free coherent sheaf $``$ over $`Spec(A(SL(2)))=SL(2)`$. We now think of $`SL(2)`$ as a complex manifold $`SL(2)_{an}`$ instead as an affine scheme.
Proposition 1. $`SL(2)_{an}`$ is locally free of rank $`l^3`$ and the corresponding vector bundle $`F`$ is trivial.
Proof. We cover $`SL(2)_{an}`$ with two open sets $`U_\alpha ,U_\beta `$ where $`\alpha `$ or $`\beta `$ is invertible. Since $`a^1=\alpha ^1a^{l1}`$, also $`a^1`$ exists on $`U_\alpha `$ and similarly $`b^1`$ on $`U_\beta `$. Thus on $`U_\alpha `$ we can set $`d=a^1(1+qbc)`$ and use the $`l^3`$ elements $`a^{r_a}b^{r_b}c^{r_c}`$ as generators, and on $`U_\beta `$ we set $`c=q^1b^1(ad1)`$ and use $`a^{r_a}b^{r_b}d^{r_d}`$ as generators.
This proves that $``$ is locally free of rank $`l^3`$ and hence it is a sheaf of sections of a vector bundle $`F`$ over $`SL(2)_{an}`$. But $`SL(2)_{an}`$ retracts to its maximal compact subgroup $`SU(2)S^3`$ and since for every Lie group $`G`$, $`\pi _2(G)=0`$, $`F`$ is topologically trivial. Now $`SL(2)_{an}`$ is a Stein domain, the fibre bundles over it have no moduli and are simply classified topologically. It follows that $`F`$ is also holomorphically trivial. $``$
The triviality of $`F`$ (and of $``$) implies that we should be able to find $`l^3`$ global nowhere vanishing independent sections. It is clear that the local trivializations given above are by no means global. To find out global sections of $`F`$ we work algebraically. The following theorem provides an explicit set of $`l^3`$ generators of $`A(SL_q(2))`$ over $`A(SL(2))`$.
Proposition 2. Let $`q`$ be a primitive root of unity of order $`l`$ with $`l`$ odd, or of order $`2l`$ with $`l`$ even. Then $`A(SL_q(2))`$ is a free (left and/or right) module of rank $`l^3`$ over the coordinate ring of the classical group $`A(SL(2))`$. The $`l^3`$ generators can be chosen as
$$a^mb^nc^s^{},b^nc^{s^{\prime \prime }}d^r,$$
with the integers $`m,n,r,s^{},s^{\prime \prime }`$ in the range $`1ml1;\mathrm{\hspace{0.33em}0}n,rl1;ms^{}l1;\mathrm{\hspace{0.33em}0}s^{\prime \prime }lr1`$.
Proof. For concreteness, we choose the left module structure, the right module case is analogous. Using the above expression for $`P_k(bc)`$ and noticing that $`p_{k,0}(q)=1`$, one has
$$\begin{array}{cc}\hfill b^rc^sd^{lk}=& q^{(lk)(r+s)}\delta a^kb^rc^s\underset{j=1}{\overset{k}{}}p_{k,j}(q)b^{r+j}c^{s+j}d^{lk},\hfill \\ \hfill a^{lk}b^rc^s=& q^{k(r+s)}\alpha b^rc^sd^k\underset{j=1}{\overset{k}{}}p_{k,j}(q)a^{lk}b^{r+j}c^{s+j}.\hfill \end{array}$$
The idea is to use recursively these relations to eliminate the monomials in the l.h.s.’s in terms of those in the r.h.s.’s. Of course there is a cyclicity problem, because the recursive substitutions will sooner or later bring in the r.h.s the monomial in the l.h.s. multiplied by a non constant element in $`A(SL(2))`$. We then have to put bounds on the ranges of the exponents. Let us look at the first equation. As $`d`$ occurs with the same exponent in both sides, there are no restrictions on $`k`$. Fix $`s=s_1`$ (possibly depending on $`k`$), let $`S_1(k)=\{(s_1+j)modl,\mathrm{\hspace{0.33em}0}jk\}`$, $`S_1^c(k)=Z_lS_1(k)`$. We can use the monomials $`b^rc^sd^{lk1}`$ with $`sS_1(k)`$, $`k,rZ_l`$ and $`a^kb^rc^s`$ with $`k,r,sZ_l`$ to generate the monomials in the l.h.s. with $`sS_1^c(k)`$. The same can be done for the second equation, by choosing $`s_2`$. We get a system of generators of the form $`b^rc^sd^{lk1}`$, with $`r,kZ_l`$, $`sS_1(k)`$, and $`a^kb^rc^s`$, with $`kZ_l\{0\}`$, $`rZ_l`$ and $`sS_2(lk1)`$. These are $`l^3`$ elements and we have to choose carefully $`s_1,s_2`$ to be sure that there are no relations among them. This is the case for $`s_1=0`$ and $`s_2=k`$. By a suitable relabelling of the indices, we complete the proof of the theorem. $``$
Remark 1. The results above hold for $`q`$ being a primitive root of unity of order not equal to $`2`$ mod $`4`$. In the remaining cases, $`q^{2l}=1`$, $`l`$ odd, the $`2l`$-th powers of $`a,b,c,d`$ commute but neither the determinant relation nor the coproduct close on the algebra generated by them. To remedy this mismatch one can enlarge the algebra by taking $`\alpha ,\beta ,\gamma ,\delta `$ (the $`l`$-th powers of $`a,b,c,d`$) as generators, obtaining however a non-commutative algebra.
Remark 2. In attempting to generalize the above results to $`SL_q(n)`$ we immediately face the lack of the topological argument used in the proof of Proposition 1, which in fact has been the starting point of this letter. So in principle one has to use algebra alone. Local triviality easily follows by considering the open covering of $`SL(n)`$ given by the subsets $`U_k`$ where the algebraic complements $`\widehat{T}_{k1}`$ of the elements of the first column are invertible (see for notations). On $`U_k`$ we have
$$T_{k1}=(q)^{1k}[1\underset{jk}{}(q)^{j1}T_{j1}\widehat{T}_{j1}]\widehat{T}_{k1}^1$$
and a basis of the localized module is given by the $`l^{n^21}`$ monomials $`_{(r,s)(k,1)}T_{rs}^{n_{rs}}`$, with $`n_{rs}Z_l`$. This simplifies in the special case of $`SL_q(n)`$ the proof in valid for a general quantum group. The determinant relation provides equations similar to those of Proposition 2, but the elimination problem is much subtler.
Bibliography
1. N. Andruskiewitsch, Notes on Extensions of Hopf algebras Can. J. Math. 48 3-42 (1996)
2. A. Connes, Gravity coupled with matter and the foundation of non commutative geometry, Commun. Math. Phys. 182 155–176 (1996)
3. R. Coquereaux, A. O. Garcia, R. Trinchero, Hopf stars, twisted Hopf stars and scalar products on quantum spaces, math-ph/9904037
4. L. Da̧browski, P. M. Hajac, P. Siniscalco, Explicit Hopf-Galois Description of $`SL_{e^{2\pi i/3}}(2)`$-Induced Frobenius Homomorphisms, in: Enlarged Proceedings of the ISI GUCCIA Workshop on quantum groups, non commutative geometry and fundamental physical interactions, 1997 Palermo, Italy; D.Kastler, M.Rosso, T.Schucker eds. Nova Science Pub. Inc., Commack, New-York 279–299, 1999
5. C. De Concini, V. Lyubashenko, Quantum function algebra at roots of 1, Adv. Math. 108 205-262 (1994)
6. D. Kastler, $`U_q(sl_2)`$ for $`|q|=1`$ as the complexification of a real Hopf algebra, Nachrichten der Academie der Wissenschaften in Gottingen, Vanderhoeck and Ruprecht, 1999
7. G. Lusztig, Quantum Groups at Roots of Unity, Geom. Dedicata 35 89–114 (1991)
8. S. Montgomery, H.J. Schneider, Prime ideals in Hopf Galois extensions, in prep.
9. B. Parshall, J. Wang, Quantum linear groups. American Mathematical Society Memoirs no. 439, Providence, R.I., American Mathematical Society (AMS) (1991)
10. C. Reina, A. Zampa, Quantum homogeneous spaces at roots of unity, in: Quantization, Coherent States and Poisson Structures, A. Strasburger et al. eds, Polish Sci. Pub. PWN - Warszawa 1998
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# Semi-Bounded Restrictions of Dirac Type Operators and the Unique Continuation Property
## 1 Introduction
There is a large class of unique continuation theorems for solutions to second order differential equations. For solutions to Dirac type equations on Riemannian manifolds there exists a number of methods to show the weak unique continuation property for solutions, i. e. a solution vanishing on a non-void open set vanishes identically ().
The relativistic quantum theory of a single electron is described by the Hilbert space of 4-component spinors $`:=L^2(\mathrm{}^{\mathrm{}}\mathrm{}\mathrm{}^{\mathrm{}}\mathrm{}`$ and the Dirac operator $`H_0`$ for mass $`m>0`$, which is essentially self-adjoint on the dense subspace $`C_c^{\mathrm{}}(\mathrm{}^{\mathrm{}}\mathrm{}\mathrm{}^{\mathrm{}}\mathrm{}`$. If $`P_+`$ and $`P_{}`$ are the orthogonal projections onto the spectral subspaces for the positive and negative part of the spectrum of the Dirac operator, the physical state space of the theory is the space of rays in the Hilbert space $`_+:=P_+`$. The time evolution is given by the unitary group generated by the Dirac operator restricted to $`_+`$. For such a physical state $`\psi _+,\psi =1`$, and any region $`𝒪\mathrm{}^{\mathrm{}}`$ the integral $`_𝒪\psi (x),\psi (x)d^3x`$ is the probability to find the electron in $`𝒪`$. It is known for the case of the free Dirac operator that functions in $`_+`$ have the weak unique continuation property, see \[13, Corollary 1.7\]. Physically this means, that one cannot localize a single electron in a bounded region. In the fully quantized theory of the electron field this property has important consequences. The reason for this is that the quantization procedure requires the splitting of $``$ into a direct sum $`_+_{}`$, see \[13, Section 10.1\] for details. The weak continuation property implies however that $`P_+C_c^{\mathrm{}}(𝒪,\mathrm{}^{\mathrm{}}\mathrm{}`$ is already dense in $`_+`$ for any non-void open set $`𝒪\mathrm{}^{\mathrm{}}`$. Hence the local information is destroyed by the splitting. As a consequence the vacuum vector turns out to be cyclic for the local field algebras. This Reeh-Schlieder property of the vacuum is a general property of quantum field theories in Minkowski space-time and can be proved in the Wightman framework even for interacting fields ().
Replacing $`\mathrm{}^{\mathrm{}}`$ by an arbitrary Riemannian manifold $`(M,g)`$ it is natural to ask whether the weak unique continuation property still holds not only for solutions to Dirac type equations but also, as in the case of $`\mathrm{}^{\mathrm{}}`$, for elements of certain spectral subspaces. In this note we will prove that such a property holds for a large class of generalized Dirac operators on an arbitrary connected Riemannian manifold. The physical interpretation is that single electrons cannot be localized even in curved space and the corresponding Dirac quantum field has the Reeh-Schlieder property. Our class of operators is large enough to allow coupling to arbitrary smooth Hermitian potentials. Note that in this situation only the space is curved which in a relativistic framework would correspond to an ultrastatic space-time, i. e. the space-time of the form $`\mathrm{}\times \mathrm{𝕄}`$ with metric tensor $`dt^2g`$. There already exist results on the Reeh-Schlieder property for the free scalar field over such space-times () and recently also for quite general free quantum fields over space-times admitting a time-like Killing vector field (). See also for the connection of unique continuation theorems with the Reeh-Schlieder property within non-relativistic quantum mechanics.
The weak unique continuation property for elements of spectral subspaces provides a particularly straightforward way to understand the localization problem for single particles and the Reeh-Schlieder property for the corresponding quantum fields over ultrastatic space-times. Independently such theorems play an important role in Riemannian geometry since splittings of the type $`=_+_{}`$ also occur in an essential manner in boundary value problems for Dirac equations (). This note can be considered a continuation of since we are now able to remove the assumption in that the underlying manifold be compact.
## 2 Notation and Result
Let $`(M,g)`$ be a Riemannian manifold, let $`EM`$ be a Hermitian vector bundle over $`M`$. Denote the Hermitian metric on $`E`$ by $`,`$. The space of smooth sections in $`E`$ will be denoted $`C^{\mathrm{}}(M,E)`$ and similar notation will be used for compactly supported smooth sections and for square-integrable sections. Let
$$D:C^{\mathrm{}}(M,E)C^{\mathrm{}}(M,E)$$
be a formally self-adjoint differential operator of first order. We call $`D`$ of Dirac type if its principal symbol $`\sigma _D`$ satisfies the Clifford relations, i. e.
$$\sigma _D(\xi )\sigma _D(\eta )+\sigma _D(\eta )\sigma _D(\xi )=2g(\xi ,\eta )\text{Id}_{E_p}$$
for all $`\xi `$, $`\eta T_p^{}M`$, $`pM`$. Then $`D`$ is an elliptic differential operator of first order.
Example 1 Let $`M=S^1`$ be the circle and let $`E`$ be the trivial complex line bundle over $`M`$. Then $`D=i\frac{d}{dt}`$ is of Dirac type. Fourier expansion of complex valued functions is nothing but the eigenvector expansion for $`D`$.
Example 2 More generally, any generalized Dirac operator in the sense of Gromov and Lawson () is of Dirac type.
Example 3 Let $`M=\mathrm{}^{\mathrm{}}`$ and let $`D`$ be the Dirac operator in an external field, $`D=H=H_0+V`$ in the terminology of \[13, Section 4\]. Here the potential $`V`$ to be added to the free Dirac operator $`H_0`$ may be any smooth function $`V`$ with values in Hermitian matrices. Examples include electromagnetic vector potentials, and electric and magnetic anomalous moments.
Example 4 Since the underlying manifold $`M`$ need not necessarily be complete we can also deal with singular potentials. Let us demonstrate this for the Coulomb potential. Put $`M=\mathrm{}^{\mathrm{}}\mathrm{}\{\mathrm{}\}`$, $`V(x):=\frac{\gamma }{|x|}Id`$, and $`D=H=H_0+V`$ as in the previous example. It is known (\[13, Section 4\]) that if the coupling constant $`\gamma `$ is not too large, $`|\gamma |<const`$, then $`D`$ is essentially self-adjoint on $`C_c^{\mathrm{}}(M,\mathrm{}^{\mathrm{}}\mathrm{}`$. But even for larger coupling constant for which essential self-adjointness is known to break down there still exist self-adjoint extensions and our theorem can be applied.
Example 5 Let $`(M,g)`$, $`E`$, and $`D`$ be as above with $`D`$ of Dirac type. Put $`X:=M\times \mathrm{}`$ with Riemannian metric $`ds^2=g+dt^2`$ where $`t`$ is the standard coordinate on $`\mathrm{}`$. Hence $`X`$ is the cylinder over $`M`$. The pull-back of $`E`$ to $`X`$ is again denoted by $`E`$. Define the Hermitian vector bundle $`:=EE`$ over $`X`$. Then
$$𝒟:=\left(\begin{array}{cc}D& \frac{}{t}\\ \frac{}{t}& D\end{array}\right)$$
acting on sections in $``$ is of Dirac type over $`X`$.
Now let $`D`$ be of Dirac type over the connected Riemannian manifold $`M`$. Let $`\phi `$ be in the kernel of $`D`$. By elliptic regularity theory $`\phi `$ is necessarily smooth. It is well-known that $`D`$ has the unique continuation property meaning that if $`\phi `$ vanishes on a non-empty open subset $`𝒪M`$, then $`\phi `$ vanishes everywhere, see e. g. . Of course this is also true for eigensections of $`D`$ to other eigenvalues $`\lambda 0`$ since one can simply replace $`D`$ by $`D\lambda Id`$. This unique continuation property of Dirac type operators can also be shown by passing to the square of $`D`$ and applying the classical Aronszajn theorem () for Laplace type operators.
We show here that the assumption $`D\phi =0`$ can be relaxed considerably. Roughly speaking, it is enough that $`\phi `$ lies in a $`D`$-invariant subspace of the Hilbert space of square-integrable sections $`L^2(M,E)`$ on which $`D`$ is semibounded.
For a subset $`S\mathrm{}`$ denote the characteristic function by $`\chi _S`$,
$$\chi _S(\lambda )=\{\begin{array}{cc}1,& \lambda S\hfill \\ 0,& \lambda \mathrm{}\mathrm{}\mathrm{𝕊}\mathrm{}\hfill \end{array}$$
Our main result is
Theorem Let $`M`$ be a connected Riemannian manifold , let $`E`$ be a Hermitian vector bundle over $`M`$, and let $`D`$ be of Dirac type acting on $`C_c^{\mathrm{}}(M,E)`$. Suppose $`D`$ has a self-adjoint extension in $`L^2(M,E)`$ which we denote by $`A`$. Let $`𝒪M`$ be a non-empty open subset of $`M`$. Let $`\phi L^2(M,E)`$ such that there exists $`\lambda _0\mathrm{}`$ with $`\chi _{[\lambda _0,\mathrm{})}(A)\phi =0`$ or $`\chi _{(\mathrm{},\lambda _0]}(A)\phi =0`$.
Then the following conclusion holds:
$$\phi |_𝒪=0\phi =0\text{ on all of }M.$$
Example 6 If $`\phi `$ is a square-integrable eigensection for $`D`$, then clearly, $`\chi _{[\lambda _0,\mathrm{})}(A)\varphi =0`$ for any $`\lambda _0`$ larger than the eigenvalue and we recover the classical unique continuation property of $`D`$ mentioned above.
Example 7 In case $`M=S^1`$ and $`D=i\frac{d}{dt}`$ we may apply the theorem to all $`\phi L^2(S^1,\mathrm{}\mathrm{}`$ whose Fourier coefficients $`c_n`$ vanish for all $`n`$ smaller than some $`n_0`$, i. e. all $`\phi `$ of the form
$$\phi (t)=\underset{nn_0}{}c_ne^{int}.$$
The theorem says that any such $`\phi `$ must be identically zero if it vanishes on a non-empty open subset of $`S^1`$.
Example 8 More generally, if $`M`$ is compact and connected , then the spectrum of $`A`$ is discrete. We may split $`L^2(M,E)=_0_{>0}`$ where $`_0`$ denotes the sum of all eigenspaces for the non-positive eigenvalues and similarly $`_{>0}`$ corresponds to the positive eigenvalues. This splitting is important to define the Atiyah-Patodi-Singer boundary conditions for boundary value problems for Dirac operators (). By the theorem no non-trivial $`\phi _0`$ or $`\phi _{>0}`$ can vanish on a non-empty open subset of $`M`$.
Remark Note that in the theorem no assumption of essential self-adjointness of $`D`$ on $`C_c^{\mathrm{}}(M,E)`$ is made. Dirac type operators on a complete manifold are always essentially self-adjoint on smooth compactly supported sections. The well-known proofs of this fact for generalized Dirac operators in sense of Gromov and Lawson () easily carry over to our more general class of operators, see also for a different approach. However, we only need existence of self-adjoint extensions not their uniqueness. Hence in our case the manifold $`M`$ need not necessarily be complete, compare Example 4. Also see for examples with cone singularities to which our theorem can be applied.
It should also be emphasized that $`\phi `$ need not lie in the domain of definition of $`D`$ or of $`A`$. Only square-integrability and $`\chi _{[\lambda _0,\mathrm{})}(A)\phi =0`$ or $`\chi _{(\mathrm{},\lambda _0]}(A)\phi =0`$ are assumed.
Also note that our theorem shows that the spectrum of any Dirac type operator is indeed unbounded from above and from below. Otherwise, every section in $`L^2(M,E)`$ would have the unique continuation property which is absurd.
## 3 The Proof
In this section we give the proof of the theorem. We keep the notation of the theorem and we assume $`\phi |_𝒪=0`$. We wish to show $`\phi =0`$.
Without loss of generality we may assume that $`\lambda _0=1`$ and that $`\chi _{(\mathrm{},1]}(A)\phi =0`$. In other words, $`\phi `$ lies in the image of the projection $`\chi _{(1,\mathrm{})}(A)`$. Denote the upper half plane by
$$:=\{\zeta \mathrm{}\mathrm{𝕜}\mathrm{}\mathrm{}\zeta \mathrm{}\mathrm{}\mathrm{}\}\mathrm{}$$
Step 1: The family of functions $`f_z`$,
$$f_z(\lambda )=\{\begin{array}{cc}\lambda ^1e^{iz\lambda },& \lambda 1\hfill \\ 0,& \lambda <1\hfill \end{array},$$
is uniformly bounded by 1 for all $`z\overline{}`$ and for each fixed $`\lambda `$ it is continuous in $`z`$. Therefore the family of bounded operators $`f_z(A)`$ is continuous in $`z\overline{}`$ in the strong operator topology. Hence $`zf_z(A)\phi `$ is a continuous $`L^2(M,E)`$-valued function on $`\overline{}`$.
Now fix $`z_0`$. The family of functions $`g_z`$ defined by
$$g_z(\lambda ):=\{\begin{array}{cc}\lambda ^1\frac{e^{iz\lambda }e^{iz_0\lambda }}{zz_0},& zz_0\hfill \\ ie^{iz_0\lambda },& z=z_0\hfill \end{array}$$
for $`\lambda 1`$ and $`g_z(\lambda )=0`$ for $`\lambda <1`$ is also uniformly bounded and continuous in $`z`$. This shows
$$\underset{zz_0}{lim}\frac{f_z(A)\phi f_{z_0}(A)\phi }{zz_0}$$
exists and hence $`zf_z(A)\phi `$ is holomorphic on $``$.
Step 2: Fix a non-empty relatively compact open subset $`\stackrel{~}{𝒪}𝒪`$. Let $`uC_c^{\mathrm{}}(M,E)`$ with support contained in $`\stackrel{~}{𝒪}`$. By finiteness of propagation speed, see Section 4, there exists $`\epsilon >0`$ such that the support of $`e^{itA}u`$ is contained in $`𝒪`$ for all $`t[\epsilon ,\epsilon ]`$. Hence for all $`t`$ with $`|t|\epsilon `$ we have
$`0`$ $`=`$ $`(\phi ,e^{itA}u)`$
$`=`$ $`(e^{itA}\phi ,u)`$
$`=`$ $`(Af_t(A)\phi ,u)`$
$`=`$ $`(f_t(A)\phi ,Au).`$
Here we used that by assumption $`\phi `$ lies in the image of the projection $`\chi _{(1,\mathrm{})}(A)`$ and hence $`e^{izA}\phi =Af_z(A)\phi `$. Since $`z(f_z(A)\phi ,Au)`$ is continuous on $`\overline{}`$ and holomorphic on $``$ and vanishes on $`[\epsilon ,\epsilon ]`$ an application of the Schwarz reflection principle (\[3, Lemma 2\]) shows
$$(f_z(A)\phi ,Au)=0$$
for all $`z\overline{}`$. In particular for $`z=it`$ with $`t>0`$ we obtain
$$(e^{tA}\phi ,u)=(f_{it}(A)\phi ,Au)=0$$
for all $`uC_c^{\mathrm{}}(M,E)`$ with support contained in $`\stackrel{~}{𝒪}`$. Hence $`e^{tA}\phi `$ vanishes on $`\stackrel{~}{𝒪}`$ for all $`t>0`$.
Step 3: Now look at the half cylinder $`X=M\times (0,\mathrm{})`$ with product metric $`g+dt^2`$, Hermitian vector bundle $`:=EE`$ over $`X`$, and Dirac type operator
$$𝒟:=\left(\begin{array}{cc}D& \frac{}{t}\\ \frac{}{t}& D\end{array}\right)$$
as in Example 5. We define a distributional section $`\mathrm{\Phi }_1`$ in $`E`$ over $`X`$ by
$$\mathrm{\Phi }_1(v):=_0^{\mathrm{}}(e^{tA}\phi ,v(,t))_{L^2(M,E)}𝑑t$$
for all $`vC_c^{\mathrm{}}(X,E)`$. The differential equation
$$\frac{d}{dt}e^{tA}\phi =Ae^{tA}\phi $$
shows that the distributional section
$$\mathrm{\Phi }:=\left(\begin{array}{c}\mathrm{\Phi }_1\\ \mathrm{\Phi }_1\end{array}\right)$$
in $``$ satisfies
$$𝒟\mathrm{\Phi }=0$$
in the distributional sense. In particular, $`\mathrm{\Phi }`$ is smooth by elliptic regularity theory.
Since $`\mathrm{\Phi }`$ vanishes on the open subset $`\stackrel{~}{𝒪}\times (0,\mathrm{})`$ of $`X`$ the standard unique continuation property of $`𝒟`$ implies $`\mathrm{\Phi }=0`$. Hence $`\mathrm{\Phi }_1=0`$, i. e. $`e^{tA}\phi =0`$ for all $`t>0`$, and the limit $`t0`$ yields $`\phi =0`$. $`\mathrm{}`$
## 4 Finite Propagation Speed
Inthe second step of the proof we needed what is known as “finite propagation speed” for possibly incomplete manifolds. For the sake of completeness we give a justification for this by reducing it to the compact case. A direct approach can be found in .
For any closed subset $`A`$ of $`M`$ let $`U_r(A)`$ be the closed $`r`$-neighborhood of $`A`$, i. e.
$$U_r(A)=\{xM|dist(x,A)r\}.$$
Proposition Let $`M`$ be a Riemannian manifold, let $`E`$ be a Hermitian vector bundle over $`M`$, and let $`D`$ be of Dirac type acting on $`C_c^{\mathrm{}}(M,E)`$. Suppose $`D`$ has a self-adjoint extension in $`L^2(M,E)`$ which we denote $`A`$. Let $`\phi C_c^{\mathrm{}}(M,E)`$.
Then there exists $`\epsilon >0`$ such that for all $`t[0,\epsilon )`$
$$\phi _t:=e^{itA}\phi $$
is smooth in all variables and
$$supp(\phi _t)U_t(supp(\phi )).$$
If $`M`$ is complete, then the statement is true for $`\epsilon =\mathrm{}`$.
Proof. The proposition is well-known if the underlying manifold is compact, see e. g. \[11, Prop. 5.5\] for a proof. By assumption $`supp(\phi )`$ is compact. Hence there exists $`\epsilon >0`$ such that $`U_\epsilon (supp(\phi ))`$ is still compact. We pick a compact manifold $`Y`$ containing an isometric image of $`U_\epsilon (supp(\phi ))`$ together with a Hermitian bundle and Dirac type operator extending $`E`$ and $`D`$ over $`U_\epsilon (supp(\phi ))`$. (This can e. g. be obtained by choosing an open subset $`XM`$ containing $`U_\epsilon (supp(\phi ))`$ with compact closure $`\overline{X}`$ and smooth boundary $`X`$. Then let $`Y`$ be the double of $`X`$ and smooth out all data along $`X`$.) Since $`\phi `$ has support contained in $`U_\epsilon (supp(\phi ))`$ we can consider it as a section over $`Y`$ by extending it by $`0`$. Let $`\phi _t`$ denote the solution of
$$\frac{}{t}\phi _t=iD\phi _t$$
over $`Y`$ with $`\phi _0=\phi `$. Then $`\phi _t`$ is smooth in all variables and has support contained in $`U_\epsilon (supp(\phi ))`$ for $`t[0,\epsilon )`$. By this last property we can consider $`\phi _t`$ as a smooth section over $`M`$, again extending by zero. Considered as elements of $`L^2(M,E)`$ we have
$$\frac{d}{dt}\phi _t=iD\phi _t=iA\phi _t$$
and hence $`\phi _t=e^{itA}\phi `$.
In case $`M`$ is complete $`\epsilon `$ can be chosen arbitrarily large. $`\mathrm{}`$
## 5 Concluding Remarks
Our main argument makes use of the unique continuation theorem for solutions to Dirac type equations on Riemannian manifolds (on the half cylinder $`X`$). In order to use this theorem we had to analytically continue the functions $`t(f_t(A)\phi ,Au)`$ (see Step 2 in the proof) to “imaginary time”. This was possible due to semi-boundedness of the spectrum of the Dirac operator on the spectral subspaces under consideration. In quantum field theory this method is commonly referred to as “Wick rotation”.
It should be noted that our theorem can be extended to even more general subspaces. Clearly analyticity of the function $`(f_z(A)\phi ,Au)`$ in the whole upper half plane is not necessary, but it suffices to have boundedness and analyticity in a strip of the form $`\{\zeta \mathrm{}\mathrm{𝕜}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\zeta \mathrm{}\mathrm{}\beta \}`$ for some $`\beta >0`$. Such subspaces arise naturally in thermal quantum field theory.
Fachbereich Mathematik
Universität Hamburg
Bundesstr. 55
20146 Hamburg
Germany
E-Mail: baer@math.uni-hamburg.de
WWW: http://www.math.uni-hamburg.de/home/baer/
and
Institut für theoretische Physik
Universität Leipzig
Augustusplatz 10/11
04109 Leipzig
Germany
E-Mail: alexander.strohmaier@itp.uni-leipzig.de
WWW: http://www.physik.uni-leipzig.de/~strohmai/
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# References
1. Introduction
In the Standard Model, photon and neutrino have no masses. However, the masslessness of photon is ensured by gauge invariance while the masslessness of neutrinos is only an artificial supposition. Since it is found that, the observed solar neutrino fluxes are all below the predictions based on the Standard Solar Model (SSM) , and it is very difficult to explain these flux deficits by modifying the SSM , people began to guess that neutrinos may have non-zero masses, and can oscillate from one flavor to another like that occuring in the sector of quarks.
Recently, the Super-Kamiokande experiment has provided a strong evidence for non-zero masses and oscillations of neutrinos . Because it is the first sign for new physics beyond the Standard Model, it has brought up a turbulent shock in the research field of particle physics after the announcement of the Super-Kamiokande result.
In analogy to the quark mixing in the Standard Model, it is expected that, the mixing matrix of the neutrino sector has the similar structure to that of quark sector. Then, it remind us to discuss the problem of mixing and CP violation in neutrino system naturally .
Based on the postulation put forward by one of us and the collaborators , we investigate the mixing and CP violation in neutrino system in this work. Here, we suppose that
A. The postulation on the relation between weak CP phase and the other three mixing angles in Cabibbo-Kobayashi-Maskawa matrix for quark sector be also available for neutrino system. Due to the similarity between quark sector and lepton sector, we take this supposition as reasonable.
B. Neutrinos come in three families with no additional species, sterile etc., and their masses are hierarchical with $`m_{\nu _e}`$ being smallest and $`m_{\nu _\tau }`$ largest .
C. Due to the confirmation of the Liquid Scintillator Neutrino Detector (LSND) results at Los Alamos awaiting future experiments , in the simplest explanation, solar neutrino data can be understood in terms of $`\nu _e\nu _\mu `$ oscillation with a small mass splitting not to influence atmospheric data, and atmospheric data can be explained in terms of $`\nu _\mu \nu _\tau `$ large mixing with a large mass splitting compared to the $`\nu _e\nu _\mu `$ case .
In fact, supposition A is necessary in this paper, while the other two suppositions are only for the convenience. Although we need two mixing angles precisely here, we need not limit which two of the three mixing angles. The suppositions B and C are helpful to draw a clear physical picture.
2. The postulation on the weak CP phase
In previous works , we have postulated that, the weak CP phase $`\delta _{13}`$ and the other three mixing angles $`\theta _{12},\theta _{23}`$ and $`\theta _{13}`$ satisfy
$$\mathrm{sin}\delta _{13}=\frac{(1+s_{12}+s_{23}+s_{13})\sqrt{1s_{12}^2s_{23}^2s_{13}^2+2s_{12}s_{23}s_{13}}}{(1+s_{12})(1+s_{23})(1+s_{13})}$$
(1)
where, the convention $`s_{ij}=\mathrm{sin}\theta _{ij},c_{ij}=\mathrm{cos}\theta _{ij}`$ (the ”generation” labels $`i,j=1,2,3`$) are used and, $`\delta _{13}`$ and $`\theta _{ij}`$ are those present in the standard parametrization of the Cabibbo-Kaboyashi-Maskawa (CKM) matrix
$$V_{KM}=\left(\begin{array}{ccc}c_{12}c_{13}& s_{12}c_{13}& s_{13}e^{i\delta _{13}}\\ s_{12}c_{23}c_{12}s_{23}s_{13}e^{i\delta _{13}}& c_{12}c_{23}s_{12}s_{23}s_{13}e^{i\delta _{13}}& s_{23}c_{13}\\ s_{12}s_{23}c_{12}c_{23}s_{13}e^{i\delta _{13}}& c_{12}s_{23}s_{12}c_{23}s_{13}e^{i\delta _{13}}& c_{23}c_{13}\end{array}\right)$$
(2)
with the real angles $`\theta _{12},\theta _{23}`$ and $`\theta _{13}`$ can all be made to lie in the first quadrant. The phase $`\delta _{13}`$ lies in the range $`0<\delta _{13}<2\pi `$. In following, we will fix the three angles $`\theta _{12},\theta _{23}`$ and $`\theta _{13}`$ in the first quadrant.
The geometry meaning of Eq.(1) is evident. $`\delta _{13}`$ is the solid angle enclosed by $`(\pi /2\theta _{12}),(\pi /2\theta _{23})`$ and $`(\pi /2\theta _{13})`$ standing on a same point, or, the area to which the solid angle corresponding on a unit spherical surface.
Hence, to make $`(\pi /2\theta _{12}),(\pi /2\theta _{23})`$ and $`(\pi /2\theta _{13})`$ be able to enclose a solid angle, the following relation must be hold.
$$(\frac{\pi }{2}\theta _{ij})+(\frac{\pi }{2}\theta _{jk})(\frac{\pi }{2}\theta _{ki})(ijki=1,2,3.\theta _{ij}=\theta _{ji})$$
(3)
Eq.(3) and Eq.(1) are the most important constraints in this work, on which the following discussions are based.
3. The relevant experimental results on neutrino masses and mixing
The recent analysis made by Hata and Langacker gives viable solutions for the BAHCALL SSM . With the Mikheyev-Smirnov-Wolfenstein (MSW) mechanism being considered , they give the small-mixing solution
$$\delta m_{sol}^25\times 10^6eV^2sin^22\theta _{sol}8\times 10^3$$
(4)
and the large-mixing solution
$$\delta m_{sol}^21.6\times 10^5eV^2sin^22\theta _{sol}0.6.$$
(5)
Vacuum oscillation also provide solutions
$$\delta m_{sol}^2=(58)\times 10^{11}eV^2sin^22\theta _{sol}=0.651.$$
(6)
The atmospheric neutrino data from Super-Kamiokande etc. imply that the parameters of the $`\nu _\mu \nu _\tau `$ oscillation of the atmospheric neutrinos are
$$10^4\delta m_{atm}^210^2eV^2sin^22\theta _{atm}1.$$
(7)
We have known that, the small-mixing solution causes the energy-spectrum distortion while the large-mixing solution causes the day-night flux difference, and, the vacuum-oscillations cause seasonal variation of the $`{}_{}{}^{7}B_{e}^{}`$ solar neutrino flux .
In the next section, we will see that, with the Super-Kamiokande results about the $`\nu _\mu \nu _\tau `$ oscillation being admitted, the small-mixing solution differents from the large-mixing solution in CP violation greatly. Because the discussions here are only related to the mixing angles, and the MSW large-mixing solution gives about the same mixing as that given by the vacuum oscillation solution, so, we only discuss the two cases of small- and large-mixing as indicated by the MSW solutions.
4. Some predictions on the mixing and CP violation in neutrino system
Let us return to the constraint Eq.(3) and recall the clear geometry meaning of Eq.(1), what can we extract from them? As we have supposed in section 1, Eq.(7) tells us
$$\theta _{\mu \tau }\pi /4.$$
(8)
Keep this point in mind, then, by the use of Eq.(3), Eq.(4-6) will give restriction on the mixing angle between $`\nu _e`$ and $`\nu _\tau `$. In the meantime, Eq.(1) tell us some information on the CP violation in neutrino system. Let us talk more detail about these two problems.
(1) About the mixing angle between $`\nu _e`$ and $`\nu _\tau `$.
From Eq.(3), we have
$$|(\frac{\pi }{2}\theta _{e\mu })(\frac{\pi }{2}\theta _{\mu \tau })|(\frac{\pi }{2}\theta _{e\tau })Min(\pi /2,(\frac{\pi }{2}\theta _{e\mu })+(\frac{\pi }{2}\theta _{\mu \tau }))$$
(9)
Note that, Eqs.(4-5) imply for small- and large-mixing solutions
$$\theta _{e\mu }0.045\mathrm{or}\pi /20.045$$
and
$$\theta _{e\mu }0.443\mathrm{or}\pi /20.443$$
respectively. Considering Eq.(8), then we obtain
$$0\theta _{e\tau }\pi /4+0.045$$
(10)
or
$$\pi /40.045\theta _{e\tau }\pi /4+0.045$$
(11)
for the case of small-mixing solution. And
$$0\theta _{e\tau }\pi /4+0.443$$
(12)
or
$$\pi /40.443\theta _{e\tau }\pi /4+0.443$$
(13)
for the case of large-mixing solution. Eq.(10-13) are the constraints on the mixing angle $`\theta _{e\tau }`$.
(2) About the CP violation in neutrino system.
Now, we discuss the CP violation in neutrino system. As is well known, all the CP violated observables are proportional to the Jarlskog invariant
$$J_{CP}=s_{12}s_{13}s_{23}c_{12}c_{13}^2c_{23}s_{\delta _{13}}$$
(14)
where, $`s_{\delta _{13}}sin\delta _{13}`$ with $`\delta _{13}`$ the weak CP phase presenting in the CKM matrix of the lepton sector. And, we instead $`\theta _{12},\theta _{13}`$ and $`\theta _{23}`$ by $`\theta _{e\mu },\theta _{e\tau }`$ and $`\theta _{\mu \tau }`$ respectively in following.
With $`\theta _{\mu \tau }`$ given by the Super-Kamiokande data definately, the small- and large-mixing solutions and the vacuum oscillation will give the pemitted ranges for $`\theta _{e\tau }`$ correspondingly. Substitute Eq.(1) into Eq.(14), we obtain $`J_{CP}`$ as a function of $`\theta _{e\tau }`$. Then, we can draw the curve which $`J_{CP}`$ versus $`\theta _{e\tau }`$. The results are shown in Fig.(1). From the figure, we find that, the small-mixing solution corresponds to small CP violation while the large-mixing solution corresponds to large CP violation.
For the small-mixing solution, $`J_{CP}`$ is very small. The maximum of $`J_{CP}`$ is about $`5\times 10^3`$ when $`\theta _{e\mu }`$ nearing to $`0`$, and $`1.5\times 10^4`$ when $`\theta _{e\mu }`$ nearing to $`\pi /2`$.
For the large-mixing solution, $`J_{CP}`$ is relative large. The maximum of $`J_{CP}`$ is about $`1.5\times 10^2`$ when $`\theta _{e\mu }`$ nearing to $`0.443`$, and $`3.2\times 10^2`$ when $`\theta _{e\mu }`$ nearing to $`\pi /20.443`$.
Now, it is evident that, if the future experiments on the CP violation in neutrino system tell us that $`J_{CP}`$ is larger than $`5\times 10^3`$, the mixing between $`\nu _e`$ and $`\nu _\mu `$ must be large. Recalling the same observable in quark sector obtained via $`K^0\overline{K^0}`$ system is about $`10^4`$, if it is the same order of magnitude in the neutrino system, then, the mixing between $`\nu _e`$ and $`\nu _\tau `$ is either around $`0.8`$ with a relative narrow window or nearing to $`0`$ or $`(\pi /4+0.443)`$ very closely.
On the other hand, either $`\theta _{e\mu }`$ takes $`0.443`$ or it takes $`(\pi /20.443)`$ for large-mixing solution can also be distingushed to some extent. For example, if experiment tells us $`J_{CP}>0.015`$, it must be $`\theta _{e\mu }0.443`$.
Finally, with the CHOOZ result being considered, that is, $`\mathrm{sin}^22\theta _{e\tau }<0.2`$, then we have
$$0<\theta _{e\tau }<0.23\mathrm{or}\pi /20.23<\theta _{e\tau }<\pi /2.$$
Based on the above constraint, we can see from Fig.(1), firstly, the small-mixing solution with $`\theta _{e\mu }(\pi /20.045)`$ and the large-mixing solution with $`\theta _{e\mu }(\pi /20.443)`$ have been excluded. Secondly, the possible domain $`\pi /20.23<\theta _{e\tau }<\pi /2`$ should also be eliminated. Thirdly, $`J_{CP}`$ can still be large to about $`0.02`$ in the case of large-mixing solution with $`\theta _{e\mu }0.443`$, especially, the larger $`J_{CP}`$ than $`3\times 10^3`$ will exclude the possibility of small-mixing solution finally.
Maybe, the most interesting conclusion is about the bi-mixmal mixing . From the above analysis, we can see that, in most of the permitted range of the third angle - $`\theta _{e\tau }`$, there will be a relative large CP violation.
Suppose that $`\theta _{e\mu }=\theta _{\mu \tau }=\pi /4`$, similarly, we can draw the curve which $`J_{CP}`$ versus $`\theta _{e\tau }`$ in the permitted range of $`\theta _{e\tau }`$. The result is shown in Fig.(2). We find that, except for $`\theta _{e\tau }`$ nearing to $`0`$ or $`\pi /2`$, $`J_{CP}`$ is larger than $`10^3`$ in most of the permitted range of $`\theta _{e\tau }`$. And, the maximum of $`J_{CP}`$ can reach to about $`0.018`$ when the CHOOZ result is considered.
5. Conclusions
Starting from the postulation on the relation between weak CP phase and the other three mixing angles in the CKM matrix, we have investigated the mixing and CP violation in the neutrino system.
We suppose that, the solar neutrino problem be understood in terms of $`\nu _e\nu _\mu `$ oscillation with a small mass splitting. With the definite large mixing between $`\nu _\mu \nu _\tau `$ indicated by the Super-Kamiokande data, and the CHOOZ result being considered, we obtain the relevant constraints on the mixing between $`\nu _e`$ and $`\nu _\tau `$. We find, $`0\theta _{e\tau }0.23`$ is permitted by the small- and the large-mixing solutions.
Besides, the mixing between $`\nu _e`$ and $`\nu _\mu `$ is limited as $`\theta _{e\mu }0.045`$ for the small-mixing solution or $`\theta _{e\mu }0.443`$ for the large-mixing solution. And, a larger $`J_{CP}`$ than $`3\times 10^3`$ will finally exclude the possibility of small-mixing solution.
Furthermore, if the suppositions B and C in section 1 holds, a $`J_{CP}`$ larger than $`3\times 10^3`$ implys the large-mixing solution for solar neutrino problem. And, if it takes the same order for $`J_{CP}`$ in the neutrino system as the one in quark system, the mixing between $`\nu _e`$ and $`\nu _\mu `$ will be very small $`(10^2\mathrm{or}\mathrm{less})`$.
For the case of bi-maximal mixing, we predict a large CP violation in neutrino system with $`J_{CP}`$ larger than $`10^3`$, except the third mixing angle approachs to $`0`$ or $`\pi /2`$ very closely.
Finally, although we have made some suppositions in this work, the basis and the method used here is actually valid for a more general discussion.
Notes: In fact, this work has been finished and submitted before the last July. Due to some reason, we have not put it on the net in time. Just two days ago, when we noted the paper hep-ph/0004020 by Sin Kyu Kang, C. S. Kim and J. D. Kim and found their results are almost the same as those of us, although they based on some concrete model while we only started out from our postulation, we are encouraged to post this short paper.
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# 0. FORWARD
## 0. FORWARD
The energy quanta occured in 1900 in the work of Max Planck (Nobel prize, 1918) on the black body electromagnetic radiation. Planck’s “quanta of light” have been used by Einstein (Nobel prize, 1921) to explain the photoelectric effect, but the first “quantization” of a quantity having units of action (the angular momentum) belongs to Niels Bohr (Nobel Prize, 1922). This opened the road to the universalization of quanta, since the action is the basic functional to describe any type of motion. However, only in the 1920’s the formalism of quantum mechanics has been developed in a systematic manner. The remarkable works of that decade contributed in a decisive way to the rising of quantum mechanics at the level of fundamental theory of the universe, with successful technological applications. Moreover, it is quite probable that many of the cosmological misteries may be disentangled by means of various quantization procedures of the gravitational field, advancing our understanding of the origins of the universe. On the other hand, in recent years, there is a strong surge of activity in the information aspect of quantum mechanics. This aspect, which was generally ignored in the past, aims at a very attractive “quantum computer” technology.
At the philosophical level, the famous paradoxes of quantum mechanics, which are perfect examples of the difficulties of ‘quantum’ thinking, are actively pursued ever since they have been first posed. Perhaps the most famous of them is the EPR paradox (Einstein, Podolsky, Rosen, 1935) on the existence of elements of physical reality, or in EPR words: “If, without in any way disturbing a system, we can predict with certainty (i.e., with probability equal to unity) the value of a physical quantity, then there exists an element of physical reality corresponding to this physical quantity.” Another famous paradox is that of Schrödinger’s cat which is related to the fundamental quantum property of entanglement and the way we understand and detect it. What one should emphasize is that all these delicate points are the sourse of many interesting and innovative experiments (such as the so-called “teleportation” of quantum states) pushing up the technology.
Here, I present eight elementary topics in nonrelativistic quantum mechanics from a course in Spanish (“castellano”) on quantum mechanics that I taught in the Instituto de Física, Universidad de Guanajuato (IFUG), León, Mexico, during the semesters of 1998.
Haret C. Rosu
## 1. THE QUANTUM POSTULATES
The following six postulates can be considered as the basis for theory and experiment in quantum mechanics in its most used form, which is known as the Copenhagen interpretation.
1. To any physical quantity L, which is well defined at the classical level, one can associate a hermitic operator $`\widehat{L}`$.
1. To any stationary physical state in which a quantum system can be found one can associate a (normalized) wavefunction. $`\psi `$ ($`\psi _^2^2=1`$).
1. In (appropriate) experiments, the physical quantity L can take only the eigenvalues of $`\widehat{L}`$. Therefore the eigenvalues should be real, a condition which is fulfilled only by hermitic operators.
1. What one measures is always the mean value $`\overline{L}`$ of the physical quantity (i.e., operator) $`\widehat{L}`$ in a state $`\psi _n`$, which, theoretically speaking, is the corresponding diagonal matrix element
$`\psi _n\widehat{L}\psi _n=\overline{L}`$.
1. The matrix elements of the operators corresponding to the cartesian coordinate and momentum, $`\widehat{x_i}`$ and $`\widehat{p_k}`$, when calculated with the wavefunctions $`f`$ and $`g`$ satisfy the Hamilton equations of motion of classical mechanics in the form:
$$\frac{d}{dt}f\widehat{p_i}g=f\frac{\widehat{H}}{\widehat{x_i}}g$$
$$\frac{d}{dt}f\widehat{x_i}g=f\frac{\widehat{H}}{\widehat{p_i}}g,$$
where $`\widehat{H}`$ is the hamiltonian operator, whereas the derivatives with respect to operators are defined as at point 3 of this chapter.
1. The operators $`\widehat{p_i}`$ and $`\widehat{x_k}`$ have the following commutators:
$$[\widehat{p_i},\widehat{x_k}]=i\mathrm{}\delta _{ik},$$
$$[\widehat{p_i},\widehat{p_k}]=0,$$
$$[\widehat{x_i},\widehat{x_k}]=0$$
$`\mathrm{}=h/2\pi =1.0546\times 10^{27}`$ erg.sec.
1. The correspondence between classical and quantum quantities
This can be done by substituting $`x_i`$, $`p_k`$ with $`\widehat{x_i}`$ $`\widehat{p_k}`$. The function L is supposed to be analytic (i.e., it can be developed in Taylor series). If the L function does not contain mixed products $`x_kp_k`$, the operator $`\widehat{L}`$ is directly hermitic.
Exemple:
$`T=(_i^3p_i^2)/2m`$ $``$ $`\widehat{T}=(_i^3\widehat{p}^2)/2m`$.
If L contains mixed products $`x_ip_i`$ and higher powers of them, $`\widehat{L}`$ is not hermitic, and in this case L is substituted by $`\widehat{\mathrm{\Lambda }}`$, the hermitic part of $`\widehat{L}`$ ($`\widehat{\mathrm{\Lambda }}`$ is an autoadjunct operator).
Exemple:
$`w(x_i,p_i)=_ip_ix_i`$ $``$ $`\widehat{w}=1/2_i^3(\widehat{p_i}\widehat{x_i}+\widehat{x_i}\widehat{p_i})`$.
In addition, one can see that we have no time operator. In quantum mechanics, time is only a parameter that can be introduced in many ways. This is so because time does not depend on the canonical variables, merely the latter depend on time.
1. Probability in the discrete part of the spectrum
If $`\psi _n`$ is an eigenfunction of the operator $`\widehat{L}`$, then:
$`\overline{L}=<n\widehat{L}n>=<n\lambda _nn>=\lambda _n<nn>=\delta _{nn}\lambda _n=\lambda _n`$.
Moreover, one can prove that $`\overline{L}^k=(\lambda _n)^k`$.
If the function $`\varphi `$ is not an eigenfunction of $`\widehat{L}`$, one can make use of the expansion in the complete system of eigenfunctions of $`\widehat{L}`$ to get:
$`\widehat{L}\psi _n=\lambda _n\psi _n`$, $`\varphi =_na_n\psi _n`$
and combining these two relationships one gets:
$`\widehat{L}\varphi =_n\lambda _na_n\psi _n`$.
In this way, one is able to calculate the matrix elements of the operator $`\widehat{L}`$:
$`\varphi \widehat{L}\varphi =_{n,m}a_m^{}a_n\lambda _nmn=_ma_m^2\lambda _m`$,
telling us that the result of the experiment is $`\lambda _m`$ with a probability $`a_m^2`$.
If the spectrum is discrete, according to P4 this means that $`a_m^2`$, that is the coefficients of the expansion in a complete set of eigenfunctions, determine the probabilitities to observe the eigenvalue $`\lambda _n`$.
If the spectrum is continuous, using the following definition
$`\varphi (\tau )=a(\lambda )\psi (\tau ,\lambda )𝑑\lambda `$,
one can calculate the matrix elements in the continuous part of the spectrum
$`\varphi \widehat{L}\varphi `$
$`=𝑑\tau a^{}(\lambda )\psi ^{}(\tau ,\lambda )𝑑\lambda \mu a(\mu )\psi (\tau ,\mu )𝑑\mu `$
$`=a^{}a(\mu )\mu \psi ^{}(\tau ,\lambda )\psi (tau,\mu )𝑑\lambda 𝑑\mu 𝑑\tau `$
$`=a^{}(\lambda )a(\mu )\mu \delta (\lambda \mu )𝑑\lambda 𝑑\mu `$
$`=a^{}(\lambda )a(\lambda )\lambda 𝑑\lambda `$
$`=a(\lambda )^2\lambda 𝑑\lambda `$.
In the continuous case, $`a(\lambda )^2`$ should be understood as the probability density for observing the eigenvalue $`\lambda `$ belonging to the continuous spectrum. Moreover, the following holds
$`\overline{L}=\varphi \widehat{L}\varphi `$.
One usually says that $`\mu \mathrm{\Phi }`$ is the representation of $`\mathrm{\Phi }`$ in the representation $`\mu `$, where $`\mu `$ is an eigenvector of $`\widehat{M}`$.
1. Definition of the derivate with respect to an operator
$`\frac{F(\widehat{L})}{\widehat{L}}=\mathrm{lim}_ϵ\mathrm{}\frac{F(\widehat{L}+ϵ\widehat{I})F(\widehat{L})}{ϵ}.`$
1. The operators of cartesian momenta
Which is the explicit form of $`\widehat{p_1}`$, $`\widehat{p_2}`$ and $`\widehat{p_3}`$, if the arguments of the wavefunctions are the cartesian coordinates $`x_i`$ ?
Let us consider the following commutator:
$`[\widehat{p_i},\widehat{x_i}^2]=\widehat{p_i}\widehat{x_i}^2\widehat{x_i}^2\widehat{p_i}`$
$`=\widehat{p_i}\widehat{x_i}\widehat{x_i}\widehat{x_i}\widehat{p_i}\widehat{x_i}+\widehat{x_i}\widehat{p_i}\widehat{x_i}\widehat{x_i}\widehat{x_i}\widehat{p_i}`$
$`=(\widehat{p_i}\widehat{x_i}\widehat{x_i}\widehat{p_i})\widehat{x_i}+\widehat{x_i}(\widehat{p_i}\widehat{x_i}\widehat{x_i}\widehat{p_i})`$
$`=[\widehat{p_i},\widehat{x_i}]\widehat{x_i}+\widehat{x_i}[\widehat{p_i},\widehat{x_i}]`$
$`=i\mathrm{}\widehat{x_i}i\mathrm{}\widehat{x_i}=2i\mathrm{}\widehat{x_i}`$.
In general, the following holds:
$`\widehat{p_i}\widehat{x_i}^n\widehat{x_i}^n\widehat{p_i}=ni\mathrm{}\widehat{x_i}^{n1}.`$
Then, for all analytic functions we have:
$`\widehat{p_i}\psi (x)\psi (x)\widehat{p_i}=i\mathrm{}\frac{\psi }{x_i}`$.
Now, let $`\widehat{p_i}\varphi =f(x_1,x_2,x_3)`$ be the manner in which $`\widehat{p_i}`$ acts on $`\varphi (x_1,x_2,x_3)=1`$. Then:
$`\widehat{p_i}\psi =i\mathrm{}\frac{\psi }{x_1}+f_1\psi `$ and similar relationships hold for $`x_2`$ and $`x_3`$.
From the commutator $`[\widehat{p_i},\widehat{p_k}]=0`$ it is easy to get $`\times \stackrel{}{f}=0`$ and therefore $`f_i=_iF`$.
The most general form of $`\widehat{p_i}`$ is $`\widehat{p_i}=i\mathrm{}\frac{}{x_i}+\frac{F}{x_i}`$, where $`F`$ is an arbitrary function. The function $`F`$ can be eliminated by the unitary transformaton $`\widehat{U}^{}=\mathrm{exp}(\frac{i}{\mathrm{}}F)`$.
$`\widehat{p_i}=\widehat{U}^{}(i\mathrm{}\frac{}{x_i}+\frac{F}{x_i})\widehat{U}`$
$`=\mathrm{exp}^{\frac{i}{\mathrm{}}F}(i\mathrm{}\frac{}{x_i}+\frac{F}{x_i})\mathrm{exp}^{\frac{i}{\mathrm{}}F}`$
$`=i\mathrm{}\frac{}{x_i}`$
leading to $`\widehat{p_i}=i\mathrm{}\frac{}{x_i}`$ $``$ $`\widehat{p}=i\mathrm{}`$.
1. Calculation of the normalization constant
Any wavefunction $`\psi (x)`$ $``$ $`^2`$ of variable $`x`$ can be written in the form:
$`\psi (x)=\delta (x\xi )\psi (\xi )𝑑\xi `$
that can be considered as the expansion of $`\psi `$ in eigenfunction of the operator position (cartesian coordinate) $`\widehat{x}\delta (x\xi )=\xi (x\xi )`$. Thus, $`\psi (x)^2`$ is the probability density of the coordinate in the state $`\psi (x)`$. From here one gets the interpretation of the norm
$`\psi (x)^2=\psi (x)^2𝑑x=1`$.
Intuitively, this relationship tells us that the system described by $`\psi (x)`$ should be encountered at a certain point on the real axis, although we can know only approximately the location.
The eigenfunctions of the momentum operator are:
$`i\mathrm{}\frac{\psi }{x_i}=p_i\psi `$, and by integrating one gets $`\psi (x_i)=A\mathrm{exp}^{\frac{i}{\mathrm{}}p_ix_i}`$. $`x`$ and $`p`$ have continuous spectra and therefore the normalization is performed by means of the Dirac delta function.
Which is the explicit way of getting the normalization constant ?
This is a matter of the following Fourier transforms:
$`f(k)=g(x)\mathrm{exp}^{ikx}dx`$, $`g(x)=\frac{1}{2\pi }f(k)\mathrm{exp}^{ikx}dk.`$
It can also be obtained with the following procedure. Consider the unnormalized wavefunction of the free particle
$`\varphi _p(x)=A\mathrm{exp}^{\frac{ipx}{\mathrm{}}}`$ and the formula
$`\delta (xx^{^{}})=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\mathrm{exp}^{ik(xx^{^{}})}dx.`$
One can see that
$`_{\mathrm{}}^{\mathrm{}}\varphi _p^{^{}}^{}(x)\varphi _p(x)𝑑x`$
$`=_{\mathrm{}}^{\mathrm{}}A^{}\mathrm{exp}^{\frac{ip^{^{}}x}{\mathrm{}}}A\mathrm{exp}^{\frac{ipx}{\mathrm{}}}dx`$
$`=_{\mathrm{}}^{\mathrm{}}A^2\mathrm{exp}^{\frac{ix(pp^{^{}})}{\mathrm{}}}dx`$
$`=A^2\mathrm{}_{\mathrm{}}^{\mathrm{}}\mathrm{exp}^{\frac{ix(pp^{^{}})}{\mathrm{}}}d\frac{x}{\mathrm{}}`$
$`=2\pi \mathrm{}A^2\delta (pp^{^{}})`$
and therefore the normalization constant is:
$`A=\frac{1}{\sqrt{2\pi \mathrm{}}}`$.
Moreover, the eigenfunctions of the momentum form a complete system (in the sense of the continuous case) for all functions of the $`^2`$ class.
$`\psi (x)=\frac{1}{\sqrt{2\pi \mathrm{}}}a(p)\mathrm{exp}^{\frac{ipx}{\mathrm{}}}dp`$
$`a(p)=\frac{1}{\sqrt{2\pi \mathrm{}}}\psi (x)\mathrm{exp}^{\frac{ipx}{\mathrm{}}}dx`$.
These formulae provide the connection between the x and p representations.
1. The momentum (p) representation
The explicit form of the operators $`\widehat{p_i}`$ and $`\widehat{x_k}`$ can be obtained either from the commutation relationships or through the usage of the kernels
$`x(p,\beta )=U^{}xU=\frac{1}{2\pi \mathrm{}}\mathrm{exp}^{\frac{ipx}{\mathrm{}}}x\mathrm{exp}^{\frac{i\beta x}{\mathrm{}}}dx`$
$`=\frac{1}{2\pi \mathrm{}}\mathrm{exp}^{\frac{ipx}{\mathrm{}}}(i\mathrm{}\frac{}{\beta }\mathrm{exp}^{\frac{i\beta x}{\mathrm{}}})`$.
The integral is of the form: $`M(\lambda ,\lambda ^{^{}})=U^{}(\lambda ,x)\widehat{M}U(\lambda ^{^{}},x)𝑑x`$, and using $`\widehat{x}f=x(x,\xi )f(\xi )𝑑\xi `$, the action of $`\widehat{x}`$ on $`a(p)`$ $``$ $`^2`$ is:
$`\widehat{x}a(p)=x(p,\beta )a(\beta )𝑑\beta `$
$`=(\frac{1}{2\pi \mathrm{}}\mathrm{exp}^{\frac{ipx}{\mathrm{}}}(i\mathrm{}\frac{}{\beta }\mathrm{exp}^{\frac{i\beta x}{\mathrm{}}})𝑑x)a(\beta )𝑑\beta `$
$`=\frac{i}{2\pi }\mathrm{exp}^{\frac{ipx}{\mathrm{}}}\frac{}{\beta }\mathrm{exp}^{\frac{i\beta x}{\mathrm{}}}a(\beta )𝑑x𝑑\beta `$
$`=\frac{i\mathrm{}}{2\pi }\mathrm{exp}^{\frac{ipx}{\mathrm{}}}\frac{}{\beta }\mathrm{exp}^{\frac{i\beta x}{\mathrm{}}}a(\beta )𝑑\frac{x}{\mathrm{}}𝑑\beta `$
$`=\frac{i\mathrm{}}{2\pi }\mathrm{exp}^{\frac{ix(\beta p)}{\mathrm{}}}\frac{}{\beta }a(\beta )𝑑\frac{x}{\mathrm{}}𝑑\beta `$
$`=i\mathrm{}\frac{a(p)}{\beta }\delta (\beta p)𝑑\beta =i\mathrm{}\frac{a(p)}{p}`$,
where $`\delta (\beta p)=\frac{1}{2\pi }\mathrm{exp}^{\frac{ix(\beta p)}{\mathrm{}}}d\frac{x}{\mathrm{}}`$.
The momentum operator in the p representation is defined by the kernel:
$`p(p,\beta )=\widehat{U}^{}p\widehat{U}`$
$`=\frac{1}{2\pi \mathrm{}}\mathrm{exp}^{\frac{ipx}{\mathrm{}}}(i\mathrm{}\frac{}{x})\mathrm{exp}^{\frac{i\beta x}{\mathrm{}}}dx`$
$`=\frac{1}{2\pi \mathrm{}}\mathrm{exp}^{\frac{ipx}{\mathrm{}}}\beta \mathrm{exp}^{\frac{i\beta x}{\mathrm{}}}dx=\beta \lambda (p\beta )`$
leading to $`\widehat{p}a(p)=pa(p)`$.
It is worth noting that $`\widehat{x}`$ and $`\widehat{p}`$, although hermitic operators for all f(x) $``$ $`^2`$, are not hermitic for their own eigenfunctions.
If $`\widehat{p}a(p)=p_oa(p)`$ and $`\widehat{x}=\widehat{x}^{}`$ $`\widehat{p}=\widehat{p}^{}`$, then
$`<a\widehat{p}\widehat{x}a><a\widehat{x}\widehat{p}a>=i\mathrm{}<aa>`$
$`p_o[<a\widehat{x}a><a\widehat{x}a>]=i\mathrm{}<aa>`$
$`p_o[<a\widehat{x}a><a\widehat{x}a>]=0`$
The left hand side is zero, whereas the right hand side is indefinite, which is a contradiction.
1. Schrödinger and Heisenberg representations
The equations of motion given by P5 have different interpretations because in the expression $`\frac{d}{dt}f\widehat{L}f`$ one can consider the temporal dependence as belonging either to the wavefunctions or operators, or both to wavefunctions and operators. We shall consider herein only the first two cases.
* For an operator depending on time $`\widehat{O}=\widehat{O(t)}`$ we have:
$`\widehat{p_i}=\frac{\widehat{H}}{\widehat{x_i}}`$, $`\widehat{x_i}=\frac{\widehat{H}}{\widehat{p_i}}`$
$`[\widehat{p},f]=\widehat{p}ff\widehat{p}=i\mathrm{}\frac{f}{\widehat{x_i}}`$
$`[\widehat{x},f]=\widehat{x}ff\widehat{x}=i\mathrm{}\frac{f}{\widehat{p_i}}`$
and the Heisenberg equations of motion are easily obtained:
$`\widehat{p_i}=\frac{i}{\mathrm{}}[\widehat{p},\widehat{H}]`$, $`\widehat{x_i}=\frac{i}{\mathrm{}}[\widehat{x},\widehat{H}]`$.
* If the wavefunctions are time dependent one can still use $`\widehat{p_i}=\frac{i}{\mathrm{}}[\widehat{p_i},\widehat{H}]`$, because being a consequence of the commutation relations it does not depend on representation
$`\frac{d}{dt}<f\widehat{p_i}g>=\frac{i}{\mathrm{}}<f[\widehat{p},\widehat{H}]g>`$.
If now $`\widehat{p_i}`$ and $`\widehat{H}`$ do not depend on time, taking into account the hermiticity, one gets:
$`(\frac{f}{t},\widehat{p_i}g)+(\widehat{p_i}f,\frac{g}{t})`$
$`=\frac{i}{\mathrm{}}(f,\widehat{p_i}\widehat{H}g)+\frac{i}{\mathrm{}}(f,\widehat{H}\widehat{p_i}g)`$
$`=\frac{i}{\mathrm{}}(\widehat{p}f,\widehat{H}g)+\frac{i}{\mathrm{}}(\widehat{H}f,\widehat{p_i}g)`$
$`(\frac{f}{t}+\frac{i}{\mathrm{}}\widehat{H}f,\widehat{p_i}g)+(\widehat{p_i}f,\frac{g}{t}\frac{i}{\mathrm{}}\widehat{H}g)=0`$
The latter relationship holds for any pair of functions $`f(x)`$ and $`g(x)`$ at the initial moment if each of them satisfies the equation
$`i\mathrm{}\frac{\psi }{t}=H\psi `$.
This is the Schrödinger equation. It describes the system by means of time-independent operators and makes up the so-called Schrödinger representation.
In both representations the temporal evolution of the system is characterized by the operator $`\widehat{H}`$, which can be obtained from Hamilton’s function of classical mechanics.
Exemple: $`\widehat{H}`$ for a particle in a potential $`U(x_1,x_2,x_3)`$ we have:
$`\widehat{H}=\frac{\widehat{p^2}}{2m}+U(x_1,x_2,x_3)`$, which in the x representation is:
$`\widehat{H}=\frac{\mathrm{}^2}{2m}_x^2+U(x_1,x_2,x_3)`$.
1. The connection between the S and H representations
P5 is correct in both Schrödinger’s representation and Heisenberg’s. This is why, the mean value of any observable coincides in the two representations. Thus, there is a unitary transformation that can be used for passing from one to the other. Such a transformation is of the form $`\widehat{s}^{}=\mathrm{exp}^{\frac{i\widehat{H}t}{\mathrm{}}}`$. In order to pass to the Schrödinger representation one should use the Heisenberg transform $`\psi =\widehat{s^{}}f`$ with $`f`$ and $`\widehat{L}`$, whereas to pass to Heisenberg’s representation the Schrödinger transform $`\widehat{\mathrm{\Lambda }}=\widehat{s^{}}\widehat{L}\widehat{s}`$ with $`\psi `$ and $`\widehat{\mathrm{\Lambda }}`$ is of usage. One can obtain the Schrödinger equation as follows: since in the transformation $`\psi =\widehat{s^{}}f`$ the function $`f`$ does not depend on time, we shall derivate the transformation with respect to time to get:
$`\frac{\psi }{t}=\frac{s^{}}{t}f=\frac{}{t}(\mathrm{exp}^{\frac{i\widehat{H}t}{\mathrm{}}})f=\frac{i}{\mathrm{}}\widehat{H}\mathrm{exp}^{\frac{i\widehat{H}t}{\mathrm{}}}f=\frac{i}{\mathrm{}}\widehat{H}\widehat{s^{}}f=\frac{i}{\mathrm{}}\widehat{H}\psi `$.
Therefore:
$`i\mathrm{}\frac{\psi }{t}=\widehat{H}\psi `$.
Next we get the Heisenberg equations: putting the Schrödinger transform in the form $`\widehat{s}\widehat{\mathrm{\Lambda }}\widehat{s^{}}=\widehat{L}`$ and performing the derivatives with respect to time one gets Heisenberg’s equation
$`\frac{\widehat{L}}{t}=\frac{\widehat{s}}{t}\widehat{\mathrm{\Lambda }}\widehat{s^{}}+\widehat{s}\widehat{\mathrm{\Lambda }}\frac{\widehat{s^{}}}{t}=\frac{i}{\mathrm{}}\widehat{H}\mathrm{exp}^{\frac{i\widehat{H}t}{\mathrm{}}}\widehat{\mathrm{\Lambda }}\widehat{s^{}}\frac{i}{\mathrm{}}\widehat{s}\widehat{\lambda }\mathrm{exp}^{\frac{i\widehat{H}t}{\mathrm{}}}\widehat{H}`$
$`=\frac{i}{\mathrm{}}(\widehat{H}\widehat{s}\widehat{\mathrm{\Lambda }}\widehat{s^{}}\widehat{s}\widehat{\mathrm{\Lambda }}\widehat{s^{}}\widehat{H})=\frac{i}{\mathrm{}}(\widehat{H}\widehat{L}\widehat{L}\widehat{H})=\frac{i}{\mathrm{}}[\widehat{H},\widehat{L}]`$.
Thus, we have:
$`\frac{\widehat{L}}{t}=\frac{i}{\mathrm{}}[\widehat{H},\widehat{L}]`$.
Moreover, Heisenberg’s equation can be written in the form:
$`\frac{\widehat{L}}{t}=\frac{i}{\mathrm{}}\widehat{s}[\widehat{H},\widehat{\mathrm{\Lambda }}]\widehat{s^{}}`$.
$`\widehat{L}`$ is known as an integral of motion, which, if $`\frac{d}{dt}<\psi \widehat{L}\psi >=0`$, is characterized by the following commutators:
$`[\widehat{H},\widehat{L}]=0`$, $`[\widehat{H},\widehat{\mathrm{\Lambda }}]=0`$.
1. Stationary states
The states of a quantum system described by the eigenfunctions of $`\widehat{H}`$ are called stationary states and the corresponding set of eigenvalues is known as the energy spectrum of the system. In such cases, the Schroedinger equation is:
$`i\mathrm{}\frac{\psi _n}{t}=E_n\psi _n=\widehat{H}\psi _n`$.
The solutions are of the form: $`\psi _n(x,t)=\mathrm{exp}^{\frac{iE_nt}{\mathrm{}}}\varphi _n(x)`$.
* The probability is the following:
$`\delta (x)=\psi _n(x,t)^2=\mathrm{exp}^{\frac{iE_nt}{\mathrm{}}}\varphi _n(x)^2`$
$`=\mathrm{exp}^{\frac{iE_nt}{\mathrm{}}}\mathrm{exp}^{\frac{iE_nt}{\mathrm{}}}\varphi _n(x)^2=\varphi _n(x)^2`$.
Thus, the probability is constant in time.
* In the stationary states, the mean value of any commutator of the form $`[\widehat{H},\widehat{A}]`$ is zero, where $`\widehat{A}`$ is an arbitrary operator:
$`<n\widehat{H}\widehat{A}\widehat{A}\widehat{H}n>=<n\widehat{H}\widehat{A}n><n\widehat{A}\widehat{H}n>`$
$`=<nE_n\widehat{A}n><n\widehat{A}E_nn>`$
$`=E_n<n\widehat{A}n>E_n<n\widehat{A}n>=0`$.
* The virial theorem in quantum mechanics \- if $`\widehat{H}`$ is a hamiltonian operator of a particle in the field $`U(r)`$, using
$`\widehat{A}=1/2_{i=1}^3(\widehat{p_i}\widehat{x_i}\widehat{x_i}\widehat{p_i})`$ one gets:
$`<\psi [\widehat{A},\widehat{H}]\psi >=0=<\psi \widehat{A}\widehat{H}\widehat{H}\widehat{A}\psi >`$
$`=_{i=1}^3<\psi \widehat{p_i}\widehat{x_i}\widehat{H}\widehat{H}\widehat{p_i}\widehat{x_i}\psi >`$
$`=_{i=1}^3<\psi [\widehat{H},\widehat{x_i}]\widehat{p_i}+\widehat{x_i}[\widehat{H},\widehat{p_i}]\psi >`$.
Using several times the commutators and $`\widehat{p_i}=i\mathrm{}_i`$, $`\widehat{H}=\widehat{T}+U(r)`$, one can get:
$`<\psi [\widehat{A},\widehat{H}]\psi >=0`$
$`=i\mathrm{}(2<\psi \widehat{T}\psi ><\psi \stackrel{}{r}U(r)\psi >)`$.
This is the virial theorem. If the potential is $`U(r)=U_or^n`$, then a form of the virial theorem similar to that in classical mechanics can be obtained with the only difference that it refers to mean values
$`\overline{T}=\frac{n}{2}\overline{U}`$.
* For a Hamiltonian $`\widehat{H}=\frac{\mathrm{}^2}{2m}^2+U(r)`$ and $`[\stackrel{}{r},H]=\frac{i\mathrm{}}{m}\stackrel{}{p}`$, calculating the matrix elements one finds:
$`(E_kE_n)<n\stackrel{}{r}k>=\frac{i\mathrm{}}{m}<n\widehat{p}k>`$.
1. The nonrelativistic probability current density
The following integral:
$`\psi _n(x)^2𝑑x=1`$,
is the normalization of an eigenfunction of the discrete spectrum in the coordinate representation. It appears as a condition on the microscopic motion in a finite region of space.
For the wavefunctions of the continuous spectrum $`\psi _\lambda (x)`$ one cannot give a direct probabilistic interpretation.
Let us consider a given wavefunction $`\varphi `$ $``$ $`^2`$, that we write as a linear combination of eigenfunctions of the continuum:
$`\varphi =a(\lambda )\psi _\lambda (x)𝑑x.`$
One says that $`\varphi `$ corresponds to an infinite motion.
In many cases, the function $`a(\lambda )`$ is not zero only in a small neighborhood of a point $`\lambda =\lambda _o`$. In such a case, $`\varphi `$ is known as a wavepacket.
We shall calculate now the rate of change of the probability of finding the system in the volume $`\mathrm{\Omega }`$.
$`P=_\mathrm{\Omega }\psi (x,t)^2𝑑x=_\mathrm{\Omega }\psi ^{}(x,t)\psi (x,t)𝑑x`$.
Derivating the integral with respect to time leads to
$`\frac{dP}{dt}=_\mathrm{\Omega }(\psi \frac{\psi ^{}}{t}+\psi ^{}\frac{\psi }{t})𝑑x`$.
Using now the Schrödinger equation in the integral of the right hand side, one gets:
$`\frac{dP}{dt}=\frac{i}{\mathrm{}}_\mathrm{\Omega }(\psi \widehat{H}\psi ^{}\psi ^{}\widehat{H}\psi )𝑑x`$.
Using the identity $`f^2gg^2f=div[(f)grad(g)(g)grad(f)]`$ and also the Schrödinger equation in the form:
$`\widehat{H}\psi =\frac{\mathrm{}^2}{2m}^2\psi `$
and subtituting in the integral, one gets:
$`\frac{dP}{dt}=\frac{i}{\mathrm{}}_\mathrm{\Omega }[\psi (\frac{\mathrm{}^2}{2m}\psi ^{})\psi ^{}(\frac{\mathrm{}^2}{2m}\psi )]𝑑x`$
$`=_\mathrm{\Omega }\frac{i\mathrm{}}{2m}(\psi \psi ^{}\psi ^{}\psi )𝑑x`$
$`=_\mathrm{\Omega }𝑑iv\frac{i\mathrm{}}{2m}(\psi \psi ^{}\psi ^{}\psi )𝑑x`$.
By means of the divergence theorem, the volume integral can be transformed in a surface one leading to:
$`\frac{dP}{dt}=\frac{i\mathrm{}}{2m}(\psi \psi ^{}\psi ^{}\psi )𝑑x`$.
The quantity $`\stackrel{}{J}(\psi )=\frac{i\mathrm{}}{2m}(\psi \psi ^{}\psi ^{}\psi )`$ is known as the probability density current, for which one can easily get the following continuity equation
$`\frac{d\rho }{dt}+div(\stackrel{}{J})=0`$.
* If $`\psi (x)=AR(x)`$, where $`R(x)`$ is a real function, then: $`\stackrel{}{J}(\psi )=0`$.
* For momentum eigenfunctions $`\psi (x)=\frac{1}{(2\pi \mathrm{})^3/2}\mathrm{exp}^{\frac{i\stackrel{}{p}\stackrel{}{x}}{\mathrm{}}}`$, one gets:
$`J(\psi )=\frac{i\mathrm{}}{2m}(\frac{1}{(2\pi \mathrm{})^3/2}\mathrm{exp}^{\frac{i\stackrel{}{p}\stackrel{}{x}}{\mathrm{}}}(\frac{i\stackrel{}{p}}{\mathrm{}(2\pi \mathrm{})^3/2}\mathrm{exp}^{\frac{i\stackrel{}{p}\stackrel{}{x}}{\mathrm{}}})`$
$`(\frac{1}{(2\pi \mathrm{})^3/2}\mathrm{exp}^{\frac{i\stackrel{}{p}\stackrel{}{x}}{\mathrm{}}}\frac{i\stackrel{}{p}}{\mathrm{}(2\pi \mathrm{})^3/2}\mathrm{exp}^{\frac{i\mathrm{}\stackrel{}{p}\stackrel{}{x}}{\mathrm{}}}))`$
$`=\frac{i\mathrm{}}{2m}(\frac{2i\stackrel{}{p}}{\mathrm{}(2\pi \mathrm{})^3})=\frac{\stackrel{}{p}}{m(2\pi \mathrm{})^3}`$,
which shows that the probability density current does not depend on the coordinate.
1. Operator of spatial transport
If $`\widehat{H}`$ is invariant at translations of arbitrary vector $`\stackrel{}{a}`$,
$`\widehat{H}(\stackrel{}{r}+\stackrel{}{a})=\widehat{H}\stackrel{}{(r)}`$ ,
then there is an operator $`\widehat{T}(\stackrel{}{a})`$ which is unitary $`\widehat{T}^{}(\stackrel{}{a})\widehat{H}(\stackrel{}{r})\widehat{T}(\stackrel{}{a})=\widehat{H}(\stackrel{}{r}+\stackrel{}{a})`$.
Commutativity of translations
$`\widehat{T}(\stackrel{}{a})\widehat{T}(\stackrel{}{b})=\widehat{T}(\stackrel{}{b})\widehat{T}(\stackrel{}{a})=\widehat{T}(\stackrel{}{a}+\stackrel{}{b})`$,
implies that $`\widehat{T}`$ is of the form $`\widehat{T}=\mathrm{exp}^{i\widehat{k}a}`$, where $`\widehat{k}=\frac{\widehat{p}}{\mathrm{}}`$.
In the infinitesimal case:
$`\widehat{T}(\delta \stackrel{}{a})\widehat{H}\widehat{T}(\delta \stackrel{}{a})(\widehat{I}+i\widehat{k}\delta \stackrel{}{a})\widehat{H}(\widehat{I}i\widehat{k}\delta \stackrel{}{a})`$,
$`\widehat{H}(\stackrel{}{r})+i[\widehat{K},\widehat{H}]\delta \stackrel{}{a}=\widehat{H}(\stackrel{}{r})+(\widehat{H})\delta \stackrel{}{a}`$.
Moreover, $`[\widehat{p},\widehat{H}]=0`$, where $`\widehat{p}`$ is an integral of the motion. The sistem of wavefunctions of the form $`\psi (\stackrel{}{p},\stackrel{}{r})=\frac{1}{(2\pi \mathrm{})^3/2}\mathrm{exp}^{\frac{i\stackrel{}{p}\stackrel{}{r}}{\mathrm{}}}`$ and the unitary transformation leads to $`\mathrm{exp}^{\frac{i\stackrel{}{p}\stackrel{}{a}}{\mathrm{}}}\psi (\stackrel{}{r})=\psi (\stackrel{}{r}+\stackrel{}{a})`$. The operator of spatial transport $`\widehat{T}^{}=\mathrm{exp}^{\frac{i\stackrel{}{p}\stackrel{}{a}}{\mathrm{}}}`$ is the analog of $`\widehat{s}^{}=\mathrm{exp}^{\frac{i\widehat{H}t}{\mathrm{}}}`$, which is the operator of time ‘transport’ (shift).
1. Exemple: The ‘crystal’ (lattice) Hamiltonian
If $`\widehat{H}`$ is invariant for a discrete translation (for exemple, in a crystal lattice) $`\widehat{H}(\stackrel{}{r}+\stackrel{}{a})=\widehat{H}(\stackrel{}{r})`$, where $`\stackrel{}{a}=_i\stackrel{}{a_i}n_i`$, $`n_i`$ $``$ $`N`$ and $`a_i`$ are baricentric vectors, then:
$`\widehat{H}(\stackrel{}{r})\psi (\stackrel{}{r})=E\psi (\stackrel{}{r})`$,
$`\widehat{H}(\stackrel{}{r}+\stackrel{}{a})\psi (\stackrel{}{r}+\stackrel{}{a})=E\psi (\stackrel{}{r}+\stackrel{}{a})=\widehat{H}(\stackrel{}{r})\psi (\stackrel{}{r}+\stackrel{}{a})`$.
Consequently, $`\psi (\stackrel{}{r})`$ and $`\psi (\stackrel{}{r}+\stackrel{}{a})`$ are wavefunctions for the same eigenvalue of $`\widehat{H}`$. The relationship between $`\psi (\stackrel{}{r})`$ and $`\psi (\stackrel{}{r}+\stackrel{}{a})`$ can be saught for in the form $`\psi (\stackrel{}{r}+\stackrel{}{a})=\widehat{c}(\stackrel{}{a})\psi (\stackrel{}{r})`$, where $`\widehat{c}(\stackrel{}{a})`$ is a gxg matrix (g is the order of degeneration of level E). Two column matrices, $`\widehat{c}(\stackrel{}{a})`$ and $`\widehat{c}(\stackrel{}{b})`$ commute and therefore they are diagonalizable simultaneously.
Moreover, for the diagonal elements, $`c_{ii}(\stackrel{}{a})c_{ii}(\stackrel{}{b})=c_{ii}(\stackrel{}{a}+\stackrel{}{b})`$ holds for i=1,2,….,g, having solutions of the type $`c_{ii}(a)=\mathrm{exp}^{ik_ia}`$. Thus, $`\psi _k(\stackrel{}{r})=U_k(\stackrel{}{r})\mathrm{exp}^{i\stackrel{}{k}\stackrel{}{a}}`$, where $`\stackrel{}{k}`$ is a real arbitrary vector and the function $`U_k(\stackrel{}{r})`$ is periodic of period $`\stackrel{}{a}`$, $`U_k(\stackrel{}{r}+\stackrel{}{a})=U_k(\stackrel{}{r})`$.
The assertion that the eigenfunctions of a periodic $`\widehat{H}`$ of the lattice type $`\widehat{H}(\stackrel{}{r}+\stackrel{}{a})=\widehat{H}(\stackrel{}{r})`$ can be written $`\psi _k(\stackrel{}{r})=U_k(\stackrel{}{r})\mathrm{exp}i\stackrel{}{k}\stackrel{}{a}`$, where $`U_k(\stackrel{}{r}+\stackrel{}{a})=U_k(\stackrel{}{r})`$ is known as Bloch’s theorem. In the continuous case, $`U_k`$ should be constant, because the constant is the only function periodic for any $`\stackrel{}{a}`$. The vector $`\stackrel{}{p}=\mathrm{}\stackrel{}{k}`$ is called quasimomentum (by analogy with the continuous case). The vector $`\stackrel{}{k}`$ is not determined univoquely, because one can add any vector $`\stackrel{}{g}`$ for which $`ga=2\pi n`$, where $`n`$ $``$ $`N`$.
The vector $`\stackrel{}{g}`$ can be written $`\stackrel{}{g}=_{i=1}^3\stackrel{}{b_i}m_i`$, where $`m_i`$ are integers and $`b_i`$ are given by
$`\stackrel{}{b_i}=2\pi \frac{\widehat{a_j}\times \stackrel{}{a_k}}{\stackrel{}{a_i}(\stackrel{}{a_j}\times \stackrel{}{a_k})}`$,
for $`ijk`$. $`\stackrel{}{b_i}`$ are the baricentric vectors of the lattice.
Recommended references
1. E. Farhi, J. Goldstone, S. Gutmann, “How probability arises in quantum mechanics”, Annals of Physics 192, 368-382 (1989)
2. N.K. Tyagi in Am. J. Phys. 31, 624 (1963) gives a very short proof of the Heisenberg uncertainty principle, which asserts that the simultaneous measurement of two noncommuting hermitic operators results in an uncertainty given by the value of their commutator.
3. H.N. Núñez-Yépez et al., “Simple quantum systems in the momentum representation”, physics/0001030 (Europ. J. Phys., 2000).
4. J.C. Garrison, “Quantum mechanics of periodic systems”, Am. J. Phys. 67, 196 (1999).
5. F. Gieres, “Dirac’s formalism and mathematical surprises in quantum mechanics”, quant-ph/9907069 (in English); quant-ph/9907070 (in French).
1N. Notes
1. For “the creation of quantum mechanics…”, Werner Heisenberg has been awarded the Nobel prize in 1932 (delivered in 1933). The paper “Zür Quantenmechanik. II”, \[“On quantum mechanics.II”, Zf. f. Physik 35, 557-615 (1926) (received by the Editor on 16 November 1925) by M. Born, W. Heisenberg and P. Jordan, is known as the “work of the three people”, being considered as the work that really opened the vast horizons of quantum mechanics.
2. For “the statistical interpretation of the wavefunction” Max Born was awarded the Nobel prize in 1954.
## 1P. Problems
Problema 1.1: Let us consider two operators, A and B, which commutes by hypothesis. In this case, one can derive the following relationship:
$`e^Ae^B=e^{(A+B)}e^{(1/2[A,B])}`$.
Solution
Defining an operator F(t), as a function of real variable t, of the form: $`F(t)=e^{(At)}e^{(Bt)}`$,
then: $`\frac{dF}{dt}=Ae^{At}e^{Bt}+e^{At}Be^{Bt}=(A+e^{At}Be^{At})F(t)`$.
Applying now the formula $`[A,F(B)]=[A,B]F^{^{}}(B)`$, we have
$`[e^{At},B]=t[A.B]e^{At}`$, and therefore: $`e^{At}B=Be^{At}+t[A,B]e^{At}.`$
Multiplying both sides of the latter equation by $`\mathrm{exp}^{At}`$ and substituting in the first equation, we get:
$`\frac{dF}{dt}=(A+B+t[A,B])F(t)`$.
The operators A , B and \[A,B\] commutes by hypothesis. Thus, we can integrate the differential equation as if $`A+B`$ and $`[A,B]`$ would be scalar numbers.
We shall have:
$`F(t)=F(0)e^{(A+B)t+1/2[A,B]t^2}`$.
Putting $`t=0`$, one can see that $`F(0)=1`$ and therefore :
$`F(t)=e^{(A+B)t+1/2[A,B]t^2}`$.
Putting now $`t=1`$, we get the final result.
Problem 1.2: Calculate the commutator $`[X,D_x]`$.
Solution
The calculation is performed by applying the commutator to an arbitrary function $`\psi (\stackrel{}{r})`$:
$`[X,D_x]\psi (\stackrel{}{r})=(x\frac{}{x}\frac{}{x}x)\psi (\stackrel{}{r})=x\frac{}{x}\psi (\stackrel{}{r})\frac{}{x}[x\psi (\stackrel{}{r})]=x\frac{}{x}\psi (\stackrel{}{r})\psi (\stackrel{}{r})x\frac{}{x}\psi (\stackrel{}{r})=\psi (\stackrel{}{r})`$.
Since this relationship is satisfied for any $`\psi (\stackrel{}{r})`$, one can conclude that $`[X,D_x]=1`$.
Problem 1.3: Check that the trace of a matrix is invariant of changes of discrete orthonormalized bases.
Solution
The sum of the diagonal elements of a matrix representation of an operator A in an arbitrary basis does not depend on the basis.
This important property can be obtained by passing from an orthonormalized discrete basis $`u_i>`$ to another orthonormalized discrete basis $`t_k>`$. We have:
$`_i<u_iAu_i>=_i<u_i\left(_kt_k><t_k\right)Au_i>`$
(where we have used the completeness relationship for the states $`t_k`$). The right hand side is:
$`_{i,j}<u_it_k><t_kAu_i>=_{i,j}<t_kAu_i><u_it_k>`$,
(the change of the order in the product of two scalar numbers is allowed). Thus, we can replace $`_iu_i><u_i`$ with unity (i.e., the completeness relationship for the states $`u_i>`$), in order to get finally:
$$\underset{i}{}<u_iAu_i>=\underset{k}{}<t_kAt_k>.$$
Thus, we have proved the invariance property for matriceal traces.
Problem 1.4: If for the hermitic operator $`N`$ there are the hermitic operators $`L`$ and $`M`$ such that : $`[M,N]=0`$, $`[L,N]=0`$, $`[M,L]0`$, then the eigenfunctions of $`N`$ are degenerate.
Solution
Let $`\psi (x;\mu ,\nu )`$ be the common eigenfunctions of $`M`$ and $`N`$ (since they commute they are simultaneous observables). Let $`\psi (x;\lambda ,\nu )`$ be the common eigenfunctions of $`L`$ and $`N`$ (again, since they commute they are simultaneous observables). The Greek parameters denote the eigenvalues of the corresponding operators. Let us consider for simplicity sake that $`N`$ has a discrete spectrum. Then:
$$f(x)=\underset{\nu }{}a_\nu \psi (x;\mu ,\nu )=\underset{\nu }{}b_\nu \psi (x;\lambda ,\nu ).$$
We calculate now the matrix element $`<f|ML|f>`$:
$$<f|ML|f>=\underset{\nu }{}\mu _\nu a_\nu \psi ^{}(x;\mu ,\nu )\underset{\nu ^{^{}}}{}\lambda _\nu ^{^{}}b_\nu ^{^{}}\psi (x;\lambda ,\nu ^{^{}})dx.$$
If all the eigenfunctions of $`N`$ are nondegenerate then $`<f|ML|f>=_\nu \mu _\nu a_\nu \lambda _\nu b_\nu `$. But the same result can be obtained if one calculates $`<f|LM|f>`$ and the commutator would be zero. Thus, at least some of the eigenfunctions of $`N`$ should be degenerate.
## 2. ONE DIMENSIONAL RECTANGULAR BARRIERS AND WELLS
### Regions of constant potential
In the case of a rectangular potential, $`V(x)`$ is a constant function $`V(x)=V`$ in a certain region of the one-dimensional space. In such a region, the Schrödinger eq. can be written:
$$\frac{d^2}{dx^2}\psi (x)+\frac{2m}{\mathrm{}^2}(EV)\psi (x)=0$$
(1)
One can distinguish several cases:
(i) $`E>V`$
Let us introduce the positive constant $`k`$, defined by
$$k=\frac{\sqrt{2m(EV)}}{\mathrm{}}$$
(2)
Then, the solution of eq. (1) can be written:
$$\psi (x)=Ae^{ikx}+A^{}e^{ikx}$$
(3)
where $`A`$ and $`A^{}`$ are complex constants.
(ii) $`E<V`$
This condition corresponds to segments of the real axis which would be prohibited to any particle from the viewpoint of classical mechanics. In this case, one introduces the positive constant $`q`$ defined by:
$$q=\frac{\sqrt{2m(VE)}}{\mathrm{}}$$
(4)
and the solution of (1) can be written:
$$\psi (x)=Be^{qx}+B^{}e^{qx},$$
(5)
where $`B`$ and $`B^{}`$ are complex constants.
(iii) $`E=V`$
In this special case, $`\psi (x)`$ is a linear function of $`x`$.
### The behaviour of $`\psi (x)`$ at a discontinuity of the potential
One might think that at the point $`x=x_1`$, where the potential $`V(x)`$ is discontinuous, the wavefunction $`\psi (x)`$ behaves in a more strange way, maybe discontinuously for example. This is not so: $`\psi (x)`$ and $`\frac{d\psi }{dx}`$ are continuous, and only the second derivative is discontinuous at $`x=x_1`$.
### General look to the calculations
The procedure to determine the stationary states in rectangular potentials is the following: in all regions in which $`V(x)`$ is constant we write $`\psi (x)`$ in any of the two forms (3) or (5) depending on application; next, we join smoothly these functions according to the continuity conditions for $`\psi (x)`$ and $`\frac{d\psi }{dx}`$ at the points where $`V(x)`$ is discontinuous.
## Examination of several simple cases
Let us make explicite calculations for some simple stationary states according to the proposed method.
### The step potential
a. $`E>V_0`$ case; partial reflexion
Let us put eq. (2) in the form:
$`k_1`$ $`=`$ $`{\displaystyle \frac{\sqrt{2mE}}{\mathrm{}}}`$ (6)
$`k_2`$ $`=`$ $`{\displaystyle \frac{\sqrt{2m(EV_0)}}{\mathrm{}}}`$ (7)
The solution of eq. (1) has the form of eq. (3) in the regions $`I(x<0)`$ and $`II(x>0)`$:
$`\psi _I`$ $`=`$ $`A_1e^{ik_1x}+A_1^{}e^{ik_1x}`$
$`\psi _{II}`$ $`=`$ $`A_2e^{ik_2x}+A_2^{}e^{ik_2x}`$
In region I eq. (1) takes the form
$`\psi ^{\prime \prime }(x)+{\displaystyle \frac{2mE}{\mathrm{}^2}}\psi (x)=\psi ^{\prime \prime }(x)+k^2\psi (x)=0`$
and in the region II:
$`\psi ^{\prime \prime }(x){\displaystyle \frac{2m}{\mathrm{}^2}}[V_0E]\varphi (x)=\psi ^{\prime \prime }(x)q^2\psi (x)=0`$
If we limit ourselves to the case of an incident particle ‘coming’ from $`x=\mathrm{}`$, we have to choose $`A_2^{}=0`$ and we can determine the ratios $`A_1^{}/A_1`$ and $`A_2/A_1`$. The joining conditions give then:
* $`\psi _I=\psi _{II}`$, at $`x=0:`$
$$A_1+A_1^{}=A_2$$
(8)
* $`\psi _I^{}=\psi _{II}^{}`$, at $`x=0:`$
$$A_1ik_1A_1^{}ik_1=A_2ik_2$$
(9)
Substituting $`A_1`$ and $`A_1^{}`$ from (8) in (9):
$`A_1^{}`$ $`=`$ $`{\displaystyle \frac{A_2(k_1k_2)}{2k_1}}`$ (10)
$`A_1`$ $`=`$ $`{\displaystyle \frac{A_2(k_1+k_2)}{2k_1}}`$ (11)
From the two expressions of the constant $`A_2`$ in (10) and (11) one gets
$$\frac{A_1^{}}{A_1}=\frac{k_1k_2}{k_1+k_2}$$
(12)
and from (11) it follows:
$$\frac{A_2}{A_1}=\frac{2k_1}{k_1+k_2}.$$
(13)
$`\psi (x)`$ is a superposition of two waves. The first (the $`A_1`$ part) corresponds to an incident wave of momentum $`p=\mathrm{}k_1`$, propagating from the left to the right. The second (the $`A_1^{}`$ part) corresponds to a reflected particle of momentum $`\mathrm{}k_1`$ propagating in opposite direction. Since we have already chosen $`A_2^{}=0`$, it follows that $`\psi _{II}(x)`$ contains a single wave, which is associated to a transmitted particle. (We will show later how it is possible by employing the concept of probability current to define the transmission coefficient T as well as the reflection coefficient R for the step potential). These coefficients give the probability that a particle coming from $`x=\mathrm{}`$ can pass through or get back from the step at $`x=0`$. Thus, we obtain:
$$R=|\frac{A_1^{}}{A_1}|^2,$$
(14)
whereas for $`T`$:
$$T=\frac{k_2}{k_1}|\frac{A_2}{A_1}|^2.$$
(15)
Taking into account (12) and (13) one is led to:
$`R`$ $`=`$ $`1{\displaystyle \frac{4k_1k_2}{(k_1+k_2)^2}}`$ (16)
$`T`$ $`=`$ $`{\displaystyle \frac{4k_1k_2}{(k_1+k_2)^2}}.`$ (17)
It is easy to check that $`R+T=1`$. It is thus sure that the particle will be either transmitted or reflected. Contrary to the predictions of classical mechanics, the incident particle has a nonzero probability of not going back.
It is also easy to check using (6), (7) and (17), that if $`EV_0`$ then $`T1`$: when the energy of the particle is sufficently big in comparison with the height of the step, everything happens as if the step does not exist for the particle.
Consider the following natural form of the solution in region I:
$`\psi _I=A_1e^{ik_1x}+Ae^{ik_1x}`$
$$j=\frac{i\mathrm{}}{2m}(\varphi ^{}\varphi \varphi \varphi ^{})$$
(18)
with $`A_1e^{ik_1x}`$ and its conjugate $`A_1^{}e^{ik_1x}`$:
$`j`$ $`=`$ $`{\displaystyle \frac{i\mathrm{}}{2m}}[(A_1^{}e^{ik_1x})(A_1ik_1e^{ik_1x})(A_1e^{ik_1x})(A_1^{}ik_1e^{ik_1x})]`$
$`j`$ $`=`$ $`{\displaystyle \frac{\mathrm{}k_1}{m}}|A_1|^2.`$
Now with $`Ae^{ik_1x}`$ and its conjugate $`A^{}e^{ik_1x}`$ one is led to:
$`j={\displaystyle \frac{\mathrm{}k_1}{m}}|A|^2.`$
In the following we wish to check the proportion of reflected current with respect to the incident current (or more exactly, we want to check the relative probability that the particle is returned back):
$`R={\displaystyle \frac{|j(\varphi _{})|}{|j(\varphi _+)|}}={\displaystyle \frac{|\frac{\mathrm{}k_1}{m}|A|^2|}{|\frac{\mathrm{}k_1}{m}|A_1|^2|}}=|{\displaystyle \frac{A}{A_1}}|^2.`$ (19)
Similarly, the proportion of transmission with respect to incidence (that is the probability that the particle is transmitted) is, taking now into account the solution in the region II:
$`T={\displaystyle \frac{|\frac{\mathrm{}k_2}{m}|A_2|^2|}{|\frac{\mathrm{}k_1}{m}|A_1|^2|}}={\displaystyle \frac{k_2}{k_1}}|{\displaystyle \frac{A_2}{A_1}}|^2.`$ (20)
b. $`E<V_0`$ case; total reflection
In this case we have:
$`k_1`$ $`=`$ $`{\displaystyle \frac{\sqrt{2mE}}{\mathrm{}}}`$ (21)
$`q_2`$ $`=`$ $`{\displaystyle \frac{\sqrt{2m(V_0E)}}{\mathrm{}}}`$ (22)
In the region $`I(x<0)`$, the solution of eq. (1) \[written as $`\psi (x)^{\prime \prime }+k_1^2\psi (x)=0`$\] has the form given in eq. (3):
$$\psi _I=A_1e^{ik_1x}+A_1^{}e^{ik_1x},$$
(23)
whereas in the region $`II(x>0)`$, the same eq. (1) \[now written as $`\psi (x)^{\prime \prime }q_2^2\psi (x)=0`$\] has the form of eq. (5):
$$\psi _{II}=B_2e^{q_2x}+B_2^{}e^{q_2x}.$$
(24)
In order that the solution be kept finite when $`x+\mathrm{}`$, it is necessary that:
$$B_2=0.$$
(25)
The joining condition at $`x=0`$ give now:
* $`\psi _I=\psi _{II}`$, at $`x=0:`$
$$A_1+A_1^{}=B_2^{}$$
(26)
* $`\psi _I^{}=\psi _{II}^{}`$, at $`x=0:`$
$$A_1ik_1A_1^{}ik_1=B_2^{}q_2.$$
(27)
Substituting $`A_1`$ and $`A_1^{}`$ from (26) in (27) we get:
$`A_1^{}`$ $`=`$ $`{\displaystyle \frac{B_2^{}(ik_1+q_2)}{2ik_1}}`$ (28)
$`A_1`$ $`=`$ $`{\displaystyle \frac{B_2^{}(ik_1q_2)}{2ik1}}.`$ (29)
Equating the expressions for the constant $`B_2^{}`$ from (28) and (29) leads to:
$$\frac{A_1^{}}{A_1}=\frac{ik_1+q_2}{ik_1q_2}=\frac{k_1iq_2}{k_1+iq_2},$$
(30)
so that from (29) we have:
$$\frac{B_2^{}}{A_1}=\frac{2ik_1}{ik_1q_2}=\frac{2k_1}{k_1iq_2}.$$
(31)
Therefore, the reflection coefficient $`R`$ is:
$$R=|\frac{A_1^{}}{A_1}|^2=|\frac{k_1iq_2}{k_1+iq_2}|^2=\frac{k_1^2+q_2^2}{k_1^2+q_2^2}=1.$$
(32)
As in classical mechanics, the microparticle is always reflected (total reflexion). However, there is an important difference, namely, because of the existence of the so-called evanescent wave $`e^{q_2x}`$, the particle has a nonzero probability to find itself in a spatial region which is classicaly forbidden. This probability decays exponentially with $`x`$ and turns to be negligible when $`x`$ overcome $`1/q_2`$ corresponding to the evanescent wave. Notice also that $`A_1^{}/A_1`$ is a complex quantity. A phase difference occurs as a consequence of the reflexion, which physically is due to the fact that the particle is slowed down when entering the region $`x>0`$. There is no analog phenomenon for this in classical mechanics (but there is of course such an analog in optical physics).
### Rectangular barrier
a. $`E>V_0`$ case; resonances
Here we put eq. (2) in the form:
$`k_1`$ $`=`$ $`{\displaystyle \frac{\sqrt{2mE}}{\mathrm{}}}`$ (33)
$`k_2`$ $`=`$ $`{\displaystyle \frac{\sqrt{2m(EV_0)}}{\mathrm{}}}.`$ (34)
The solution of eq. (1) is as in eq. (3) in the regions $`I(x<0)`$, $`II(0<x<a`$) and $`III(x>a):`$
$`\psi _I`$ $`=`$ $`A_1e^{ik_1x}+A_1^{}e^{ik_1x}`$
$`\psi _{II}`$ $`=`$ $`A_2e^{ik_2x}+A_2^{}e^{ik_2x}`$
$`\psi _{III}`$ $`=`$ $`A_3e^{ik_1x}+A_3^{}e^{ik_1x}.`$
If we limit ourselves to the case of an incident particle coming from $`x=\mathrm{}`$, we have to choose $`A_3^{}=0`$.
* $`\psi _I=\psi _{II}`$, at $`x=0:`$
$$A_1+A_1^{}=A_2+A_2^{}$$
(35)
* $`\psi _I^{}=\psi _{II}^{}`$, at $`x=0:`$
$$A_1ik_1A_1^{}ik_1=A_2ik_2A_2^{}ik_2$$
(36)
* $`\psi _{II}=\psi _{III}`$, at $`x=a:`$
$$A_2e^{ik_2a}+A_2^{}e^{ik_2a}=A_3e^{ik_1a}$$
(37)
* $`\psi _{II}^{}=\psi _{III}^{}`$, at $`x=a:`$
$$A_2ik_2e^{ik_2a}A_2^{}ik_2e^{ik_2a}=A_3ik_1e^{ik_1a}.$$
(38)
The joining conditions at $`x=a`$ give $`A_2`$ and $`A_2^{}`$ as functions of $`A_3`$, whereas those at $`x=0`$ give $`A_1`$ and $`A_1^{}`$ as functions of $`A_2`$ and $`A_2^{}`$ (thus, as functions of $`A_3`$). This procedure is shown in detail in the following.
Substituting $`A_2^{}`$ from eq. (37) in (38) leads to:
$$A_2=\frac{A_3e^{ik_1a}(k_2+k_1)}{2k_2e^{ik_2a}}.$$
(39)
Substituting $`A_2`$ from eq. (37) in (38) leads to:
$$A_2^{}=\frac{A_3e^{ik_1a}(k_2k_1)}{2k_2e^{ik_2a}}.$$
(40)
Substituting $`A_1`$ from eq. (35) in (36) leads to:
$$A_1^{}=\frac{A_2(k_2k_1)A_2^{}(k_2+k_1)}{2k_1}.$$
(41)
Substituting $`A_1^{}`$ from eq. (35) in (36) gives:
$$A_1=\frac{A_2(k_2+k_1)A_2^{}(k_2k_1)}{2k_1}.$$
(42)
Now, substituting the eqs. (39) and (40) in (41), we have:
$$A_1^{}=i\frac{(k_2^2k_1^2)}{2k_1k_2}(\mathrm{sin}k_2a)e^{ik_1a}A_3.$$
(43)
Finally, substituting the eqs. (39) and (40) in (42) we get:
$$A_1=[\mathrm{cos}k_2ai\frac{k_1^2+k_2^2}{2k_1k_2}\mathrm{sin}k_2a]e^{ik_1a}A_3.$$
(44)
$`A_1^{}/A_1`$ and $`A_3/A_1`$ \[these ratios can be obtained by equating (43) and (44), and by separating, respectively, in eq. (44)\] allow the calculation of the reflexion coefficient $`R`$ as well as of the transmission one $`T`$. For this type of barrier, they are given by the following formulas:
$$R=|A_1^{}/A_1|^2=\frac{(k_1^2k_2^2)^2\mathrm{sin}^2k_2a}{4k_1^2k_2^2+(k_1^2k_2^2)^2\mathrm{sin}^2k_2a},$$
(45)
$$T=|A_3/A_1|^2=\frac{4k_1^2k_2^2}{4k_1^2k_2^2+(k_1^2k_2^2)^2\mathrm{sin}^2k_2a}.$$
(46)
It is easy to see that they check $`R+T=1`$.
b. $`E<V_0`$ case; the tunnel effect
Now, let us take the eqs. (2) and (4):
$`k_1`$ $`=`$ $`{\displaystyle \frac{\sqrt{2mE}}{\mathrm{}}}`$ (47)
$`q_2`$ $`=`$ $`{\displaystyle \frac{\sqrt{2m(V_0E)}}{\mathrm{}}}.`$ (48)
The solution of eq. (1) has the form given in eq. (3) in the regions $`I(x<0)`$ and $`III(x>a)`$, while in the region $`II(0<x<a`$) has the form of eq. (5):
$`\psi _I`$ $`=`$ $`A_1e^{ik_1x}+A_1^{}e^{ik_1x}`$
$`\psi _{II}`$ $`=`$ $`B_2e^{q_2x}+B_2^{}e^{q_2x}`$
$`\psi _{III}`$ $`=`$ $`A_3e^{ik_1x}+A_3^{}e^{ik_1x}.`$
The joining conditions at $`x=0`$ and $`x=a`$ allow the calculation of the transmission coefficient of the barrier. As a matter of fact, it is not necessary to repeat the calculation: merely, it is sufficient to replace $`k_2`$ by $`iq_2`$ in the equation obtained in the first case of this section.
### Bound states in rectangular well
a. Well of finite depth
We first study the case $`0<E<V_0`$ ($`E>V_0`$ is similar to the calculation in the previous section).
For the exterior regions I, $`(x<0)`$ and III, $`(x>a)`$ we employ eq. (4):
$$q=\frac{\sqrt{2m(V_0E)}}{\mathrm{}}.$$
(49)
For the central region II $`(0<x<a)`$ we use eq. (2):
$$k=\frac{\sqrt{2m(E)}}{\mathrm{}}.$$
(50)
The solution of eq. (1) has the form of eq. (5) in the exterior regions and of eq. (3) in the central region:
$`\psi _I`$ $`=`$ $`B_1e^{qx}+B_1^{}e^{qx}`$
$`\psi _{II}`$ $`=`$ $`A_2e^{ikx}+A_2^{}e^{ikx}`$
$`\psi _{III}`$ $`=`$ $`B_3e^{qx}+B_3^{}e^{qx}`$
In the region $`(0<x<a)`$ eq. (1) has the form:
$$\psi (x)^{\prime \prime }+\frac{2mE}{\mathrm{}^2}\psi (x)=\psi (x)^{\prime \prime }+k^2\psi (x)=0$$
(51)
while in the exterior regions:
$$\psi (x)^{\prime \prime }\frac{2m}{\mathrm{}^2}[V_0E]\varphi (x)=\psi (x)^{\prime \prime }q^2\psi (x)=0.$$
(52)
Because $`\psi `$ should be finite in the region I, we impose:
$$B_1^{}=0.$$
(53)
The joining conditions give:
$`\psi _I=\psi _{II}`$, at $`x=0:`$
$$B_1=A_2+A_2^{}$$
(54)
$`\psi _I^{}=\psi _{II}^{}`$, at $`x=0:`$
$$B_1q=A_2ikA_2^{}ik$$
(55)
$`\psi _{II}=\psi _{III}`$, at $`x=a:`$
$$A_2e^{ika}+A_2^{}e^{ika}=B_3e^{qa}+B_3^{}e^{qa}$$
(56)
$`\psi _{II}^{}=\psi _{III}^{}`$, at $`x=a:`$
$$A_2ike^{ika}A_2^{}ike^{ika}=B_3qe^{qa}B_3^{}qe^{qa}$$
(57)
Substituting the constants $`A_2`$ and $`A_2^{}`$ from eq. (54) in eq. (55) we get
$`A_2^{}`$ $`=`$ $`{\displaystyle \frac{B_1(qik)}{2ik}}`$
$`A_2`$ $`=`$ $`{\displaystyle \frac{B_1(q+ik)}{2ik}},`$ (58)
respectively.
Substituting the constant $`A_2`$ and the constant $`A_2^{}`$ from eq. (56) in eq. (57) we get
$`B_3^{}e^{qa}(ik+q)+B_3e^{qa}(ikq)+A_2^{}e^{ika}(2ik)`$ $`=`$ $`0`$
$`2ikA_2e^{ika}+B_3^{}e^{qa}(ik+q)+B_3E^{qa}(ikq)`$ $`=`$ $`0,`$ (59)
respectively.
Equating $`B_3^{}`$ from eqs. (59) and taking into account the eqs (58) leads to
$$\frac{B_3}{B_1}=\frac{e^{qa}}{4ikq}[e^{ika}(q+ik)^2e^{ika}(qik)^2].$$
(60)
Since $`\psi (x)`$ should be finite in region III as well, we require $`B_3=0`$. Thus
$$[\frac{qik}{q+ik}]^2=\frac{e^{ika}}{e^{ika}}=e^{2ika}.$$
(61)
Because $`q`$ and $`k`$ depend on $`E`$, eq. (1) can be satisfied for some particular values of $`E`$. The condition that $`\psi (x)`$ should be finite in all spatial regions imposes the quantization of the energy. Two cases are possible:
(i) if
$$\frac{qik}{q+ik}=e^{ika},$$
(62)
equating in both sides the real and the imaginary parts, respectively, we have
$$\mathrm{tan}(\frac{ka}{2})=\frac{q}{k}.$$
(63)
Putting
$$k_0=\sqrt{\frac{2mV_0}{\mathrm{}}}=\sqrt{k^2+q^2}$$
(64)
one gets
$$\frac{1}{\mathrm{cos}^2(\frac{ka}{2})}=1+\mathrm{tan}^2(\frac{ka}{2})=\frac{k^2+q^2}{k^2}=(\frac{k_0}{k})^2$$
(65)
Eq. (63) is therefore equivalent to the system of eqs.
$`|\mathrm{cos}({\displaystyle \frac{ka}{2}})|`$ $`=`$ $`{\displaystyle \frac{k}{k_0}}`$
$`\mathrm{tan}({\displaystyle \frac{ka}{2}})`$ $`>`$ $`0`$ (66)
The energy levels are determined by the intersection of a straight line of slope $`\frac{1}{k_0}`$ with the first set of dashed cosinusoides in fig. 2.4. Thus, we get a certain number of energy levels whose wavefunctions are even. This fact becomes clearer if we substitute (62) in (58) and (60). It is easy to check that $`B_3^{}=B_1`$ and $`A_2=A_2^{}`$ leading to $`\psi (x)=\psi (x)`$.
(ii) if
$$\frac{qik}{q+ik}=e^{ika},$$
(67)
a similar calculation gives
$`|\mathrm{sin}({\displaystyle \frac{ka}{2}})|`$ $`=`$ $`{\displaystyle \frac{k}{k_0}}`$
$`\mathrm{tan}({\displaystyle \frac{ka}{2}})`$ $`<`$ $`0.`$ (68)
The energy levels are in this case determined by the intersection of the same straight line with the second set of dashed cosinusoides in fig. 2.4. The obtained levels are interlaced with those found in the case (i). One can easily show that the corresponding wavefunctions are odd.
b. Well of infinite depth
In this case it is convenient to put $`V(x)`$ equal to zero for $`0<x<a`$ and equal to infinity for the rest of the real axis. Putting
$$k=\sqrt{\frac{2mE}{\mathrm{}^2}},$$
(69)
$`\psi (x)`$ should be zero outside the interval $`[0,a]`$ and continuous at $`x=0`$ and $`x=a`$.
For $`0xa`$:
$$\psi (x)=Ae^{ikx}+A^{}e^{ikx}.$$
(70)
Since $`\psi (0)=0`$, one can infer that $`A^{}=A`$, leading to:
$$\psi (x)=2iA\mathrm{sin}(kx).$$
(71)
Moreover, $`\psi (a)=0`$ and therefore
$$k=\frac{n\pi }{a},$$
(72)
where $`n`$ is an arbitrary positive integer. If we normalize the function (71), taking into account (72), then we obtain the stationary wavefunctions
$$\psi _n(x)=\sqrt{\frac{2}{a}}\mathrm{sin}(\frac{n\pi x}{a})$$
(73)
with the energies
$$E_n=\frac{n^2\pi ^2\mathrm{}^2}{2ma^2}.$$
(74)
The quantization of the energy levels is extremely simple in this case. The stationary energies are proportional with the natural numbers squared.
## 2P. Problems
### Problem 2.1: The attractive $`\delta `$ potential
Suppose we have a potential of the form:
$`V(x)=V_0\delta (x);V_0>0;x\mathrm{}.`$
The corresponding wavefunction $`\psi (x)`$ is assumed continuous.
a) Obtain the bound states ($`E<0`$), if they exist, localized in this type of potential.
b) Calculate the dispersion of a plane wave falling on the $`\delta `$ potential and obtain the reflexion coefficient
$`R={\displaystyle \frac{|\psi _{refl}|^2}{|\psi _{inc}|^2}}|_{x=0},`$
where $`\psi _{refl}`$, $`\psi _{inc}`$ are the reflected and incoming waves, respectively.
Suggestion: To determine the behavior of $`\psi (x)`$ in x=0, it is better to proceed by integrating the Schrödinger equation in the interval ($`\epsilon ,+\epsilon `$), and then to apply the limit $`\epsilon `$ $``$ $`0`$.
Solution. a) The Schrödinger eq. is:
$$\frac{d^2}{dx^2}\psi (x)+\frac{2m}{\mathrm{}^2}(E+V_0\delta (x))\psi (x)=0.$$
(75)
Far from the origin we have a differential eq. of the form
$$\frac{d^2}{dx^2}\psi (x)=\frac{2mE}{\mathrm{}^2}\psi (x).$$
(76)
Consequently, the wavefunctions are of the form
$$\psi (x)=Ae^{qx}+Be^{qx}\mathrm{for}x>0\mathrm{and}x<0,$$
(77)
where $`q=\sqrt{2mE/\mathrm{}^2}`$ $`\mathrm{}.`$ Since $`|\psi |^2`$ should be $`^2`$ integrable , we cannot accept that a part of it grows exponentially. Moreover, the wavefunction should be continuous at the origin. With these conditions, we have
$`\psi (x)`$ $`=`$ $`Ae^{qx};(x<0),`$
$`\psi (x)`$ $`=`$ $`Ae^{qx};(x>0).`$ (78)
Integrating the Schrödinger eq. between $`\epsilon `$ and $`+\epsilon `$, we get
$$\frac{\mathrm{}^2}{2m}[\psi ^{}(\epsilon )\psi ^{}(\epsilon )]V_0\psi (0)=E_\epsilon ^{+\epsilon }\psi (x)𝑑x2\epsilon E\psi (0)$$
(79)
Introducing now the result (78) and taking into account the limit $`\epsilon 0`$, we have
$$\frac{\mathrm{}^2}{2m}(qAqA)V_0A=0,$$
(80)
or $`E=m(V_0^2/2\mathrm{}^2)`$ \[$`\frac{V_0^2}{4}`$ in units of $`\frac{\mathrm{}^2}{2m}`$\]. Clearly, there is a single discrete energy. The normalization constant is found to be $`A=\sqrt{mV_0/\mathrm{}^2}`$. The wavefunction of the bound state will be $`\psi _o=Ae^{V_0|x|/2}`$, where $`V_0`$ is in $`\frac{\mathrm{}^2}{2m}`$ units. b) Take now the wavefunction of a plane wave
$$\psi (x)=Ae^{ikx},k^2=\frac{2mE}{\mathrm{}^2}.$$
(81)
It moves from the left to the right and is reflected by the potential. If $`B`$ and $`C`$ are the amplitudes of the reflected and transmitted waves, respectively, then we have
$`\psi (x)`$ $`=`$ $`Ae^{ikx}+Be^{ikx};(x<0),`$
$`\psi (x)`$ $`=`$ $`Ce^{ikx};(x>0).`$ (82)
The joining conditions and the relationship $`\psi ^{}(\epsilon )\psi ^{}(\epsilon )=f\psi (0)`$ cu $`f=2mV_0/\mathrm{}^2`$ lead to
$`A+B`$ $`=`$ $`CB={\displaystyle \frac{f}{f+2ik}}A,`$
$`ik(CA+B)`$ $`=`$ $`fCC={\displaystyle \frac{2ik}{f+2ik}}A.`$ (83)
The reflection coefficient will be
$`R={\displaystyle \frac{|\psi _{refl}|^2}{|\psi _{inc}|^2}}|_{x=0}={\displaystyle \frac{|B|^2}{|A|^2}}={\displaystyle \frac{m^2V_0^2}{m^2V_0^2+\mathrm{}^4k^2}}.`$ (84)
If the potential is very strong ($`V_0\mathrm{}`$), one can see that $`R1`$, i.e., the wave is totally reflected.
The transmission coefficient, on the other hand, will be
$`T={\displaystyle \frac{|\psi _{trans}|^2}{|\psi _{inc}|^2}}|_{x=0}={\displaystyle \frac{|C|^2}{|A|^2}}={\displaystyle \frac{\mathrm{}^4k^2}{m^2V_0^2+\mathrm{}^4k^2}}.`$ (85)
Again, if the potential is very strong ($`V_0\mathrm{}`$) then $`T0`$,i.e., the transmitted wave fades rapidly on the other side of the potential.
In addition, $`R+T=1`$ as expected, which is a check of the calculation.
### Problem 2.2: Particle in a 1D potential well of finite depth
Solve the 1D Schrödinger eq. for a finite depth potential well given by
$$V(x)=\{\begin{array}{cc}V_0\hfill & \text{dacă }|x|a\hfill \\ 0\hfill & \text{dacă }|x|>a\text{ .}\hfill \end{array}$$
Consider only the bound spectrum ($`E<0`$).
Solution.
a) The wavefunction for $`|x|<a`$ and $`|x|>a`$.
The corresponding Schrödinger eq. is
$$\frac{\mathrm{}^2}{2m}\psi ^{^{\prime \prime }}(x)+V(x)\psi (x)=E\psi (x).$$
(86)
Defining
$$q^2=\frac{2mE}{\mathrm{}^2},k^2=\frac{2m(E+V_0)}{\mathrm{}^2},$$
(87)
we get:
$`1)\mathrm{for}\mathrm{x}<\mathrm{a}:\psi ^{^{\prime \prime }}_1(x)q^2\psi _1`$ $`=`$ $`0,\psi _1=A_1e^{qx}+B_1e^{qx};`$
$`2)\mathrm{for}\mathrm{a}\mathrm{x}\mathrm{a}:\psi ^{^{\prime \prime }}_2(x)+k^2\psi _2`$ $`=`$ $`0,\psi _2=A_2\mathrm{cos}(kx)+B_2\mathrm{sin}(kx);`$
$`3)\mathrm{for}\mathrm{x}>\mathrm{a}:\psi ^{^{\prime \prime }}_3(x)q^2\psi _3`$ $`=`$ $`0,\psi _3=B_3e^{qx}+B_3e^{qx}.`$
b) Formulation of the boundary conditions.
The normalization of the bound states requires solutions going to zero at infinity. This means $`B_1=A_3=0`$. Moreover, $`\psi (x)`$ should be continuously differentiable. All particular solutions are fixed in such a way that $`\psi `$ and $`\psi ^{}`$ are continuous for that value of $`x`$ corresponding to the boundary between the interior and the outside regions. The second derivative $`\psi ^{\prime \prime }`$ displays the discontinuity the ‘box’ potential imposes. Thus we are led to:
$`\psi _1(a)`$ $`=`$ $`\psi _2(a),\psi _2(a)=\psi _3(a),`$
$`\psi _1^{}(a)`$ $`=`$ $`\psi _2^{}(a),\psi _2^{}(a)=\psi _3^{}(a).`$ (88)
c) The eigenvalue equations.
From (88) we get four linear and homogeneous eqs for the coefficients $`A_1`$, $`A_2`$, $`B_2`$ and $`B_3`$:
$`A_1e^{qa}`$ $`=`$ $`A_2\mathrm{cos}(ka)B_2\mathrm{sin}(ka),`$
$`qA_1e^{qa}`$ $`=`$ $`A_2k\mathrm{sin}(ka)+B_2k\mathrm{cos}(ka),`$
$`B_3e^{qa}`$ $`=`$ $`A_2\mathrm{cos}(ka)+B_2\mathrm{sin}(ka),`$
$`qB_3e^{qa}`$ $`=`$ $`A_2k\mathrm{sin}(ka)+B_2k\mathrm{cos}(ka).`$ (89)
Adding and subtracting, one gets a system of eqs. which is easier to solve:
$`(A_1+B_3)e^{qa}`$ $`=`$ $`2A_2\mathrm{cos}(ka)`$
$`q(A_1+B_3)e^{qa}`$ $`=`$ $`2A_2k\mathrm{sin}(ka)`$
$`(A_1B_3)e^{qa}`$ $`=`$ $`2B_2\mathrm{sin}(ka)`$
$`q(A_1B_3)e^{qa}`$ $`=`$ $`2B_2k\mathrm{cos}(ka).`$ (90)
Assuming $`A_1+B_30`$ and $`A_20`$, the first two eqs give
$$q=k\mathrm{tan}(ka),$$
(91)
which inserted in the last two eqs gives
$$A_1=B_3;B_2=0.$$
(92)
The result is the symmetric solution $`\psi (x)=\psi (x)`$, also called of positive parity.
A similar calculation for $`A_1B_30`$ and $`B_20`$ leads to
$$q=k\mathrm{cot}(ka)yA_1=B_3;A_2=0.$$
(93)
The obtained wavefunction is antisymmetric, corresponding to a negative parity
d) Quantitative solution of the eigenvalue problem.
The equation connecting $`q`$ and $`k`$, already obtained previously, gives the condition to get the eigenvalues. Using the notation
$$\xi =ka,\eta =qa,$$
(94)
from the definition (87) we get
$$\xi ^2+\eta ^2=\frac{2mV_0a^2}{\mathrm{}^2}=r^2.$$
(95)
On the other hand, using (91) and (93) we get the equations
$`\eta =\xi \mathrm{tan}(\xi ),\eta =\xi \mathrm{cot}(\xi ).`$
Thus, the sought energy eigenvalues can be obtained from the intersections of these two curves with the circle defined by (95) in the plane $`\xi `$-$`\eta `$ (see fig. 2.6).
There is at least one solution for arbitrary values of the parameter $`V_0`$, in the positive parity case, because the tangent function passes through the origin. For the negative parity, the radius of the circle should be greater than a certain lower bound for the two curves to intersect. Thus, the potential should have a certain depth related to a given spatial scale $`a`$ and a given mass scale $`m`$, to allow for negative parity solutions. The number of energy levels grows with $`V_0`$, $`a`$, and $`m`$. For the case in which $`mVa^2\mathrm{}`$, the intersections are obtained from
$`\mathrm{tan}(ka)`$ $`=`$ $`\mathrm{}ka={\displaystyle \frac{2n1}{2}}\pi ,`$
$`\mathrm{cot}(ka)`$ $`=`$ $`\mathrm{}ka=n\pi ,`$ (96)
where $`n=1,2,3,\mathrm{}`$; by combining the previous relations
$$k(2a)=n\pi .$$
(97)
For the energy spectrum this fact means that
$$E_n=\frac{\mathrm{}^2}{2m}(\frac{n\pi }{2a})^2V_0.$$
(98)
Widening the well and/or the mass of the particle $`m`$, the diference between two neighbour eigenvalues will decrease. The lowest level ($`n=1`$) is not localized at $`V_0`$, but slightly upper. This ‘small’ difference is called zero point energy.
e) The forms of the wavefunctions are shown in fig. 2.7.
### Problem 2.3: Particle in 1D rectangular well of infinite depth
Solve the 1D Schrödinger eq. for a particle in a potential well of infinite depth as given by:
$$V(x)=\{\begin{array}{cc}0\hfill & \text{for }x^{}<x<x^{}+2a\hfill \\ \mathrm{}\hfill & \text{for }x^{}x\mathrm{o}xx^{}+2a\text{.}\hfill \end{array}$$
The solution in its general form is
$$\psi (x)=A\mathrm{sin}(kx)+B\mathrm{cos}(kx),$$
(99)
where
$$k=\sqrt{\frac{2mE}{\mathrm{}^2}}.$$
(100)
Since $`\psi `$ should fulfill $`\psi (x^{})=\psi (x^{}+2a)=0`$, we get:
$`A\mathrm{sin}(kx^{})+B\mathrm{cos}(kx^{})=0`$ (101)
$`A\mathrm{sin}[k(x^{}+2a)]+B\mathrm{cos}[k(x^{}+2a)]=0.`$ (102)
Multiplying (101) by $`\mathrm{sin}[k(x^{}+2a)]`$ and (102) by $`\mathrm{sin}(kx^{})`$ and next subtracting the latter result from the first we get:
$$B[\mathrm{cos}(kx^{})\mathrm{sin}[k(x^{}+2a)]\mathrm{cos}[k(x^{}+2a)]\mathrm{sin}(kx^{})]=0,$$
(103)
and by means of a trigonometric identity:
$$B\mathrm{sin}(2ak)=0$$
(104)
Multiplying (101) by $`\mathrm{cos}[k(x^{}+2a)]`$ and subtracting (102) multiplied by $`\mathrm{cos}(kx^{})`$ leads to:
$$A[\mathrm{sin}(kx^{})\mathrm{cos}[k(x^{}+2a)]\mathrm{sin}[k(x^{}+2a)]\mathrm{cos}(kx^{})]=0,$$
(105)
and by means of the same trigonometric identity:
$$A\mathrm{sin}[k(2ak)]=A\mathrm{sin}[k(2ak)]=0.$$
(106)
Since we do not take into account the trivial solution $`\psi =0`$, using (104) and (106) one has $`\mathrm{sin}(2ak)=0`$ that takes place only if $`2ak=n\pi `$, with $`n`$ an integer. Accordingly, $`k=n\pi /2a`$ and since $`k^2=2mE/\mathrm{}^2`$ then it comes out that the eigenvalues are given by the following expression:
$$E=\frac{\mathrm{}^2\pi ^2n^2}{8a^2m}.$$
(107)
The energy is quantized because only for each $`k_n=n\pi /2a`$ one gets a well-defined energy $`E_n=[n^2/2m][\pi \mathrm{}/2a]^2`$.
The general form of the solution is:
$$\psi _n=A\mathrm{sin}(\frac{n\pi x}{2a})+B\mathrm{cos}(\frac{n\pi x}{2a}),$$
(108)
and it can be normalized
$$1=_x^{}^{x^{}+2a}\psi \psi ^{}𝑑x=a(A^2+B^2),$$
(109)
wherefrom:
$$A=\pm \sqrt{1/aB^2}.$$
(110)
Substituting this value of $`A`$ in (101) one gets:
$$B=\frac{1}{\sqrt{a}}\mathrm{sin}(\frac{n\pi x^{}}{2a}),$$
(111)
and plugging $`B`$ in (110) we get
$$A=\pm \frac{1}{\sqrt{a}}\mathrm{cos}(\frac{n\pi x^{}}{2a}).$$
(112)
Using the upper signs for $`A`$ and $`B`$, by substituting their values in (108) we obtain:
$$\psi _n=\frac{1}{\sqrt{a}}\mathrm{sin}(\frac{n\pi }{2a})(xx^{}).$$
(113)
Using the lower signs for $`A`$ and $`B`$, one gets
$$\psi _n=\frac{1}{\sqrt{a}}\mathrm{sin}(\frac{n\pi }{2a})(xx^{}).$$
(114)
3. ANGULAR MOMENTUM AND SPIN
## Introduction
It is known from Classical Mechanics that the angular momentum $`𝐥`$ for macroscopic particles is given by
$$𝐥=𝐫\times 𝐩,$$
(1)
where $`𝐫`$ and $`𝐩`$ are the radius vector and the linear momentum, respectively.
However, in Quantum Mechanics, one can find operators of angular momentum type (OOAMT), which are not compulsory expressed only in terms of the coordinate $`\widehat{x}_j`$ and the momentum $`\widehat{p}_k`$ and acting only on the eigenfunctions in the x representation. Consequently, it is very important to settle first of all general commutation relations for the OOAMT components.
In Quantum Mechanics $`𝐥`$ is expressed by the operator
$$𝐥=i\mathrm{}𝐫\times ,$$
(2)
whose components are operators satisfying the following commutation rules
$$[l_x,l_y]=il_z,[l_y,l_z]=il_x,[l_z,l_x]=il_y.$$
(3)
Moreover, each of the components commutes with the square of the angular momentum, i.e.
$$l^2=l_x^2+l_y^2+l_z^2,[l_i,l^2]=0,i=1,2,3.$$
(4)
These relations, besides being correct for the angular momentum, are fulfilled for the important OOAMT class of spin operators, which miss exact analogs in classical mechanics.
These commutation relations are fundamental for getting the spectra of the aforementioned operators as well as for their differential representations.
## The angular momentum
For an arbitrary point of a fixed space (FS), one can introduce a function $`\psi (x,y,z)`$, for which let’s consider two cartesian systems $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$, where $`\mathrm{\Sigma }^{}`$ is obtained by the rotation of the $`z`$ axis of $`\mathrm{\Sigma }`$.
In the general case, a FS refers to a coordinate system, which is different of $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$.
Now, let’s compare the values of $`\psi `$ at two points of the FS with the same coordinates (x,y,z) in $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$, which is equivalent to the vectorial rotation
$$\psi (x^{},y^{},z^{})=R\psi (x,y,z)$$
(5)
where $`R`$ is a rotation matrix in R<sup>3</sup>
$$\left(\begin{array}{c}x^{}\\ y^{}\\ z^{}\end{array}\right)=\left(\begin{array}{ccc}\mathrm{cos}\varphi & \mathrm{sin}\varphi & 0\\ \mathrm{sin}\varphi & \mathrm{cos}\varphi & 0\\ 0& 0& z\end{array}\right)\left(\begin{array}{c}x\\ y\\ z\end{array}\right).$$
(6)
Then
$$R\psi (x,y,z)=\psi (x\mathrm{cos}\varphi y\mathrm{sin}\varphi ,x\mathrm{sin}\varphi +y\mathrm{cos}\varphi ,z).$$
(7)
On the other hand, it is important to recall that the wavefunctions are frame independent and that the transformation at rotations within the FS is achieved by means of unitary operators. Thus, to determine the form of the unitary operator $`U^{}(\varphi )`$ that passes $`\psi `$ to $`\psi ^{}`$, one usually considers an infinitesimal rotation $`d\varphi `$, keeping only the linear terms in $`d\varphi `$ when one expands $`\psi ^{}`$ in Taylor series in the neighborhood of $`x`$
$`\psi (x^{},y^{},z^{})`$ $``$ $`\psi (x+yd\varphi ,xd\varphi +y,z),`$ (8)
$``$ $`\psi (x,y,z)+d\varphi \left(y{\displaystyle \frac{\psi }{x}}x{\displaystyle \frac{\psi }{y}}\right),`$
$``$ $`(1id\varphi l_z)\psi (x,y,z),`$
where we have used the notation<sup>1</sup><sup>1</sup>1The proof of (8) is displayed as problem 3.1
$$l_z=\mathrm{}^1(\widehat{x}\widehat{p}_y\widehat{y}\widehat{p}_x).$$
(9)
As one will see later, this corresponds to the projection operator onto $`z`$ of the angular momentum according to the definition (2) unless the factor $`\mathrm{}^1`$. In this way, the rotations of finite angle $`\varphi `$ can be represented as exponentials of the form
$$\psi (x^{},y^{},z)=e^{il_z\varphi }\psi (x,y,z),$$
(10)
where
$$\widehat{U}^{}(\varphi )=e^{il_z\varphi }.$$
(11)
In order to reassert the concept of rotation, we will consider it in a more general approach with the help of the vectorial operator $`\widehat{\stackrel{}{A}}`$ acting on $`\psi `$, assuming that $`\widehat{A}_x`$, $`\widehat{A}_y`$, $`\widehat{A}_z`$ have the same form in $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$, that is, the mean values of $`\widehat{\stackrel{}{A}}`$ as calculated in $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$ should be equal when they are seen from the FS
$`{\displaystyle \psi ^{}(\stackrel{}{r}^{})(\widehat{A}_x\widehat{ı}^{}+\widehat{A}_y\widehat{ȷ}^{}+\widehat{A}_z\widehat{k}^{})\psi ^{}(\stackrel{}{r}^{})𝑑\stackrel{}{r}}`$
$`={\displaystyle \psi ^{}(\stackrel{}{r})(\widehat{A}_x\widehat{ı}+\widehat{A}_y\widehat{ȷ}+\widehat{A}_z\widehat{k})\psi ^{}(\stackrel{}{r})𝑑\stackrel{}{r}},`$ (12)
where
$$\widehat{ı}^{}=\widehat{ı}\mathrm{cos}\varphi +\widehat{ȷ}\mathrm{sin}\varphi ,\widehat{ȷ}^{}=\widehat{ı}\mathrm{sin}\varphi +\widehat{ȷ}\mathrm{cos}\varphi ,\widehat{k}^{}=\widehat{k}.$$
(13)
Thus, by combining (10), (12) and (13) we get
$`e^{il_z\varphi }\widehat{A}_xe^{il_z\varphi }`$ $`=`$ $`\widehat{A}_x\mathrm{cos}\varphi \widehat{A}_y\mathrm{sin}\varphi ,`$
$`e^{il_z\varphi }\widehat{A}_ye^{il_z\varphi }`$ $`=`$ $`\widehat{A}_x\mathrm{sin}\varphi \widehat{A}_y\mathrm{cos}\varphi ,`$
$`e^{il_z\varphi }\widehat{A}_ze^{il_z\varphi }`$ $`=`$ $`\widehat{A}_z.`$ (14)
Again, considering infinitesimal rotations and expanding the left hand sides in (14), one can determine the commutation relations of $`\widehat{A}_x`$, $`\widehat{A}_y`$ and $`\widehat{A}_z`$ with $`\widehat{l}_z`$
$$[l_z,A_x]=iA_y,[l_z,A_y]=iA_x,[l_z,A_z]=0,$$
(15)
and similarly for $`l_x`$ and $`l_y`$.
The basic conditions to obtain these commutation relations are
* The eigenfunctions transform as in (7) when $`\mathrm{\Sigma }\mathrm{\Sigma }^{}`$.
* The components $`\widehat{A}_x`$, $`\widehat{A}_y`$, $`\widehat{A}_z`$ have the same form in $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$.
* The kets corresponding to the mean values of $`\widehat{A}`$ in $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$ coincide (are the same) for a FS observer.
One can also use another representation in which $`\psi (x,y,z)`$ does not change when $`\mathrm{\Sigma }\mathrm{\Sigma }^{}`$ and the vectorial operators transform as ordinary vectors. In order to pass to such a representation when we rotate by $`\varphi `$ around $`z`$ one makes use of the operator $`\widehat{U}(\varphi )`$, that is
$$e^{il_z\varphi }\psi ^{}(x,y,z)=\psi (x,y,z),$$
(16)
and therefore
$$e^{il_z\varphi }\widehat{\stackrel{}{A}}e^{il_z\varphi }=\widehat{\stackrel{}{A}}.$$
(17)
Using the relationships (14) we obtain
$`\widehat{A}_x^{}`$ $`=`$ $`\widehat{A}_x\mathrm{cos}\varphi +\widehat{A}_y\mathrm{sin}\varphi =e^{il_z\varphi }\widehat{A}_xe^{il_z\varphi },`$
$`\widehat{A}_y^{}`$ $`=`$ $`\widehat{A}_x\mathrm{sin}\varphi +\widehat{A}_y\mathrm{cos}\varphi =e^{il_z\varphi }\widehat{A}_ye^{il_z\varphi },`$
$`\widehat{A}_z^{}`$ $`=`$ $`e^{il_z\varphi }\widehat{A}_ze^{il_z\varphi }.`$ (18)
Since the transformations of the new representation are performed by means of unitary operators, the commutation relations do not change.
### Remarks
* The operator $`\widehat{A}^2`$ is invariant at rotations, that is
$$e^{il_z\varphi }\widehat{A}^2e^{il_z\varphi }=\widehat{A}^2=\widehat{A}^2.$$
(19)
* It follows that
$$[\widehat{l}_i,\widehat{A}^2]=0.$$
(20)
* If the Hamiltonian operator is of the form
$$\widehat{H}=\frac{1}{2m}\widehat{p}^2+U(|\stackrel{}{r}|),$$
(21)
then it remains invariant under rotations in any axis passing through the coordinate origin
$$[\widehat{l}_i,\widehat{H}]=0,$$
(22)
where $`\widehat{l}_i`$ are integrals of the motion.
### Definition
If $`\widehat{A}_i`$ are the components of a vectorial operator acting on a wavefunction depending only on the coordinates and if there are operators $`\widehat{l}_i`$ that satisfy the following commutation relations
$$[\widehat{l}_i,\widehat{A}_j]=i\epsilon _{ijk}\widehat{A}_k,[\widehat{l}_i,\widehat{l}_j]=i\epsilon _{ijk}\widehat{l}_k,$$
(23)
then $`\widehat{l}_i`$ are known as the components of the angular momentum operator and we can infer from (20) and (23) that
$$[\widehat{l}_i,\widehat{l}^2]=0.$$
(24)
Consequently the three operatorial components associated to the components of a classical angular momentum satisfy commutation relations of the type (23), (24). Moreover, one can prove that these relations lead to specific geometric properties of the rotations in a 3D euclidean space. This takes place if we adopt a more general point of view by defining an angular momentum operator $`𝐉`$ (we shall not use the hat symbol for simplicity of writing) as any set of three observables $`J_x`$, $`J_y`$ şi $`J_z`$ which fulfill the commutation relations
$$[J_i,J_j]=i\epsilon _{ijk}J_k.$$
(25)
Moreover, let us introduce the operator
$$𝐉^2=J_x^2+J_y^2+J_z^2,$$
(26)
the scalar square of the angular momentum $`𝐉`$. This operator is hermitic because $`J_x`$, $`J_y`$ and $`J_z`$ are hermitic and it is easy to show that $`𝐉^\mathrm{𝟐}`$ commutes with the three components of $`𝐉`$
$$[𝐉^2,J_i]=0.$$
(27)
Since $`𝐉^\mathrm{𝟐}`$ commutes with each of the components it follows that there is a complete system of eigenfunctions, i.e.
$$𝐉^\mathrm{𝟐}\psi _{\gamma \mu }=f(\gamma ^2)\psi _{\gamma \mu },J_i\psi _{\gamma \mu }=g(\mu )\psi _{\gamma \mu },$$
(28)
where, as it will be shown in the following, the eigenfunctions depend on two subindices, which will be determined together with the form of the functions $`f(\gamma )`$ and $`g(\mu )`$. The operators $`J_i`$ and $`J_k`$ $`(ik)`$ do not commute, i.e. they do not have common eigenfunctions. For physical and mathematical reasons, we are interested to determine the common eigenfunctions of $`𝐉^\mathrm{𝟐}`$ and $`J_z`$, that is, we shall take $`i=z`$ in (28).
Instead of using the components $`J_x`$ and $`J_y`$ of the angular momentum $`𝐉`$, it is more convenient to work with the following linear combinations
$$J_+=J_x+iJ_y,J_{}=J_xiJ_y.$$
(29)
Contrary to the operators $`a`$ and $`a^{}`$ of the harmonic oscillator (see chapter 5), these operators are not hermitic, they are only adjunct to each other. The following properties are easy to prove
$$[J_z,J_\pm ]=\pm J_\pm ,[J_+,J_{}]=2J_z,$$
(30)
$$[J^2,J_+]=[J^2,J_{}]=[J^2,J_z]=0.$$
(31)
$$J_z(J_\pm \psi _{\gamma \mu })=\{J_\pm J_z+[J_z,J_\pm ]\}\psi _{\gamma \mu }=(\mu \pm 1)(J_\pm \psi _{\gamma \mu }).$$
(32)
Therefore $`J_\pm \psi _{\gamma \mu }`$ are eigenfunctions of $`J_z`$ corresponding to the eigenvalues $`\mu \pm 1`$, that is these functions are identical up to the constant factors $`\alpha _\mu `$ and $`\beta _\mu `$ (to be determined)
$`J_+\psi _{\gamma \mu 1}`$ $`=`$ $`\alpha _\mu \psi _{\gamma \mu },`$
$`J_{}\psi _{\gamma \mu }`$ $`=`$ $`\beta _\mu \psi _{\gamma \mu 1}.`$ (33)
On the other hand
$$\alpha _\mu ^{}=(J_+\psi _{\gamma \mu 1},\psi _{\gamma \mu })=(\psi _{\gamma \mu 1}J_{}\psi _{\gamma \mu })=\beta _\mu .$$
(34)
Therefore, taking a phase of the type $`e^{ia}`$ (where $`a`$ is real) for the function $`\psi _{\gamma \mu }`$ one can put $`\alpha _\mu `$ real and equal to $`\beta _\mu `$, which means
$$J_+\psi _{\gamma ,\mu 1}=\alpha \mu \psi _{\gamma \mu },J_{}\psi _{\gamma \mu }=\alpha \mu \psi _{\gamma ,\mu 1},$$
(35)
and therefore
$`\gamma `$ $`=`$ $`(\psi _{\gamma \mu },[J_x^2+J_y^2+J_z^2]\psi _{\gamma \mu })=\mu ^2+a+b,`$
$`a`$ $`=`$ $`(\psi _{\gamma \mu },J_x^2\psi _{\gamma \mu })=(J_x\psi _{\gamma \mu },J_x\psi _{\gamma \mu })0,`$
$`b`$ $`=`$ $`(\psi _{\gamma \mu },J_y^2\psi _{\gamma \mu })=(J_y\psi _{\gamma \mu },J_y\psi _{\gamma \mu })0.`$ (36)
The normalization constant cannot be negative. This implies
$$\gamma \mu ^2,$$
(37)
for a fixed $`\gamma `$; thus, $`\mu `$ has both superior and inferior limits (it takes values in a finite interval).
Let $`\mathrm{\Lambda }`$ and $`\lambda `$ be these limits, respectively, for a given $`\gamma `$
$$J_+\psi _{\gamma \mathrm{\Lambda }}=0,J_{}\psi _{\gamma \lambda }=0.$$
(38)
Using the following operatorial identities
$`J_{}J_+`$ $`=`$ $`𝐉^\mathrm{𝟐}J_z^2+J_z=𝐉^\mathrm{𝟐}J_z(J_z1),`$
$`J_+J_{}`$ $`=`$ $`𝐉^\mathrm{𝟐}J_z^2+J_z=𝐉^\mathrm{𝟐}J_z(J_z+1),`$ (39)
acting on $`\psi _{\gamma \mathrm{\Lambda }}`$ as well as on $`\psi _{\gamma \lambda }`$ one gets
$`\gamma \mathrm{\Lambda }^2\mathrm{\Lambda }`$ $`=`$ $`0,`$
$`\gamma \lambda ^2+\lambda `$ $`=`$ $`0,`$
$`(\lambda \lambda +1)(\lambda +\lambda )`$ $`=`$ $`0.`$ (40)
In addition,
$$\mathrm{\Lambda }\lambda \mathrm{\Lambda }=\lambda =J\gamma =J(J+1).$$
(41)
For a given $`\gamma `$ the projection $`\mu `$ of the momentum takes $`2J+1`$ values that differ by unity, from $`J`$ to $`J`$. Therefore, the difference $`\mathrm{\Lambda }\lambda =2J`$ should be an integer and consequently the eigenvalues of $`J_z`$ that are labelled by $`m`$ are integer
$$m=k,k=0,\pm 1,\pm 2,\mathrm{},$$
(42)
or half-integer
$$m=k+\frac{1}{2},k=0,\pm 1,\pm 2,\mathrm{}.$$
(43)
A state having a given $`\gamma =J(J+1)`$ presents a degeneration of order $`g=2J+1`$ with regard to the eigenvalues $`m`$ (this is so because $`J_i,J_k`$ commute with $`J^2`$ but do not commute between themselves.
By a “state of angular momentum $`J`$” one usually understands a state of $`\gamma =J(J+1)`$ having the maximum projection of its momentum, i.e. $`J`$. Quite used notations for angular momentum states are $`\psi _{jm}`$ and the Dirac ket one $`|jm`$.
Let us now obtain the matrix elements of $`J_x,J_y`$ in the representation in which $`J^2`$ and $`J_z`$ are diagonal. In this case, one obtains from (35) and (39) the following relations
$`J_{}J_+\psi _{jm1}=\alpha _mJ_{}\psi _{jm}=\alpha _m\psi _{jm1},`$
$`J(J+1)(m1)^2(m1)=\alpha _m^2,`$
$`\alpha _m=\sqrt{(J+m)(Jm+1)}.`$ (44)
Combining (44) and (35) leads to
$$J_+\psi _{jm1}=\sqrt{(J+m)(Jm+1)}\psi _{jm}.$$
(45)
It follows that the matrix element of $`J_+`$ is
$$jm|J_+|jm1=\sqrt{(J+m)(Jm+1)}\delta _{nm},$$
(46)
and analogously
$$jn|J_{}|jm=\sqrt{(J+m)(Jm+1)}\delta _{nm1}.$$
(47)
Finally, from the definitions (29) for $`J_+,J_{}`$ one easily gets
$`jm|J_x|jm1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sqrt{(J+m)(Jm+1)},`$
$`jm|J_y|jm1`$ $`=`$ $`{\displaystyle \frac{i}{2}}\sqrt{(J+m)(Jm+1)}.`$ (48)
### Partial conclusions
* Properties of the eigenvalues of $`𝐉`$ and $`J_z`$
If $`j(j+1)\mathrm{}^2`$ and $`m\mathrm{}`$ are eigenvalues of $`𝐉`$ and $`J_z`$ associated to the eigenvectors $`|kjm`$, then $`j`$ and $`m`$ satisfy the inequality
$$jmj.$$
* Properties of the vector $`J_{}|kjm`$
Let $`|kjm`$ be an eigenvector of $`𝐉^\mathrm{𝟐}`$ and $`J_z`$ with the eigenvalues $`j(j+1)\mathrm{}^2`$ and $`m\mathrm{}`$
+ (i) If $`m=j`$, then $`J_{}|kjj=0`$.
+ (ii) If $`m>j`$, then $`J_{}|kjm`$ is a nonzero eigenvector of $`J^2`$ and $`J_z`$ with the eigenvalues $`j(j+1)\mathrm{}^2`$ and $`(m1)\mathrm{}`$.
* Properties of the vector $`J_+|kjm`$
Let $`|kjm`$ be a (ket) eigenvector of $`𝐉^\mathrm{𝟐}`$ and $`J_z`$ for the eigenvalues $`j(j+1)\mathrm{}^2`$ and $`m\mathrm{}`$
+ If $`m=j`$, then $`J_+|kjm=0.`$
+ If $`m<j`$, then $`J_+|kjm`$ is a nonzero eigenvector of $`𝐉^\mathrm{𝟐}`$ and $`J_z`$ with the eigenvalues $`j(j+1)\mathrm{}^2`$ and $`(m+1)\mathrm{}`$
* Consequences of the previous properties
$`J_z|kjm`$ $`=`$ $`m\mathrm{}|kjm,`$
$`J_+|kjm`$ $`=`$ $`m\mathrm{}\sqrt{j(j+1)m(m+1)}|kjm+1,`$
$`J_{}|kjm`$ $`=`$ $`m\mathrm{}\sqrt{j(j+1)m(m1)}|kjm+1.`$
## Applications of the orbital angular momentum
Until now we have considered those properties of the angular momentum that could be derived only from the commutation relations. Let us go back to the orbital momentum $`𝐥`$ of a particle without intrinsic rotation and let us examine how one can apply the theory of the previous section in the important particular case
$$[\widehat{l}_i,\widehat{p}_j]=i\epsilon _{ijk}\widehat{p}_k.$$
(49)
First, $`\widehat{l}_z`$ and $`\widehat{p}_j`$ have a common system of eigenfunctions. On the other hand, the Hamiltonian of a free particle
$$\widehat{H}=\left(\frac{\widehat{\stackrel{}{p}}}{\sqrt{2m}}\right)^2,$$
being the square of a vectorial operator has a complete system of eigenfunctions with $`\widehat{L^2}`$ and $`\widehat{l}_z`$. In addition, this implies that the free particle can be found in a state of well-defined $`E`$, $`l`$, and $`m`$.
### Eigenvalues and eigenfunctions of $`𝐥^\mathrm{𝟐}`$ and $`𝐥_𝐳`$
It is more convenient to work in spherical coordinates because, as we will see, various angular momentum operators act only on the angle variables $`\theta ,\varphi `$ and not on $`r`$. Thus, instead of describing $`r`$ by its cartesian components $`x,y,z`$ we determine the arbitrary point $`M`$ of vector radius $`𝐫`$ by the spherical 3D coordinates
$$x=r\mathrm{cos}\varphi \mathrm{sin}\theta ,y=r\mathrm{sin}\varphi \mathrm{sin}\theta ,z=r\mathrm{cos}\theta ,$$
(50)
where
$$r0,0\theta \pi ,0\varphi 2\pi .$$
Let $`\mathrm{\Phi }(r,\theta ,\varphi )`$ and $`\mathrm{\Phi }^{}(r,\theta ,\varphi )`$ be the wavefunctions of a particle in $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$, respectively, in which the infinitesimal rotation is given by $`\delta \alpha `$ around the $`z`$ axis
$`\mathrm{\Phi }^{}(r,\theta ,\varphi )`$ $`=`$ $`\mathrm{\Phi }(r,\theta ,\varphi +\delta \alpha ),`$ (51)
$`=`$ $`\mathrm{\Phi }(r,\theta ,\varphi )+\delta \alpha {\displaystyle \frac{\mathrm{\Phi }}{\varphi }},`$
or
$$\mathrm{\Phi }^{}(r,\theta ,\varphi )=(1+i\widehat{l}_z\delta \alpha )\mathrm{\Phi }(r,\theta ,\varphi ).$$
(52)
It follows that
$$\frac{\mathrm{\Phi }}{\varphi }=i\widehat{l_z}\mathrm{\Phi },\widehat{l}_z=i\frac{}{\varphi }.$$
(53)
For an inifinitesimal rotation in $`x`$
$`\mathrm{\Phi }^{}(r,\theta ,\varphi )`$ $`=`$ $`\mathrm{\Phi }+\delta \alpha \left({\displaystyle \frac{\mathrm{\Phi }}{\theta }}{\displaystyle \frac{\theta }{\alpha }}+{\displaystyle \frac{\mathrm{\Phi }}{\theta }}{\displaystyle \frac{\varphi }{\alpha }}\right),`$ (54)
$`=`$ $`(1+i\widehat{l}_x\delta \alpha )\mathrm{\Phi }(r,\theta ,\varphi ),`$
but in this rotation
$$z^{}=z+y\delta \alpha ;z^{}=z+y\delta \alpha ;x^{}=x$$
(55)
and from (50) one gets
$`r\mathrm{cos}(\theta +d\theta )`$ $`=`$ $`r\mathrm{cos}\theta +r\mathrm{sin}\theta \mathrm{sin}\varphi \delta \alpha ,`$
$`r\mathrm{sin}\varphi \mathrm{sin}(\theta +d\theta )`$ $`=`$ $`r\mathrm{sin}\theta \mathrm{sin}\varphi +r\mathrm{sin}\theta \mathrm{sin}\varphi r\mathrm{cos}\theta \delta \alpha ,`$ (56)
i.e.
$$\mathrm{sin}\theta d\theta =\mathrm{sin}\theta \mathrm{sin}\varphi \delta \alpha \frac{d\theta }{d\alpha }=\mathrm{sin}\varphi ,$$
(57)
and
$`\mathrm{cos}\theta \mathrm{sin}\varphi d\theta +\mathrm{sin}\theta \mathrm{cos}\varphi d\varphi `$ $`=`$ $`\mathrm{cos}\theta \delta \alpha ,`$
$`\mathrm{cos}\varphi \mathrm{sin}\theta {\displaystyle \frac{d\varphi }{d\alpha }}`$ $`=`$ $`\mathrm{cos}\theta \mathrm{cos}\theta \mathrm{sin}\varphi {\displaystyle \frac{d\theta }{d\alpha }}.`$ (58)
Substituting (57) in (56) leads to
$$\frac{d\varphi }{d\alpha }=\mathrm{cot}\theta \mathrm{cos}\varphi .$$
(59)
With (56) and (58) substituted in (51) and comparing the right hand sides of (51) one gets
$$\widehat{l}_x=i\left(\mathrm{sin}\varphi \frac{}{\theta }+\mathrm{cot}\theta \mathrm{cos}\varphi \frac{}{\varphi }\right).$$
(60)
For the rotation in $`y`$, the result is similar
$$\widehat{l}_y=i\left(\mathrm{cos}\varphi \frac{}{\theta }+\mathrm{cot}\theta \mathrm{sin}\varphi \frac{}{\varphi }\right).$$
(61)
Using $`\widehat{l}_x,\widehat{l}_y`$ one can also obtain $`\widehat{l}_\pm ,\widehat{l}^2`$
$`\widehat{l}_\pm `$ $`=`$ $`\mathrm{exp}\left[\pm i\varphi \left(\pm {\displaystyle \frac{}{\theta }}+i\mathrm{cot}\theta {\displaystyle \frac{}{\varphi }}\right)\right],`$
$`\widehat{l}^2`$ $`=`$ $`\widehat{l}_{}\widehat{l}_++\widehat{l}^2+\widehat{l}_z,`$ (62)
$`=`$ $`\left[{\displaystyle \frac{1}{\mathrm{sin}^2\theta }}{\displaystyle \frac{^2}{\varphi ^2}}+{\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{}{\theta }}\left(\mathrm{sin}\theta {\displaystyle \frac{}{\theta }}\right)\right].`$
From (62) one can see that $`\widehat{l}^2`$ is identical up to a constant to the angular part of the Laplace operator at a fixed radius
$$^2f=\frac{1}{r^2}\frac{}{r}\left(r^2\frac{f}{r}\right)+\frac{1}{r^2}\left[\frac{1}{\mathrm{sin}\theta }\frac{}{\theta }\left(\mathrm{sin}\theta \frac{f}{\theta }\right)+\frac{1}{\mathrm{sin}^2\theta }\frac{^2}{\varphi ^2}\right].$$
(63)
### The eigenfunctions of $`l_z`$
$`\widehat{l}_z\mathrm{\Phi }_m=m\mathrm{\Phi }=i{\displaystyle \frac{\mathrm{\Phi }_m}{\varphi }},`$
$`\mathrm{\Phi }_m={\displaystyle \frac{1}{\sqrt{2\pi }}}e^{im\varphi }.`$ (64)
### Hermiticity conditions of $`\widehat{l}_z`$
$$_0^{2\pi }f^{}\widehat{l}_zg𝑑\varphi =\left(_0^{2\pi }g^{}\widehat{l}_zf𝑑\varphi \right)^{}+f^{}g(2\pi )f^{}g(0).$$
(65)
It follows that $`\widehat{l}_z`$ is hermitic in the class of functions for which
$$f^{}g(2\pi )=f^{}g(0).$$
(66)
The eigenfunctions $`\mathrm{\Phi }_m`$ of $`\widehat{l}_z`$ belong to the integrable class $`^2(0,2\pi )`$ and they fulfill (66), as it happens for any function that can be expanded in $`\mathrm{\Phi }_m(\varphi )`$
$`F(\varphi )`$ $`=`$ $`{\displaystyle \stackrel{k}{}}a_ke^{ik\varphi },k=0,\pm 1,\pm 2,\mathrm{},`$
$`G(\varphi )`$ $`=`$ $`{\displaystyle \stackrel{k}{}}b_ke^{ik\varphi },k=\pm 1/2,\pm 3/2,\pm 5/2\mathrm{},`$ (67)
with $`k`$ only integers or half-integers, but not for combinations of $`F(\varphi )`$ and $`G(\varphi )`$. The correct choice of $`m`$ is based on the common eigenfunctions of $`\widehat{l}_z`$ and $`\widehat{l}^2`$.
### Spherical harmonics
In the $`\{\stackrel{}{𝐫}\}`$ representation, the eigenfunctions associated to the eigenvalues $`l(l+1)\mathrm{}^2`$ of $`𝐥^\mathrm{𝟐}`$ and $`m\mathrm{}`$ of $`l_z`$ are solutions of the partial differential equations
$`\left({\displaystyle \frac{^2}{\theta ^2}}+{\displaystyle \frac{1}{\mathrm{tan}\theta }}{\displaystyle \frac{}{\theta }}+{\displaystyle \frac{1}{\mathrm{sin}^2\theta }}{\displaystyle \frac{^2}{\varphi ^2}}\right)\psi (r,\theta ,\varphi )`$ $`=`$ $`l(l+1)\mathrm{}^2\psi (r,\theta ,\varphi ),`$
$`i{\displaystyle \frac{}{\varphi }}\psi (r,\theta ,\varphi )`$ $`=`$ $`m\mathrm{}\psi (r,\theta ,\varphi ).`$ (68)
Taking into account that the general results presented above can be applied to the orbital momentum, we infer that $`l`$ can be an integer or half-integer and that, for fixed $`l`$, $`m`$ can only take the values
$$l,l+1,\mathrm{},l1,l.$$
In (68), $`r`$ is not present in the differential operator, so that it can be considered as a parameter. Thus, considering only the dependence on $`\theta ,\varphi `$ of $`\psi `$, one uses the notation $`Y_{lm}(\theta ,\varphi )`$ for these common eigenfunctions of $`𝐥^\mathrm{𝟐}`$ and $`l_z`$, corresponding to the eigenvalues $`l(l+1)\mathrm{}^2,m\mathrm{}`$. They are known as spherical harmonics
$`𝐥^\mathrm{𝟐}Y_{lm}(\theta ,\varphi )`$ $`=`$ $`l(l+1)\mathrm{}^2Y_{lm}(\theta ,\varphi ),`$
$`l_zY_{lm}(\theta ,\varphi )`$ $`=`$ $`m\mathrm{}Y_{lm}(\theta ,\varphi ).`$ (69)
For more rigorousness, one should introduce one more index in order to distinguish among the various solutions of (69) corresponding to the same $`(l,m)`$ pairs for particles with spin. If the spin is not taken into account, these equations have a unique solution (up to a constant factor) for each allowed pair of $`(l,m)`$; this is so because the subindices $`l,m`$ are sufficient in this context. The solutions $`Y_{lm}(\theta ,\varphi )`$ have been found by the method of the separation of variables in spherical variables (see also the chapter The hydrogen atom)
$$\psi _{lm}(r,\theta ,\varphi )=f(r)\psi _{lm}(\theta ,\varphi ),$$
(70)
where $`f(r)`$ is a function of $`r`$, which looks as an integration constant from the viewpoint of the partial differential equations in (68). The fact that $`f(r)`$ is arbitrary proves that $`𝐋^\mathrm{𝟐}`$ and $`l_z`$ do not form a complete set of observables<sup>2</sup><sup>2</sup>2By definition, the hermitic operator A is an observable if the orthogonal system of eigenvectors form a base in the space of states. in the space $`\epsilon _r`$<sup>3</sup><sup>3</sup>3Each quantum state of a particle is characterized by a vectorial state belonging to an abstract vectorial space $`\epsilon _r`$. of functions of $`\stackrel{}{r}`$ ($`r,\theta ,\varphi `$).
In order to normalize $`\psi _{lm}(r,\theta ,\varphi )`$, it is convenient to normalize $`Y_{lm}(\theta ,\varphi )`$ and $`f(r)`$ separately
$`{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _0^\pi }\mathrm{sin}\theta |\psi _{lm}(\theta ,\varphi )|^2d\theta `$ $`=`$ $`1,`$ (71)
$`{\displaystyle _0^{\mathrm{}}}r^2|f(r)|^2𝑑r`$ $`=`$ $`1.`$ (72)
### The values of the pair $`(l,m)`$
($`\alpha `$): $`l,m`$ should be integers
Using $`l_z=\frac{\mathrm{}}{i}\frac{}{\varphi }`$, we can write (69) as follows
$$\frac{\mathrm{}}{i}\frac{}{\varphi }Y_{lm}(\theta ,\varphi )=m\mathrm{}Y_{lm}(\theta ,\varphi ).$$
(73)
Thus,
$$Y_{lm}(\theta ,\varphi )=F_{lm}(\theta ,\varphi )e^{im\varphi }.$$
(74)
If $`0\varphi <2\pi `$, then we should tackle the condition of covering all space according to the requirement of dealing with a function continuous in any angular zone, i.e. că
$$Y_{lm}(\theta ,\varphi =0)=Y_{lm}(\theta ,\varphi =2\pi ),$$
(75)
implying
$$e^{im\pi }=1.$$
(76)
As has been seen, $`m`$ is either an integer or a half-integer; for the application to the orbital momentum, $`m`$ should be an integer. ($`e^{2im\pi }`$ would be $`1`$ if $`m`$ is a half-integer).
($`\beta `$): For a given value of $`l`$, all the corresponding $`Y_{lm}`$ can be obtained by algebraic means using
$$l_+Y_{ll}(\theta ,\varphi )=0,$$
(77)
which combined with eq. (62) for $`l_+`$ leads to
$$\left(\frac{d}{d\theta }l\mathrm{cot}\theta \right)F_{ll}(\theta )=0.$$
(78)
This equation can be immediately integrated if we notice the relationship
$$\mathrm{cot}\theta d\theta =\frac{d(\mathrm{sin}\theta )}{\mathrm{sin}\theta }.$$
(79)
Its general solution is
$$F_{ll}=c_l(\mathrm{sin}\theta )^l,$$
(80)
where $`c_l`$ is a normalization constant.
It follows that for any positive or zero value of $`l`$, there is a function $`Y_{ll}(\theta ,\varphi )`$, which up to a constant factor is
$$Y^{ll}(\theta ,\varphi )=c_l(\mathrm{sin}\theta )^le^{il\varphi }.$$
(81)
Using repeatedly the action of $`l_{}`$, one can build the whole set of functions $`Y_{ll1}(\theta ,\varphi ),\mathrm{},Y_{l0}(\theta ,\varphi ),`$ $`\mathrm{},Y_{ll}(\theta ,\varphi )`$. Next, we look at the way in which these functions can be put into correspondence with the eigenvalue pair $`l(l+1)\mathrm{},m\mathrm{}`$ (where $`l`$ is an arbitrary positive integer such that $`lml`$ ). Using (78), we can make the conclusion that any other eigenfunction $`Y_{lm}(\theta ,\varphi )`$ can be unambigously obtained from $`Y_{ll}`$.
### Properties of spherical harmonics
$`\alpha `$ Iterative relationships
From the general results of this chapter, we have
$$l_\pm Y_{lm}(\theta ,\varphi )=\mathrm{}\sqrt{l(l+1)m(m\pm 1)}Y_{lm\pm 1}(\theta ,\varphi ).$$
(82)
Using (62) for $`l_\pm `$ and the fact that $`Y_{lm}(\theta ,\varphi )`$ is the product of a $`\theta `$-dependent function and $`e^{\pm i\varphi }`$, one gets
$$e^{\pm i\varphi }\left(\frac{}{\theta }m\mathrm{cot}\theta \right)Y_{lm}(\theta ,\varphi )=\sqrt{l(l+1)m(m\pm 1)}Y_{lm\pm 1}(\theta ,\varphi )$$
(83)
$`\beta `$ Orthonormalization and completeness relationships
Equation (68) determines the spherical harmonics only up to a constant factor. We shall now choose this factor such that to have the orthonormalization of $`Y_{lm}(\theta ,\varphi )`$ (as functions of the angular variables $`\theta ,\varphi `$)
$$_0^{2\pi }𝑑\varphi _0^\pi \mathrm{sin}\theta d\theta Y_{lm}^{}(\theta ,\varphi )Y_{lm}(\theta ,\varphi )=\delta _{l^{}l}\delta _{m^{}m}.$$
(84)
In addition, any continuous function of $`\theta ,\varphi `$ can be expressed by means of the spherical harmonics as follows
$$f(\theta ,\varphi )=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}c_{lm}Y_{lm}(\theta ,\varphi ),$$
(85)
where
$$c_{lm}=_0^{2\pi }𝑑\varphi _0^\pi \mathrm{sin}\theta d\theta Y_{lm}^{}(\theta ,\varphi )f(\theta ,\varphi ).$$
(86)
The results (85), (86) are consequences of defining the spherical harmonics as an orthonormalized and complete base in the space $`\epsilon _\mathrm{\Omega }`$ of functions of $`\theta ,\varphi `$. The completeness relationship is
$`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}Y_{lm}(\theta ,\varphi )Y_{lm}^{}(\theta ^{},\varphi )`$ $`=`$ $`\delta (\mathrm{cos}\theta \mathrm{cos}\theta ^{})\delta (\varphi ,\varphi ),`$ (87)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{sin}\theta }}\delta (\theta \theta ^{})\delta (\varphi ,\varphi ).`$
The ‘function’ $`\delta (\mathrm{cos}\theta \mathrm{cos}\theta ^{})`$ occurs because the integral over the variable $`\theta `$ is performed by using the differential element $`\mathrm{sin}\theta d\theta =d(\mathrm{cos}\theta )`$.
### Parity operator $`𝒫`$ for spherical harmonics
The behavior of $`𝒫`$ in 3D is rather close to that in 1D. When it is applied to a function of cartesian coordinates $`\psi (x,y,z)`$ changes the sign of each of the coordinates
$$𝒫\psi (x,y,z)=\psi (x,y,z).$$
(88)
$`𝒫`$ has the properties of a hermitic operator, being also a unitary operator, as well as a projector since $`𝒫^2`$ is an identity operator
$`\psi (𝐫)|𝒫|\psi (𝐫)=\psi (𝐫)|\psi (𝐫)=\psi (𝐫)|\psi (𝐫^{}),`$
$`𝒫^2|𝐫=𝒫(𝒫|𝐫)=𝒫|𝐫=|𝐫.`$ (89)
Therefore
$$𝒫^2=\widehat{1},$$
(90)
for which the eigenvalues are $`P=\pm 1`$. The eigenfunctions are called even if $`P=1`$ and odd if $`P=1`$. In nonrelativistic quantum mechanics, the operator $`\widehat{H}`$ for a conservative system is invariant with regard to discrete unitary transformations, i.e.
$$𝒫\widehat{H}𝒫=𝒫^1\widehat{H}𝒫=\widehat{H}.$$
(91)
Thus, $`\widehat{H}`$ commutes with $`𝒫`$ and the parity of the state is a constant of the motion. In addition, $`𝒫`$ commutes with the operators $`\widehat{l}`$ and $`\widehat{l}_\pm `$
$$[𝒫,\widehat{l}_i]=0,[𝒫,\widehat{l}_\pm ]=0.$$
(92)
Because of all these properties, one can have the important class of wave functions which are common eigenfunctions of the triplet $`𝒫,\widehat{l}^2`$ and $`\widehat{l}_z`$. It follows from (92) that the parities of the states which difer only in $`\widehat{l}_z`$ coincide. In this way, one can identify the parity of a particle of definite orbital angular momentum $`\widehat{l}`$.
In spherical coordinates, we shall consider the following change of variables
$$rr,\theta \pi \theta \varphi \pi +\varphi .$$
(93)
Thus, using a standard base in the space of wavefunctions of a particle without ‘intrinsic rotation’, the radial part of the base functions $`\psi _{klm}(\stackrel{}{r})`$ is not changed by the parity operator. Only the spherical harmonics will change. From the trigonometric standpoint, the transformations (93) are as follows
$$\mathrm{sin}(\pi \theta )\mathrm{sin}\theta ,\mathrm{cos}(\pi \theta )\mathrm{cos}\theta e^{im(\pi +\varphi }(1)^me^{im\varphi }$$
(94)
leading to the following transformation of the function $`Y_{ll}(\theta ,\varphi )`$
$$Y_{ll}(\varphi \theta ,\pi +\varphi )=(1)^lY_{ll}(\theta ,\varphi ).$$
(95)
From (95) it follows that the parity of $`Y_{ll}`$ is $`(1)^l`$. On the other hand, $`l_{}`$ (as well as $`l_+`$ is invariant to the transformations
$$\frac{}{(\pi \theta )}\frac{}{\theta },\frac{}{(\pi +\varphi )}\frac{}{\varphi }e^{i(\pi +\varphi )}e^{i\varphi }\mathrm{cot}(\pi \theta )\mathrm{cot}\theta .$$
(96)
In other words, $`l_\pm `$ are even. Therefore, we infer that the parity of any spherical harmonics is $`(1)^l`$, that is it is invariant under azimuthal changes
$$Y_{lm}(\varphi \theta ,\pi +\varphi )=(1)^lY_{lm}(\theta ,\varphi ).$$
(97)
In conclusion, the spherical harmonics are functions of well-defined parity, which is independent of $`m`$, even if $`l`$ is even and odd if $`l`$ is odd.
## The spin operator
Some particles have not only orbital angular momentum with regard to external axes but also a proper momentum , which is known as spin denoted here by $`\widehat{S}`$. This operator is not related to normal rotation with respect to ‘real’ axes in space, although it fulfills commutation relations of the same type as those of the orbital angular momentum, i.e.
$$[\widehat{S}_i,\widehat{S}_j]=i\epsilon _{ijk}\widehat{S}_k,$$
(98)
together with the following properties
* For the spin operator all the formulas of the orbital angular momentum from (23) till (48) are satisfied.
* The spectrum of the spin projections is a sequence of either integer or half-integer numbers differing by unity.
* The eigenvalues of $`\widehat{S}^2`$ are the following
$$\widehat{S}^2\psi _s=S(S+1)\psi _s.$$
(99)
* For a given $`S`$, the components $`S_z`$ can take only $`2S+1`$ values, from $`S`$ to $`+S`$.
* Besides the usual dependence on $`\stackrel{}{r}`$ and/or $`\stackrel{}{p}`$, the eigenfunctions of the particles with spin depend also on a discrete variable, (characteristic for the spin) $`\sigma `$ denoting the projection of the spin on the $`z`$ axis.
* The wavefunctions $`\psi (\stackrel{}{r},\sigma )`$ of a particle with spin can be expanded in eigenfunctions of given spin projection $`S_z`$, i.e.
$$\psi (\stackrel{}{r},\sigma )=\underset{\sigma =S}{\overset{S}{}}\psi _\sigma (\stackrel{}{r})\chi (\sigma ),$$
(100)
where $`\psi _\sigma (\stackrel{}{r})`$ is the orbital part and $`\chi (\sigma )`$ is the spinorial part.
* The spin functions (the spinors) $`\chi (\sigma _i)`$ are orhtogonal for any pair $`\sigma _i\sigma _k`$. The functions $`\psi _\sigma (\stackrel{}{r})\chi (\sigma )`$ in the sum of (100) are the components of a wavefunction of a particle with spin.
* The function $`\psi _\sigma (\stackrel{}{r})`$ is called the orbital part of the spinor, or shortly orbital.
* The normalization of the spinors is done as follows
$$\underset{\sigma =S}{\overset{S}{}}\psi _\sigma (\stackrel{}{r})=1.$$
(101)
The commutation relations allow to determine the explicit form of the spin operators (spin matrices) acting in the space of the eigenfunctions of definite spin projections.
Many ‘elementary’ particles, such as the electron, the neutron, the proton, etc. have a spin of $`1/2`$ (in units of $`\mathrm{}`$) and therefore the projection of their spin takes only two values, ($`S_z=\pm 1/2`$ (in $`\mathrm{}`$ units), respectively. They belong to the fermion class because of their statistics when they form many-body systems.
On the other hand, the matrices $`S_x,S_y,S_z`$ in the space of $`\widehat{S}^2,\widehat{S}_z`$ are
$`S_x={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),`$ $`S_y={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),`$ (106)
$`S_z={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),`$ $`S^2={\displaystyle \frac{3}{4}}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).`$ (111)
### Definition of the Pauli matrices
The matrices
$$\sigma _i=2S_i$$
(112)
are called the Pauli matrices. They are hermitic and have the same characteristic eq.
$$\lambda ^21=0.$$
(113)
Therefore, the eigenvalues of $`\sigma _x,\sigma _y`$ and $`\sigma _z`$ are
$$\lambda =\pm 1.$$
(114)
The algebra of these matrices is the following
$$\sigma _i^2=\widehat{I},\sigma _k\sigma _j=\sigma _j\sigma _k=i\sigma _z,\sigma _j\sigma _k=i\underset{l}{}\epsilon _{jkl}\sigma _l.+\delta _{jk}I.$$
(115)
In the case in which the spin system has spherical symmetry
$$\psi _1(r,+\frac{1}{2}),\psi _1(r,\frac{1}{2}),$$
(116)
are different solutions because of the different projections $`S_z`$. The value of the probability of one or another projection is determined by the square moduli $`\psi _1^2`$ or $`\psi _2^2`$, respectively, such that
$$\psi _1^2+\psi _2^2=1.$$
(117)
Since the eigenfunctions of $`S_z`$ have two components, then
$$\chi _1=\left(\begin{array}{c}1\\ 0\end{array}\right),\chi _2=\left(\begin{array}{c}0\\ 1\end{array}\right),$$
(118)
so that the eigenfunction of a spin one-half particle can be written as a column matrix
$$\psi =\psi _1\chi _1+\psi _2\chi _2=\left(\begin{array}{c}\psi _1\\ \psi _2\end{array}\right).$$
(119)
In the following, the orbitals will be replaced by numbers because we are interested only in the spin part.
## Transformations of spinors to rotations
Let $`\psi `$ be the wavefunction of a spin system in $`\mathrm{\Sigma }`$. We want to determine the probability of the spin projection in a arbitrary direction in 3D space, which one can always chose as the $`z^{}`$ of $`\mathrm{\Sigma }^{}`$. As we have already seen in the case of the angular momentum there are two viewpoints in trying to solve this problem
* $`\psi `$ does not change when $`\mathrm{\Sigma }\mathrm{\Sigma }^{}`$ and the operator $`\widehat{\mathrm{\Lambda }}`$ transforms as a vector. We have to find the eigenfunctions of the projections $`S_z^{}`$ and to expand $`\psi `$ in these eigenfunctions. The square moduli of the coefficients give the result
$`\widehat{S}_x^{}=\widehat{S}_x\mathrm{cos}\varphi +\widehat{S}_y\mathrm{sin}\varphi `$ $`=`$ $`e^{il\varphi }\widehat{S}_xe^{il\varphi },`$
$`\widehat{S}_y^{}=\widehat{S}_x\mathrm{sin}\varphi +\widehat{S}_y\mathrm{cos}\varphi `$ $`=`$ $`e^{il\varphi }\widehat{S}_ye^{il\varphi },`$
$`\widehat{S}_z^{}=\widehat{S}_z=e^{il\varphi }\widehat{S}_z,`$ (120)
for infinitesimal rotations. Then, from the commutation relations for spin one can find
$$\widehat{l}=\widehat{S}_z,$$
(121)
where $`\widehat{l}`$ is the infinitesimal generator.
* The second representation is:
$`\widehat{S}`$ does not change when $`\mathrm{\Sigma }\mathrm{\Sigma }^{}`$ and the components of $`\psi `$ does change. The transformation to this representation can be performed through a unitary transformation of the form
$`\widehat{V}^{}\widehat{S}^{}\widehat{V}`$ $`=`$ $`\widehat{\mathrm{\Lambda }},`$
$`\left(\begin{array}{c}\psi _1^{}\\ \psi _2^{}\end{array}\right)`$ $`=`$ $`\widehat{V}^{}\left(\begin{array}{c}\psi _1\\ \psi _2\end{array}\right).`$ (126)
Using (111) and (113) one gets
$`\widehat{V}^{}e^{i\widehat{S}_z\varphi }\widehat{S}e^{i\widehat{S}_z\varphi }\widehat{V}`$ $`=`$ $`\widehat{S},`$
$`\widehat{V}^{}`$ $`=`$ $`e^{i\widehat{S}_z\varphi },`$ (127)
and from (114) we are led to
$$\left(\begin{array}{c}\psi _1^{}\\ \psi _2^{}\end{array}\right)=e^{i\widehat{S}_z\varphi }\left(\begin{array}{c}\psi _1\\ \psi _2\end{array}\right).$$
(128)
Using the explicit form of $`\widehat{S}_z`$ and the properties of the Pauli matrices one can find the explicit form of $`\widehat{V}_z^{}`$, such that
$$\widehat{V}_z^{}(\varphi )=\left(\begin{array}{cc}e^{\frac{i}{2}\varphi }& 0\\ 0& e^{\frac{i}{2}\varphi }\end{array}\right).$$
(129)
## A result of Euler
One can reach any reference frame $`\mathrm{\Sigma }^{}`$ of arbitrary orientation with regard to $`\mathrm{\Sigma }`$ through only three rotations; the first of angle $`\varphi `$ around $`z`$, the next of angle $`\theta `$ around $`x^{}`$ and the last of angle $`\psi _a`$ around $`z^{}`$, i.e. This important result belongs to Euler. The parameters $`(\phi ,\theta ,\psi _a)`$ are called Euler’s angles. Thus
$$\widehat{V}^{}(\phi ,\theta ,\psi _a)=\widehat{V}_z^{}^{}(\psi _a)\widehat{V}_x^{}^{}(\theta )\widehat{V}_z^{}(\phi ).$$
(130)
The matrices $`\widehat{V}_z^{}`$ are of the form (116), whereas $`\widehat{V}_x^{}`$ is of the form
$$\widehat{V}_x^{}(\phi )=\left(\begin{array}{cc}\mathrm{cos}\frac{\theta }{2}& i\mathrm{sin}\frac{\theta }{2}\\ i\mathrm{sin}\frac{\theta }{2}& \mathrm{cos}\frac{\theta }{2}\end{array}\right),$$
(131)
so that
$$\widehat{V}^{}(\phi ,\theta ,\psi _a)=\left(\begin{array}{cc}e^{i\frac{\phi +\psi _a}{2}}\mathrm{cos}\frac{\theta }{2}& ie^{i\frac{\psi _a\phi }{2}}\mathrm{sin}\frac{\theta }{2}\\ ie^{i\frac{\phi \psi _a}{2}}\mathrm{sin}\frac{\theta }{2}& e^{i\frac{\phi +\psi _a}{2}}\mathrm{cos}\frac{\theta }{2}\end{array}\right).$$
(132)
It comes out in this way that by the rotation of $`\mathrm{\Sigma }`$, the components of the spinorial function transforms as follows
$`\psi _1^{}`$ $`=`$ $`\psi _1e^{i\frac{\phi +\psi _a}{2}}\mathrm{cos}{\displaystyle \frac{\theta }{2}}+i\psi _2e^{i\frac{\psi _a\phi }{2}}\mathrm{sin}{\displaystyle \frac{\theta }{2}},`$
$`\psi _2^{}`$ $`=`$ $`i\psi _1e^{i\frac{\phi \psi _a}{2}}\mathrm{sin}{\displaystyle \frac{\theta }{2}}+\psi _2e^{i\frac{\phi +\psi _a}{2}}\mathrm{cos}{\displaystyle \frac{\theta }{2}}.`$ (133)
From (120) one can infer that there is a one-to-one mapping between any rotation in $`E_3`$ and a linear transformation of $`E_2`$, the two-dimensional Euclidean space. This mapping is related to the two components of the spinorial wavefunction. The rotation in $`E_3`$ does not imply a rotation in $`E_2`$, which means that
$$\mathrm{\Phi }^{}|\psi ^{}=\mathrm{\Phi }|\psi =\mathrm{\Phi }_1^{}\psi _1+\mathrm{\Phi }_2^{}\psi _2.$$
(134)
From (119) one finds that (121) does not hold; nevertheless there is an invariance in the transformations (119) in the space $`E_2`$ of spinorial wavefunctions
$$\{\mathrm{\Phi }|\psi \}=\psi _1\mathrm{\Phi }_2\psi _2\mathrm{\Phi }_1.$$
(135)
The linear transformations that preserve invariant bilinear forms invariant are called binary transformations.
A physical quantity with two components for which a rotation of the coordinate system is a binary transformation is know as a spin of first order or shortly spin.
### The spinors of a system of two fermions
The eigenfunctions of $`{}_{i}{}^{}\widehat{s}_{i}^{2}\widehat{s}_z`$, with $`i=1,2`$ have the following form
$$i|+=\left(\begin{array}{c}1\\ 0\end{array}\right)_i,i|=\left(\begin{array}{c}0\\ 1\end{array}\right)_i.$$
(136)
A very used operator in a two-fermion system is the total spin
$$\widehat{S}=_1\widehat{S}+_2\widehat{S}.$$
(137)
The spinors of $`\widehat{s}^2\widehat{s}_z`$ are kets $`|\widehat{S},\sigma `$, which are linear combinations of $`{}_{i}{}^{}\widehat{s}_{i}^{2}\widehat{s}_z`$
$`|++=(\begin{array}{c}1\\ 0\end{array}\left)_1\right(\begin{array}{c}1\\ 0\end{array})_1,`$ $`|+=(\begin{array}{c}1\\ 0\end{array}\left)_1\right(\begin{array}{c}0\\ 1\end{array})_2,`$ (146)
$`|+=(\begin{array}{c}0\\ 1\end{array}\left)_2\right(\begin{array}{c}1\\ 0\end{array})_1,`$ $`|=(\begin{array}{c}0\\ 1\end{array}\left)_2\right(\begin{array}{c}0\\ 1\end{array})_2.`$ (155)
The spinorial functions in (125) are assumed orthonormalized. In $`E_n`$ the ket $`|++`$ has $`S_z=1`$ and at the same time it is an eigenfunction of the operator
$$\widehat{S}=_1\widehat{s}^2+2(_1\widehat{s})(_2\widehat{s})+_2\widehat{s}^2,$$
(156)
as one can see from
$`\widehat{S}^2`$ $`=`$ $`|++=\frac{3}{2}|+++2(_1\widehat{s}_x_2\widehat{s}_x+_1\widehat{s}_y_2\widehat{s}_y+_1\widehat{s}_z_2\widehat{s}_z)|++,`$ (157)
$`\widehat{S}^2`$ $`=`$ $`|++=2|++=1(1+1)|++.`$ (158)
If we introduce the operator
$$\widehat{S}_{}=_1\widehat{s}_{}+_2\widehat{s}_{},$$
(159)
one gets
$$[\widehat{S}_{},\widehat{S}^2]=0.$$
(160)
Then $`(\widehat{S}_{})^k|1,1`$ can be written in terms of the eigenfunctions of the operator $`\widehat{S}^2`$, i.e.
$$\widehat{S}_{}|1,1=\widehat{S}_{}|++=\sqrt{2}|++\sqrt{2}|+.$$
(161)
Thus, $`S_z=0`$ in the state $`\widehat{S}_{}|1,1`$. On the other hand, from the normalization condition, we have
$`|1,0=\frac{1}{\sqrt{2}}(|++|+)`$ (162)
$`\widehat{S}_{}|1,0=|+|=\alpha |1,1.`$ (163)
In addition, the normalization condition gives
$$|1,1=|,.$$
(164)
There is only one other linear-independent combination of functions of the type (125), which is different of $`|1,1,|1,0`$ and $`|1,1`$, which is
$`\psi _4=\frac{1}{\sqrt{2}}(|+|+),`$ (165)
$`\widehat{S}_z\psi _4=0,\widehat{S}^2\psi _4.`$ (166)
Therefore
$$\psi _4=|0,0.$$
(167)
$`\psi _4`$ describes the state of a system of two fermions having the total spin equal to zero. The latter type of state is called singlet. On the other hand, the state of two fermions of total spin one can be called triplet having a degree of degeneration $`g=3`$.
## Total angular momentum
The total angular momentum is an operator defined as the sum of the angular and spin momenta, i.e.
$$\widehat{J}=\widehat{l}+\widehat{S},$$
(168)
where $`\widehat{l}`$ and $`\widehat{S}`$, as we have seen, act in different spaces, though the square of $`\widehat{l}`$ and $`\widehat{S}`$ commute with $`\widehat{J}`$
$$[\widehat{J}_i,\widehat{J}_j]=i\epsilon _{ijk}\widehat{J}_k,[\widehat{J}_i,\widehat{l}^2]=0,[\widehat{J}_i,\widehat{S}^2]=0,$$
(169)
From (139) one finds that $`\widehat{l}^2`$ and $`\widehat{S}^2`$ have a common eigenfunction system with $`\widehat{J}^2`$ and $`\widehat{J}_z`$.
Let us determine the spectrum of the projections $`\widehat{J}_z`$ for a fermion. The state of maximum projection $`\widehat{J}_z`$ can be written
$`\overline{\psi }`$ $`=`$ $`\psi _{ll}\left(\begin{array}{c}1\\ 0\end{array}\right)=|l,l,+`$ (172)
$`\widehat{ȷ}_z\psi `$ $`=`$ $`(l+\frac{1}{2})\overline{\psi },j=l+\frac{1}{2}.`$ (173)
We introduce the operator $`\widehat{J}_{}`$ defined as
$$\widehat{J}_{}=\widehat{l}_{}+\widehat{S}_{}=\widehat{l}_{}+\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).$$
(174)
On account of the normalization $`\alpha =\sqrt{(J+M)(JM+1)}`$, one gets
$$\widehat{J}_{}\psi _{ll}\left(\begin{array}{c}1\\ 0\end{array}\right)=\sqrt{2l}|l,l1,++|l,l1,,$$
(175)
so that the value of the projection of $`\widehat{j}_{}`$ in $`\widehat{j}_{}\overline{\psi }`$ will be
$$\widehat{ȷ}_z=(l1)+\frac{1}{2}=l\frac{1}{2}.$$
(176)
It follows that $`\widehat{ȷ}_{}`$ lowers by one unit the action of $`\widehat{J}_z`$.
In the general case we have
$$\widehat{ȷ}_{}^k=\widehat{l}_{}^k+k\widehat{l}_{}^{k1}\widehat{S}_{}.$$
(177)
One can see that (145) is obtained from the binomial expansion considering that $`\widehat{s}_{}^2`$ and all higher-order powers of $`\widehat{s}`$ are zero.
$$\widehat{ȷ}_{}^k|l,l,+=\widehat{l}_{}^k|l,l,++k\widehat{l}_{}^{k1}|l,l,.$$
(178)
Using
$$(\widehat{l}_{})^k\psi _{l,l}=\sqrt{\frac{k!(2l)!}{(2lk)!}}\psi _{l,lk}$$
we get
$$\widehat{ȷ}_{}^k|l,l,+=\sqrt{\frac{k!(2l)!}{(2lk)!}}|l,lk,++\sqrt{\frac{(k+1)!(2l)!}{(2lk+1)!}}k|l,lk+1,.$$
(179)
Now noticing that $`m=lk`$
$$\widehat{ȷ}_{}^{lm}|l,l,+=\sqrt{\frac{(lm)!(2l)!}{(l+m)!}}|l,m,++\sqrt{\frac{(lm1)!(2l)!}{(2l+m+1)!}}(lm)|l,m+1,.$$
(180)
The eigenvalues of the projections of the total angular momentum are given by the sequence of numbers differing by one unit from $`j=l+\frac{1}{2}`$ pînă to $`j=l\frac{1}{2}`$. All these states belong to the same eigenfunction of $`\widehat{J}`$ as $`|l,l,+`$ because $`[\widehat{J}_{},\widehat{J}^2]=0`$:
$`\widehat{J}^2|l,l,+`$ $`=`$ $`(\widehat{l}^2+2\widehat{l}\widehat{S}+\widehat{S}^2)|l,l,+,`$ (181)
$`=`$ $`[l(l+1)+2l\frac{1}{2}+\frac{3}{4}]|l,l,+`$
where $`j(j+1)=(l+\frac{1}{2})(l+\frac{3}{2})`$.
In the left hand side of (149) a contribution different of zero gives only $`j=\widehat{l}_z\widehat{S}_z`$. Thus, the obtained eigenfunctions correspond to the pair $`j=l+\frac{1}{2}`$, $`m_j=m+\frac{1}{2}`$; they are of the form
$$|l+\frac{1}{2},m+\frac{1}{2}=\sqrt{\frac{l+m+1}{2l+1}}|l,m,++\sqrt{\frac{lm}{2l+1}}|l,m+1,.$$
(182)
The total number of linearly independent states is
$$N=(2l+1)(2s+1)=4l+2,$$
(183)
of which in (150) only (2j+1)=2l+3 have been built. The rest of $`2l1`$ eigenfunctions can be obtained from the orthonormalization condition:
$$|l\frac{1}{2},m\frac{1}{2}=\sqrt{\frac{lm}{2l+1}}|l,m,+\sqrt{\frac{l+m+1}{2l+1}}|l,m+1,.$$
(184)
If two subsystems are in interaction in such a way that each of the angular momenta $`\widehat{j}_i`$ is conserved, then the eigenfunctions of the total angular momentum
$$\widehat{J}=\widehat{ȷ}_1+\widehat{ȷ}_2,$$
(185)
can be obtained by a procedure similar to the previous one. For fixed eigenvalues of $`\widehat{ȷ}_1`$ and $`\widehat{ȷ}_2`$ there are $`(2j_1+1)(2j_2+1)`$ orthonormalized eigenfunctions of the projection of the total angular momentum $`\widehat{J}_z`$; the one corresponding to the maximum value of the projection $`\widehat{J}_z`$, i.e. $`M_J=j_1+j_2`$, can be built in a unique way and therefore $`J=j_1+j_2`$ is the maximum value of the total angular momentum of the system. Applying the operator $`\widehat{J}=\widehat{ȷ}_1+\widehat{ȷ}_2`$ repeatingly to the function
$$|j_1+j_2,j_1+j_2,j_1+j_2=|j_1,j_1|j_2,j_2,$$
(186)
one can obtain all the $`2(j_1+j_2)+1`$ eigenfunctions of $`\widehat{J}=j_1+j_2`$ with different $`M`$s:
$$(j_1+j_2)M(j_1+j_2).$$
For example, the eigenfunction of $`M=j_1+j_21`$ is
$$|j_1+j_2,j_1+j_21,j_1,j_2=\sqrt{\frac{j_1}{j_1+j_2}}|j_1,j_11,j_2,j_2+\sqrt{\frac{j_2}{j_1+j_2}}|j_1,j_1,j_2,j_21.$$
(187)
Applying iteratively the operator $`\widehat{J}_{}`$, all the $`2(j_1+j_21)1`$ eigenfunctions of $`J=j_1+j_21`$ can be obtained.
One can prove that
$$|j_1j_2|Jj_1+j_2,$$
so that
$$\underset{\mathrm{min}J}{\overset{\mathrm{max}J}{}}(2J+1)=(2J_1+1)(2J_2+1).$$
(188)
Thus
$$|J,M,j_1,j_2=\underset{m_1+m_2=M}{}(j_1m_1j_2m_2|JM)|j_1,m_1,j_2,m_2,$$
(189)
where the coefficients $`(j_1m_1j_2m_2|JM)`$ determine the contribution of the various kets $`|j_1,m_1,j_2,m_2`$ to the eigenfunctions of $`\widehat{J^2}`$, $`\widehat{J_z}`$ having the eigenvalues $`J(J+1)`$, $`M`$. They are called Clebsch-Gordan coefficients.
References:
1. H.A. Buchdahl, “Remark concerning the eigenvalues of orbital angular momentum”,
Am. J. Phys. 30, 829-831 (1962)
3N. Note: 1. The operator corresponding to the Runge-Lenz vector of the classical Kepler problem is written as
$$\widehat{\stackrel{}{A}}=\frac{\widehat{𝐫}}{r}+\frac{1}{2}\left[(\widehat{l}\times \widehat{p})(\widehat{p}\times \widehat{l})\right],$$
where atomic units have been used and the case $`Z=1`$ (hydrogen atom) was assumed. This operator commutes with the Hamiltonian of the atomic hydrogen $`\widehat{H}=\frac{\widehat{p^2}}{2}\frac{1}{r}`$, that is it is an integral of the atomic quantum motion. Its components have commutators of the type $`[A_i,A_j]=2iϵ_{ijk}l_kH`$; the commutators of the Runge-Lenz components with the components of the angular momentum are of the type $`[l_i,A_j]=iϵ_{ijk}A_k`$. Thus, they respect the conditions (23). Proving that can be a useful exercise.
## 3P. Problems
### Problem 3.1
Show that any translation operator, for which $`\psi (y+a)=T_a\psi (y)`$, can be written as an exponential operator. Apply the result for $`y=\stackrel{}{r}`$ and for a the finite rotation $`\alpha `$ around $`z`$.
Solution
The proof can be obtained expanding $`\psi (y+a))`$ in Taylor series in the infinitesimal neighborhood around $`x`$, that is in powers of $`a`$
$$\psi (y+a)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{a^n}{n!}\frac{d^n}{dx^n}\psi (x)$$
We notice that
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{a^n\frac{d^n}{dx^n}}{n!}=e^{a\frac{d}{dx}}$$
and therefore one has $`T_a=e^{a\frac{d}{dx}}`$ in the 1D case. In 3D, $`y=\stackrel{}{r}`$ and $`a\stackrel{}{a}`$. The result is $`T_\stackrel{}{a}=e^\stackrel{}{a}\stackrel{}{}`$.
For the finite rotation $`\alpha `$ around $`z`$ we has $`y=\varphi `$ and $`a=\alpha `$. It follows
$$T_\alpha =R_\alpha =e^{\alpha \frac{d}{d\varphi }}.$$
Another exponential form of the rotation around $`z`$ is that in terms of the angular momentum operator as was already commented in this chapter. Let $`x^{}=x+dx`$ and consider only the first order of the Taylor series
$`\psi (x^{},y^{},z^{})`$ $`=`$ $`\psi (x,y,z)+(x^{}x){\displaystyle \frac{}{x^{}}}\psi (x^{},y^{},z^{})|_{\stackrel{}{r^{}}=\stackrel{}{r}}`$
$`+(y^{}y){\displaystyle \frac{}{y^{}}}\psi (x^{},y^{},z^{})|_{\stackrel{}{r^{}}=\stackrel{}{r}}`$
$`+(z^{}z){\displaystyle \frac{}{z^{}}}\psi (x^{},y^{},z^{})|_{\stackrel{}{r^{}}=\stackrel{}{r}}.`$
Taking into account
$`{\displaystyle \frac{}{x_i^{}}}\psi (\stackrel{}{r}^{})|_\stackrel{}{r}^{}`$ $`=`$ $`{\displaystyle \frac{}{x_i}}\psi (\stackrel{}{r}),`$
$`x^{}=xyd\varphi ,y^{}`$ $`=`$ $`y+xd\varphi ,z^{}=z,`$
one can reduce the series from three to two dimensions
$`\psi (\stackrel{}{r}^{})`$ $`=`$ $`\psi (\stackrel{}{r})+(xyd\varphi x){\displaystyle \frac{\psi (\stackrel{}{r})}{x}}+(y+xd\varphi y){\displaystyle \frac{\psi (\stackrel{}{r})}{y^{}}},`$
$`=`$ $`\psi (\stackrel{}{r})yd\varphi {\displaystyle \frac{\psi (\stackrel{}{r})}{x}}+xd\varphi x{\displaystyle \frac{\psi (\stackrel{}{r})}{y}},`$
$`=`$ $`\left[1d\varphi \left(x{\displaystyle \frac{}{y}}+y{\displaystyle \frac{}{x}}\right)\right]\psi (\stackrel{}{r}).`$
Since $`i\widehat{l}_z=\left(x\frac{}{y}y\frac{}{x}\right)`$ it follows that $`R=\left[1d\varphi \left(x\frac{}{y}y\frac{}{x}\right)\right].`$ In the second order one can get $`\frac{1}{2!}(i\widehat{l}_zd\varphi )^2`$, and so forth. Thus, $`R`$ can be written as an exponential
$$R=e^{i\widehat{l}_zd\varphi }.$$
### Problem 3.2
Based on the expressions given in (14) show that one can get (15).
Solution
Let us consider only linear terms in the Taylor expansion (infinitesimal rotations)
$$e^{i\widehat{l}_zd\varphi }=1+i\widehat{l}_zd\varphi +\frac{1}{2!}(i\widehat{l}_zd\varphi )^2+\mathrm{},$$
so that
$`(1+i\widehat{l}_zd\varphi )\widehat{A}_x(1i\widehat{l}_zd\varphi )`$ $`=`$ $`\widehat{A}_x\widehat{A}_xd\varphi ,`$
$`(\widehat{A}_x+i\widehat{l}_zd\varphi \widehat{A}_x)(1i\widehat{l}_zd\varphi )`$ $`=`$ $`\widehat{A}_x\widehat{A}_xd\varphi ,`$
$`\widehat{A}_x\widehat{A}_xi\widehat{l}_zd\varphi +i\widehat{l}_zd\varphi \widehat{A}_x+\widehat{l}_zd\varphi \widehat{A}_x\widehat{l}_zd\varphi `$ $`=`$ $`\widehat{A}_x\widehat{A}_xd\varphi ,`$
$`i(\widehat{l}_z\widehat{A}_x\widehat{A}_x\widehat{l}_z)d\varphi `$ $`=`$ $`\widehat{A}_yd\varphi .`$
We easily arrive at the conclusion
$$[\widehat{l}_z,\widehat{A}_x]=i\widehat{A}_y.$$
In addition, $`[\widehat{l}_z,\widehat{A}_y]=i\widehat{A}_x`$ can be obtained from
$`(1+i\widehat{l}_zd\varphi )\widehat{A}_y(1i\widehat{l}_zd\varphi )`$ $`=`$ $`\widehat{A}_xd\varphi \widehat{A}_y,`$
$`(\widehat{A}_y+i\widehat{l}_zd\varphi \widehat{A}_y)(1i\widehat{l}_zd\varphi )`$ $`=`$ $`\widehat{A}_xd\varphi \widehat{A}_y,`$
$`\widehat{A}_y\widehat{A}_yi\widehat{l}_zd\varphi +i\widehat{l}_zd\varphi \widehat{A}_y+\widehat{l}_zd\varphi \widehat{A}_y\widehat{l}_zd\varphi `$ $`=`$ $`\widehat{A}_xd\varphi \widehat{A}_y,`$
$`i(\widehat{l}_z\widehat{A}_y\widehat{A}_y\widehat{l}_z)d\varphi `$ $`=`$ $`\widehat{A}_xd\varphi .`$
### Problem 3.3
Determine the operator $`\frac{d\widehat{\sigma }_x}{dt}`$ based on the Hamiltonian of an electron with spin in a magnetic field of induction $`\stackrel{}{B}`$.
Solution
The Hamiltonian in this case is $`\widehat{H}(\widehat{𝐩},\widehat{𝐫},\widehat{\sigma })=\widehat{H}(\widehat{𝐩},\widehat{𝐫})+\widehat{\sigma }\stackrel{}{𝐁}`$, where the latter term is the Zeeman Hamiltonian of the electron. Since $`\widehat{\sigma }_x`$ commutes with the momenta and the coordinates, applying the Heisenberg equation of motion leads to
$$\frac{d\widehat{\sigma }_x}{dt}=\frac{i}{\mathrm{}}[\widehat{H},\widehat{\sigma }_x]=\frac{i}{\mathrm{}}\frac{e\mathrm{}}{2m_e}((\widehat{\sigma }_yB_y+\widehat{\sigma }_zB_z)\widehat{\sigma }_x\widehat{\sigma }_x(\widehat{\sigma }_yB_y+\widehat{\sigma }_zB_z)).$$
Using $`[\sigma _x,\sigma _y]=i\sigma _z`$, one gets :
$$\frac{d\widehat{\sigma }_x}{dt}=\frac{e}{m_e}(\widehat{\sigma }_yB_z\widehat{\sigma }_zB_y)=\frac{e}{m_e}(\stackrel{}{\sigma }\times \stackrel{}{B})_x.$$
4. THE WKB METHOD
In order to study more realistic potentials with regard to rectangular barriers and wells, it is necessary to employ approximate methods allowing to solve the Schrödinger equation for more general classes of potentials and at the same time to give very good approximations of the exact solutions.
The aim of the various approximative methods is to offer solutions of acceptable precision and simplicity that can be used for understanding the behaviour of the system in quasianalytic terms.
Within quantum mechanics, one of the oldest and efficient approximate method for getting rather good Schrödinger solutions was developed almost simulataneously by G. Wentzel, H. A. Kramers and L. Brillouin in 1926, hence the acronym WKB under which this method is known (or JWKB as is more correctly used by many authors, see note 4N).
It is worth mentioning that the WKB method applies to 1D Schrödinger equations and that there are serious difficulties when trying to generalize it to more dimensions.
In order to solve the Schrödinger equation
$$\frac{\mathrm{}^2}{2m}\frac{d^2\psi }{dy^2}+u(y)\psi =E\psi $$
(1)
with a potential of the form
$$u(y)=u_0f\left(\frac{y}{a}\right),$$
(2)
we first perform the changes of notations and of variable
$$\xi ^2=\frac{\mathrm{}^2}{2mu_0a^2}$$
(3)
$$\eta =\frac{E}{u_0}$$
(4)
$$x=\frac{y}{a}.$$
(5)
From eq. $`(5)`$ we get
$$\frac{d}{dx}=\frac{dy}{dx}\frac{d}{dy}=a\frac{d}{dy}$$
(6)
$$\frac{d^2}{dx^2}=\frac{d}{dx}\left(a\frac{d}{dy}\right)=\left(a\frac{d}{dx}\right)\left(a\frac{d}{dx}\right)=a^2\frac{d^2}{dy^2}$$
(7)
and the Schrödinger eq. reads
$$\xi ^2\frac{d^2\psi }{dx^2}+f(x)\psi =\eta \psi .$$
(8)
Multiplying by $`1/\xi ^2`$ and defining $`r(x)=\eta f(x)`$, it is possible to write it as folows
$$\frac{d^2\psi }{dx^2}+\frac{1}{\xi ^2}r(x)\psi =0.$$
(9)
To solve (9), the following form of the solution is proposed
$$\psi (x)=\mathrm{exp}\left[\frac{i}{\xi }_a^xq(x)𝑑x\right].$$
(10)
Therefore
$$\frac{d^2\psi }{dx^2}=\frac{d}{dx}\left(\frac{d\psi }{x}\right)=\frac{d}{dx}\left\{\frac{i}{\xi }q(x)\mathrm{exp}\left[\frac{i}{\xi }_a^xq(x)𝑑x\right]\right\}$$
$$\frac{d^2\psi }{dx^2}=\frac{i}{\xi }\left\{\frac{i}{\xi }q^2(x)\mathrm{exp}\left[\frac{i}{\xi }_a^xq(x)𝑑x\right]+\frac{q(x)}{x}\mathrm{exp}\left[\frac{i}{\xi }_a^xq(x)𝑑x\right]\right\}.$$
Factorizing $`\psi `$, we have
$$\frac{d^2\psi }{dx^2}=\left[\frac{1}{\xi ^2}q^2(x)+\frac{i}{\xi }\frac{dq(x)}{dx}\right]\psi .$$
(11)
Discarding for the time being the dependence of $`x`$, the Schrödinger eq. can be written
$$\left[\frac{1}{\xi ^2}q^2+\frac{i}{\xi }\frac{q}{x}+\frac{1}{\xi ^2}r\right]\psi =0$$
(12)
and since in general $`\psi 0`$, we get:
$$i\xi \frac{dq}{dx}+rq^2=0,$$
(13)
which is a nonlinear differential eq. of the Riccati type whose solutions are sought in the form of expansions in powers of $`\xi `$ under the assumption that $`\xi `$ is very small.
More precisely, the series is taken of the form
$$q(x)=\underset{n=0}{\overset{\mathrm{}}{}}(i\xi )^nq_n(x).$$
(14)
Plugging it into the Riccati eq., we get
$$i\xi \underset{n=0}{\overset{\mathrm{}}{}}(i\xi )^n\frac{dq_n}{dx}+r(x)\underset{\mu =0}{\overset{\mathrm{}}{}}(i\xi )^\mu q_\mu \underset{\nu =0}{\overset{\mathrm{}}{}}(i\xi )^\nu q_\nu =0.$$
(15)
By a rearrangement of the terms one is led to
$$\underset{n=0}{\overset{\mathrm{}}{}}(1)^n(i\xi )^{n+1}\frac{dq_n}{dx}+r(x)\underset{\mu =0}{\overset{\mathrm{}}{}}\underset{\nu =0}{\overset{\mathrm{}}{}}(i\xi )^{\mu +\nu }q_\mu q_\nu =0.$$
(16)
Double series have the following important property
$$\underset{\mu =0}{\overset{\mathrm{}}{}}\underset{\nu =0}{\overset{\mathrm{}}{}}a_{\mu \nu }=\underset{n=0}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{n}{}}a_{m,nm},$$
where $`\mu =nm,\nu =m`$ .
Thus
$$\underset{n=0}{\overset{\mathrm{}}{}}(1)^n(i\xi )^{n+1}\frac{dq_n}{dx}+r(x)\underset{n=0}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{n}{}}(i\xi )^{nm+m}q_mq_{nm}=0.$$
(17)
Let us see explicitly the first several terms in each of the series in eq. (17):
$$\underset{n=0}{\overset{\mathrm{}}{}}(1)^n(i\xi )^{n+1}\frac{dq_n}{dx}=i\xi \frac{dq_0}{dx}+\xi ^2\frac{dq_1}{dx}i\xi ^3\frac{dq_2}{dx}+\mathrm{}$$
(18)
$$\underset{n=0}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{n}{}}(i\xi )^nq_mq_{nm}=q_0^2i2\xi q_0q_1+\mathrm{}$$
(19)
Asking that the first terms in both series contain $`i\xi `$, one should write them as
$$\underset{n=1}{\overset{\mathrm{}}{}}(1)^{n1}(i\xi )^n\frac{dq_{n1}}{dx}+r(x)q_0^2\underset{n=1}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{n}{}}(i\xi )^nq_mqnm=0,$$
which leads to
$$\underset{n=1}{\overset{\mathrm{}}{}}\left[(i\xi )^n\frac{dq_{n1}}{dx}\underset{m=0}{\overset{n}{}}(i\xi )^nq_mq_{nm}\right]+\left[r(x)q_0^2\right]=0.$$
(20)
In order that this equation be right the following conditions should be satisfied
$$r(x)q_0^2=0q_0=\pm \sqrt{r(x)}$$
(21)
$$(i\xi )^n\frac{dq_{n1}}{dx}\underset{m=0}{\overset{n}{}}(i\xi )^nq_mq_{nm}=0$$
$$\frac{dq_{n1}}{dx}=\underset{m=0}{\overset{n}{}}q_mq_{nm}n1.$$
(22)
The latter is a recurrence relatioship, which occurs naturally in the WKB method. Recalling that we have defined $`r(x)=\eta f(x),\eta =\frac{E}{u_0}\&f(x)=\frac{u}{u_0}`$, by means of eq. $`(21)`$ we get
$$q_0=\pm \sqrt{\eta f(x)}=\pm \sqrt{\frac{E}{u_0}\frac{u}{u_0}}=\pm \sqrt{\frac{2m(Eu)}{2mu_0}}.$$
(23)
This clearly indicates the classical nature of the WKB momentum of the particle of energy $`E`$ in the potential $`u`$ and units of $`\sqrt{2mu_0}`$. Thus
$$q_0=p(x)=\sqrt{\eta f(x)}$$
is not an operator. If we approximate till the second order, we get
$$q(x)=q_0i\xi q_1\xi ^2q_2$$
and using the WKB recurrence relationship (22) we calculate $`q_1`$ and $`q_2`$
$$\frac{dq_0}{dx}=2q_0q_1q_1=\frac{1}{2}\frac{\frac{dq_0}{dx}}{q_0}=\frac{1}{2}\frac{d}{dx}(\mathrm{ln}|q_0|)$$
$$q_1=\frac{1}{2}\frac{d}{dx}(\mathrm{ln}|p(x)|)$$
(24)
$$\frac{dq_1}{dx}=2q_0q_2q_1^2q_2=\frac{\frac{dq_1}{dx}q_1^2}{2q_0}.$$
(25)
A glance to eq. $`(24)`$, affords us to consider $`q_1`$ as the slope, up to a change of sign, of $`\mathrm{ln}|q_0|`$; when $`q_0`$ is very small, then $`q_10\xi q_10`$ and therefore the series diverges. To avoid this the following WKB condition is imposed
$$|q_0||\xi q_1|=\xi |q_1|.$$
It is worth noting that this WKB condition WKB is not fulfilled at those points $`x_k`$ where
$$q_0(x_k)=p(x_k)=0.$$
Since $`q_0=p=\sqrt{\frac{2m(Eu)}{2mu_0}}`$ the previous equation leads us to
$$E=u(x_k).$$
(26)
In classical mechanics the points $`x_k`$ that satisfies (26) are called turning points because the change of the sense of the motion of a macroscopic particle takes place there.
By means of these arguments, we can say that $`q_0`$ is a classical solution of the problem under examination; also that the quantities $`q_1`$ & $`q_2`$ are the first and the second quantum corrections, respectively, in the WKB problem.
To obtain the WKB wavefunctions we shall consider only the classical solution and the first quantum correction that we plug in the WKB form of $`\psi `$
$$\psi =\mathrm{exp}\left[\frac{i}{\xi }_a^xq(x)𝑑x\right]=\mathrm{exp}\left[\frac{i}{\xi }_a^x(q_0i\xi q_1)𝑑x\right]$$
$$\psi =\mathrm{exp}\left(\frac{i}{\xi }_a^xq_0𝑑x\right)\mathrm{exp}\left(_a^xq_1𝑑x\right).$$
For the second factor, we get
$$\mathrm{exp}\left(_a^xq_1𝑑x\right)=\mathrm{exp}\left[\frac{1}{2}_a^x\frac{d}{dx}(\mathrm{ln}|p(x)|)𝑑x\right]=$$
$$=\mathrm{exp}\left[\frac{1}{2}(\mathrm{ln}|p(x)|)|_a^x\right]=\frac{A}{\sqrt{p(x)}},$$
where $`A`$ is a constant, whereas for the first factor we get
$$\mathrm{exp}\left(\frac{i}{\xi }_a^xq_0𝑑x\right)=\mathrm{exp}\left[\pm \frac{i}{\xi }_a^xp(x)𝑑x\right].$$
Thus, we can write $`\psi `$ in the following form
$$\psi ^\pm =\frac{1}{\sqrt{p(x)}}\mathrm{exp}\left[\pm \frac{i}{\xi }_a^xp(x)𝑑x\right].$$
(27)
The latter are known as the WKB solutions of the 1D Schrödinger equation. The general WKB solution in the region in which the WKB condition is satisfied is written down as
$$\psi =a_+\psi ^++a_{}\psi ^{}.$$
(28)
As already mentioned there is no WKB solution at the turning points. This raises the question of the manner in which one has to do the passing from $`\psi (x<x_k)`$ to $`\psi (x>x_k)`$. The solution of this difficulty is achieved by introducing the WKB connection formulas.
### The connection formulas
We have already seen that the WKB solutions are singular at the classical turning points; however, these solutions are correct both on the left and right side of these turning points $`x_k`$. A natural question is how do we change $`\psi (x<x_k)`$ in $`\psi (x>x_k)`$ when passing through the turning points. The explicit answer is given by the connection formulas.
From the theory of differential equations of complex variable it can be proved that really there are such connection formulas and that they are the following
$$\psi _1(x)=\frac{1}{\left[r(x)\right]^{\frac{1}{4}}}\mathrm{exp}\left(_x^{x_k}\sqrt{r(x)}𝑑x\right)$$
$$\frac{2}{\left[r(x)\right]^{\frac{1}{4}}}\mathrm{cos}\left(_{x_k}^x\sqrt{r(x)}𝑑x\frac{\pi }{4}\right),$$
(29)
where $`\psi _1(x)`$ has only an attenuated exponential behavior for $`x<x_k`$. The first connection formula shows that the function $`\psi (x)`$, which at the left of the turning point behaves exponentially decaying, turns at the right of $`x_k`$ into a cosinusoide of phase $`\varphi =\frac{\pi }{4}`$ and double amplitude with regard to the amplitude of the exponential.
In the case of a more general function $`\psi (x)`$, such as a function with both rising and decaying exponential behavior, the connection formula is
$$\mathrm{sin}\left(\varphi +\frac{\pi }{4}\right)\frac{1}{\left[r(x)\right]^{\frac{1}{4}}}\mathrm{exp}\left(_x^{x_k}\sqrt{r(x)}𝑑x\right)$$
$$\frac{1}{\left[r(x)\right]^{\frac{1}{4}}}\mathrm{cos}\left(_{x_k}^x\sqrt{r(x)}𝑑x+\varphi \right),$$
(30)
under the condition that $`\varphi `$ să do not take a value that is too close to $`\frac{\pi }{4}`$. The reason is that if $`\varphi =\frac{\pi }{4}`$, then the sinus function is zero . The latter connection formula means that a function whose behavior is of the cosinusoid type at the right of a turning point changes into a growing exponential with sinusoid-modulated amplitude at the right of that point.
In order to study the details of the procedure of getting the connection formulas we recommend the book Mathematical Methods of Physics by J. Mathews & R.L. Walker.
### Estimation of the WKB error
We have found the solution of the Schrödinger equation in the regions where the WKB condition is satisfied. However, the WKB solutions are divergent at the turning points. We thus briefly analyze the error introduced by using the WKB approximation and tackling the connection formulas in a close neighbourhood of the turning points.
Considering $`x=x_k`$ as a turning point, we have $`q_0(x_k)=p(x_k)=0E=u(x_k)`$. At the left of $`x_k`$, that is on the ‘half-line’ $`x<x_k`$, we shall assume $`E<u(x)`$ leading to the WKB solution
$$\psi (x)=\frac{a}{\left[\frac{u(x)E}{u_0}\right]^{\frac{1}{4}}}\mathrm{exp}\left(\frac{1}{\xi }_x^{x_k}\sqrt{\frac{u(x)E}{u_0}}𝑑x\right)+$$
$$+\frac{b}{\left[\frac{u(x)E}{u_0}\right]^{\frac{1}{4}}}\mathrm{exp}\left(\frac{1}{\xi }_x^{x_k}\sqrt{\frac{u(x)E}{u_0}}𝑑x\right).$$
(31)
Similarly, at the right of $`x_k`$ (on the ‘half-line $`x>x_k`$) we assume $`E>u(x)`$; therefore the WKB solution in the latter region will be
$$\psi (x)=\frac{c}{\left[\frac{Eu(x)}{u_0}\right]^{\frac{1}{4}}}\mathrm{exp}\left(\frac{i}{\xi }_{x_k}^x\sqrt{\frac{Eu(x)}{u_0}}𝑑x\right)+$$
$$+\frac{d}{\left[\frac{Eu(x)}{u_0}\right]^{\frac{1}{4}}}\mathrm{exp}\left(\frac{i}{\xi }_{x_k}^x\sqrt{\frac{Eu(x)}{u_0}}𝑑x\right).$$
(32)
If $`\psi (x)`$ is a real function, it will have this property both at the right and the left of $`x_k`$. It is usually called the “reality condition”. It means that if $`a,b\mathrm{}`$, then $`c=d^{}`$.
Our problem consists in connecting the approximations on the two sides of $`x_k`$ such that they refer to the same solution. This means to find $`c`$ and $`d`$ if one knows $`a`$ and $`b`$, as well as viceversa. To achieve this connection, we have to use an approximate solution, which should be correct along a contour connecting the regions on the two sides of $`x_k`$, where the WKB solutions are also correct. A method proposed by Zwann and Kemble is very useful in this case. It consists in going out from the real axis in the neighbourhood of $`x_k`$ on a contour around $`x_k`$ in the complex plane. It is assumed that on this contour the WKB solutions are still correct. Here, we shall use this method as a means of getting the estimation of the error produced by the WKB method.
The estimation of the error is always an important matter for any approximate solutions. In the case of the WKB procedure, it is more significant because it is an approximation on large intervals of the real axis that can lead to the accuulation of the errors as well as to some artefacts due to the phase shifts that can be introduced in this way.
Let us define the associated WKB functions as follows
$$W_\pm =\frac{1}{\left[\frac{Eu(x)}{u_0}\right]^{\frac{1}{4}}}\mathrm{exp}\left(\pm \frac{i}{\xi }_{x_k}^x\sqrt{\frac{Eu(x)}{u_0}}𝑑x\right),$$
(33)
that we consider as functions of complex variable. We shall use cuts in order to avoid the discontinuities in the zeros of $`r(x)=\frac{Eu(x)}{u_0}`$. These functions satisfy the differential equation that is obtained by differentiating with respect to $`x`$, leading to
$$W_\pm ^{}=\left(\pm \frac{i}{\xi }\sqrt{r}\frac{1}{4}\frac{r^{}}{r}\right)W_\pm $$
$$W_\pm ^{\prime \prime }+\left[\frac{r}{\xi ^2}+\frac{1}{4}\frac{r^{\prime \prime }}{r}\frac{5}{16}\left(\frac{r^{}}{r}\right)^2\right]W_\pm =0.$$
(34)
Let us notice that
$$s(x)=\frac{1}{4}\frac{r^{\prime \prime }}{r}\frac{5}{16}\left(\frac{r^{}}{r}\right)^2,$$
(35)
then $`W_\pm `$ are exact solutions of the equation
$$W_\pm ^{\prime \prime }+\left[\frac{1}{\xi ^2}r(x)+s(x)\right]W_\pm =0,$$
(36)
although they satisfy only approximately the Schrödinger equation, which is a regular equation in $`x=x_k`$, whereas the same equation for the associate WKB functions is singular at that point.
We shall now define the functions $`\alpha _\pm (x)`$ satisfying the following two relationships
$$\psi (x)=\alpha _+(x)W_+(x)+\alpha _{}(x)W_{}(x)$$
(37)
$$\psi ^{}(x)=\alpha _+(x)W_+^{}(x)+\alpha _{}(x)W_{}^{}(x),$$
(38)
where $`\psi (x)`$ is a solution of the Schrödinger equation. Solving the previous equations for $`\alpha _\pm `$, we get
$$\alpha _+=\frac{\psi W_{}^{}\psi ^{}W_{}}{W_+W_{}^{}W_+^{}W_{}}\alpha _{}=\frac{\psi W_+^{}\psi ^{}W_+}{W_+W_{}^{}W_+^{}W_{}},$$
where the numerator is just the Wronskian of $`W_+`$ and $`W_{}`$. It is not difficult to prove that this takes the value $`\frac{2}{\xi }i`$, so that $`\alpha _\pm `$ simplifies to the following form
$$\alpha _+=\frac{\xi }{2}i\left(\psi W_{}^{}\psi ^{}W_{}\right)$$
(39)
$$\alpha _{}=\frac{\xi }{2}i\left(\psi W_+^{}\psi ^{}W_+\right).$$
(40)
Doing the derivative in $`x`$ in the eqs. $`(39)`$ and $`(40)`$, we have
$$\frac{d\alpha _\pm }{dx}=\frac{\xi }{2}i\left(\psi ^{}W_{}^{}+\psi W_{}^{\prime \prime }\psi ^{\prime \prime }W_{}\psi ^{}W_{}^{}\right).$$
(41)
In the brackets, the first and the fourth terms are zero; recalling that
$$\psi ^{\prime \prime }+\frac{1}{\xi ^2}r(x)\psi =0\&W_\pm ^{\prime \prime }+\left[\frac{1}{\xi ^2}r(x)+s(x)\right]W_\pm =0,$$
we can write eq. $`(41)`$ in the form
$$\frac{d\alpha _\pm }{dx}=\frac{\xi }{2}i\left[\psi \left(\frac{r}{\xi ^2}+s\right)W_{}+\frac{r}{\xi ^2}\psi W_{}\right]$$
$$\frac{d\alpha _\pm }{dx}=\frac{\xi }{2}is(x)\psi (x)W_{}(x),$$
(42)
which based on eqs. $`(33)`$ and $`(37)`$ becomes
$$\frac{d\alpha _\pm }{dx}=\frac{\xi }{2}i\frac{s(x)}{\left[r(x)\right]^{\frac{1}{2}}}\left[\alpha _\pm +\alpha _{}\mathrm{exp}\left(\frac{2}{\xi }i_{x_k}^x\sqrt{r(x)}𝑑x\right)\right].$$
(43)
Eqs. $`(42)`$ and $`(43)`$ are useful for estimating the WKB error in the 1D case.
The reason for which $`\frac{d\alpha _\pm }{dx}`$ can be considered as a measure of the WKB errors is that in the eqs. $`(31)`$ and $`(32)`$ the constants $`a`$, $`b`$ and $`c`$, $`d`$, respectively, give only approximate solutions $`\psi `$, while the functions $`\alpha _\pm `$ when introduced in the eqs. $`(37)`$ and $`(38)`$ produce exact $`\psi `$ solutions. From the geometrical viewpoint the derivative gives the slope of the tangent to these functions and indicates the measure in which $`\alpha _\pm `$ deviates from the constants $`a`$, $`b`$, $`c`$ and $`d`$.
4N. Note: The original (J)WKB papers are the following:
G. Wentzel, “Eine Verallgemeinerung der Wellenmechanik”, \[“A generalization of wave mechanics”\],
Zeitschrift für Physik 38, 518-529 (1926) \[received on 18 June 1926\]
L. Brillouin, “La mécanique ondulatoire de Schrödinger: une méthode générale de resolution par approximations successives”, \[“Schrödinger’s wave mechanics: a general method of solving by succesive approximations”\],
Comptes Rendus Acad. Sci. Paris 183, 24-26 (1926) \[received on 5 July 1926\]
H.A. Kramers, “Wellenmechanik und halbzahlige Quantisierung”, \[“Wave mechanics and half-integer quantization”\],
Zf. Physik 39, 828-840 (1926) \[received on 9 Sept. 1926\]
H. Jeffreys, “On certain approx. solutions of linear diff. eqs. of the second order”,
Proc. Lond. Math. Soc. 23, 428-436 (1925)
4P. Problems
Problem 4.1
Employ the WKB method for a particle of energy $`E`$ moving in a potential $`u(x)`$ of the form shown in fig. 4.1.
Solution
The Schrödinger equation is
$$\frac{d^2\psi }{dx^2}+\frac{2m}{\mathrm{}^2}\left[Eu(x)\right]\psi =0.$$
(44)
As one can see, we have
$$r(x)=\frac{2m}{\mathrm{}^2}\left[Eu(x)\right]\{\begin{array}{cc}\text{is positive for }a<x<b\hfill & \\ \text{is negative for }x<a,x>b\text{.}\hfill & \end{array}$$
If $`\psi (x)`$ corresponds to the region $`x<a`$, when passing to the interval $`a<x<b`$, the connection formula is given by eq. $`(29)`$ telling us that
$$\psi (x)\frac{A}{\left[Eu\right]^{\frac{1}{4}}}\mathrm{cos}\left(_a^x\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}𝑑x\frac{\pi }{4}\right),$$
(45)
where $`A`$ is an arbitrary constant.
When $`\psi (x)`$ corresponds to the region $`x>b`$, when passing to the segment $`a<x<b`$, we have in a similar way
$$\psi (x)\frac{B}{\left[Eu\right]^{\frac{1}{4}}}\mathrm{cos}\left(_x^b\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}𝑑x\frac{\pi }{4}\right),$$
(46)
where $`B`$ is another arbitrary constant. The reason why the connection formula is again given by eq. (29) is easily understood examinining what happens when the particle reaches the second classical turning point at $`x=b`$. This produces the inversion of the direction of motion. Thus, the particle appears to come from the right toward the left. In other words, we are in the first case (from the left to the right), only that as seen in a mirror placed at the point $`x=a`$.
These two expressions should be the same, independently of the constants $`A`$ and $`B`$, so that
$$\mathrm{cos}\left(_a^x\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}𝑑x\frac{\pi }{4}\right)=\mathrm{cos}\left(_x^b\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}𝑑x\frac{\pi }{4}\right)$$
$$\mathrm{cos}\left(_a^x\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}𝑑x\frac{\pi }{4}\right)+\mathrm{cos}\left(_x^b\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}𝑑x\frac{\pi }{4}\right)=0.$$
(47)
Recalling that
$$\mathrm{cos}A+\mathrm{cos}B=2\mathrm{cos}\left(\frac{A+B}{2}\right)\mathrm{cos}\left(\frac{AB}{2}\right),$$
eq. $`(47)`$ can be written
$$2\mathrm{cos}\left[\frac{1}{2}(_a^x\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}dx\frac{\pi }{4}+_x^b\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}dx\frac{\pi }{4})\right]$$
$$\mathrm{cos}\left[\frac{1}{2}(_a^x\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}dx\frac{\pi }{4}_x^b\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}dx+\frac{\pi }{4})\right]=0,$$
(48)
which implies that the arguments of the cosinusoids are integer multiples of $`\frac{\pi }{2}`$. On the other hand, the argument of the second cosinusoid do not lead to a nontrivial result. Therefore, we pay attention only to the argument of the first cosinusoid, which prove to be essential for getting an important result
$$\frac{1}{2}\left(_a^x\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}𝑑x\frac{\pi }{4}+_x^b\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}𝑑x\frac{\pi }{4}\right)=\frac{n}{2}\pi \text{for n odd}$$
$$_a^b\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}𝑑x\frac{\pi }{2}=n\pi $$
$$_a^b\sqrt{\frac{2m}{\mathrm{}^2}(Eu)}𝑑x=(n+\frac{1}{2})\pi $$
$$_a^b\sqrt{2m(Eu)}𝑑x=(n+\frac{1}{2})\pi \mathrm{}.$$
(49)
This result is very similar to the Bohr - Sommerfeld quantization rules.
We recall that Bohr’s postulate says that the orbital angular momentum of an electron moving on an “allowed atomic orbit” is quantized as $`L=n\mathrm{}`$, $`n=1,2,3,\mathrm{}`$. We also recall that the Wilson - Sommerfeld quantization rules assert that any coordinate of a system that varies periodically in time should satisfy the ‘quantum’ condition: $`p_q𝑑q=n_qh`$, where $`q`$ is a periodic coordinate, $`p_q`$ is the associated momentum, $`n_q`$ is an integer, and $`h`$ is Planck’s constant. One can see that the WKB result is indeed very similar.
Problem 4.2
Estimate the error of the WKB solution WKB at a point $`x_1x_k`$, where $`x_k`$ is a classical turning point for the differential equation $`y^{\prime \prime }+xy=0`$. The solution of this problem is of importance in the study of uniform fields, such as the gravitational and electric fields generated by large planes.
Solution:
For this differential equation we have
$$\xi =1,r(x)=x\&s(x)=\frac{5}{16}x^2.$$
$`r(x)=x`$ has a single zero at $`x_k=0`$, therefore for $`x0`$:
$$W_\pm =x^{\frac{1}{4}}\mathrm{exp}\left(\pm i_0^x\sqrt{x}𝑑x\right)=x^{\frac{1}{4}}\mathrm{exp}\left(\pm \frac{2}{3}ix^{\frac{3}{2}}\right).$$
(50)
Derivating $`W_\pm `$ up to the second derivative in $`x`$, we realize that the following differential equation is satisfied
$$W_\pm ^{\prime \prime }+(x\frac{5}{16}x^2)W_\pm =0.$$
(51)
The exact solution $`y(x)`$ of the latter differential equation can be written as a linear combination of $`W_\pm `$, as it has been shown in the corresponding section where the WKB error has been tackled; recall that the following form of the linear combination was proposed therein
$$y(x)=\alpha _+(x)W_+(x)+\alpha _{}(x)W_{}(x)$$
For large $`x`$, the general solution of the differential equation can be written in the WKB approximation as follows
$$y(x)=Ax^{\frac{1}{4}}\mathrm{cos}\left(\frac{2}{3}x^{\frac{3}{2}}+\delta \right)\text{for}x\mathrm{}.$$
(52)
Thus, $`\alpha _+\frac{A}{2}e^{i\delta }`$ and $`\alpha _{}\frac{A}{2}e^{i\delta }`$ for $`x\mathrm{}`$. We want to calculate the error due to these WKB solutions. A simple measure of this error is the deviation of $`\alpha _+`$ and of $`\alpha _{}`$ relative to the constants $`A`$. Using the equation
$$\frac{d\alpha _\pm }{dx}=\frac{\xi }{2}i\frac{s(x)}{\sqrt{r(x)}}\left[\alpha _\pm +\alpha _{}\mathrm{exp}\left(2i_{x_k}^x\sqrt{r(x)}𝑑x\right)\right]$$
and doiing the corresponding substitutions, one gets
$$\frac{d\alpha _\pm }{dx}=\frac{i}{2}\left(\frac{5}{16}x^2\right)x^{\frac{1}{2}}\left[\frac{A}{2}e^{\pm i\delta }+\frac{A}{2}e^{i\delta }\mathrm{exp}\left(2i\frac{2}{3}x^{\frac{3}{2}}\right)\right].$$
(53)
Taking $`\mathrm{\Delta }\alpha _\pm `$ as the changes displayed by $`\alpha _\pm `$ when $`x`$ varies between $`x_1`$ and $`\mathrm{}`$, we can do the required calculation by means of
$$\frac{\mathrm{\Delta }\alpha _\pm }{A/2}=\frac{2}{A}_{x_1}^{\mathrm{}}\frac{d\alpha _\pm }{dx}𝑑x=$$
$$=\pm i\frac{5}{32}e^{\pm i\delta }\left[\frac{2}{3}x_1^{\frac{3}{2}}+e^{2i\delta }_{x_1}^{\mathrm{}}x^{\frac{5}{2}}\mathrm{exp}\left(i\frac{4}{3}x^{\frac{3}{2}}\right)𝑑x\right].$$
(54)
The second term in the parentheses is less important than the first one because the complex exponential oscillates between $`1`$ and $`1`$ and therefore $`x^{\frac{5}{2}}<x^{\frac{3}{2}}`$. Consequently
$$\frac{\mathrm{\Delta }\alpha _\pm }{A/2}\pm \frac{5}{48}ie^{\pm i\delta }x_1^{\frac{3}{2}},$$
(55)
and as we can see the error introduced by the WKB approximation is indeed small if we take into account that the complex exponential oscillates between $`1`$ and $`1`$, while $`x_1^{\frac{3}{2}}`$ is also small.
## 5. THE HARMONIC OSCILLATOR (HO)
## The solution of the Schrödinger eq. for HO
The HO can be considered as a paradigm of Physics. Its utility is manifest in many areas from classical physics until quantum electrodynamics and theories of gravitational collapse.
From classical mechanics we know that many complicated potentials can be well approximated near their equilibrium positions by HO potentials
$$V(x)\frac{1}{2}V^{\prime \prime }(a)(xa)^2.$$
(1)
This is a 1D case. For this case, the classical Hamiltonian function of a particle of mass m, oscillating at the frequency $`\omega `$ has the following form:
$$H=\frac{p^2}{2m}+\frac{1}{2}m\omega ^2x^2$$
(2)
and the quantum Hamiltonian corresponding to the space of configurations is given by
$$\widehat{H}=\frac{1}{2m}(i\mathrm{}\frac{d}{dx})^2+\frac{1}{2}m\omega ^2x^2$$
(3)
$$\widehat{H}=\frac{\mathrm{}^2}{2m}\frac{d^2}{dx^2}+\frac{1}{2}m\omega ^2x^2.$$
(4)
Since we consider a time-independent potential, the eigenfunctions $`\mathrm{\Psi }_n`$ and the eigenvalues $`E_n`$ are obtained by means of the time-independent Schrödinger equation
$$\widehat{H}\mathrm{\Psi }_n=E_n\mathrm{\Psi }_n.$$
(5)
For the HO Hamiltonian, the Schrödinger eq. is
$$\frac{d^2\mathrm{\Psi }}{dx^2}+\left[\frac{2mE}{\mathrm{}^2}\frac{m^2\omega ^2}{\mathrm{}^2}x^2\right]\mathrm{\Psi }=0.$$
(6)
We cancealed the subindices of $`E`$ and $`\mathrm{\Psi }`$ because they are not of any importance here. Defining
$$k^2=\frac{2mE}{\mathrm{}^2}$$
(7)
$$\lambda =\frac{m\omega }{\mathrm{}},$$
(8)
the Schrödinger eq. becomes
$$\frac{d^2\mathrm{\Psi }}{dx^2}+[k^2\lambda ^2x^2]\mathrm{\Psi }=0,$$
(9)
which is known as Weber’s differential equation in mathematics.
We shall make now the transformation
$$y=\lambda x^2.$$
(10)
In general, by changing the variable from $`x`$ to $`y`$, the differential operators take the form
$$\frac{d}{dx}=\frac{dy}{dx}\frac{d}{dy}$$
(11)
$$\frac{d^2}{dx^2}=\frac{d}{dx}(\frac{dy}{dx}\frac{d}{dy})=\frac{d^2y}{dx^2}\frac{d}{dy}+(\frac{dy}{dx})^2\frac{d^2}{dy^2}.$$
(12)
Applying this obvious rule to the proposed transformation we obtain the following differential eq. in the $`y`$ variable
$$y\frac{d^2\mathrm{\Psi }}{dy^2}+\frac{1}{2}\frac{d\mathrm{\Psi }}{dy}+[\frac{k^2}{4\lambda }\frac{1}{4}y]\mathrm{\Psi }=0,$$
(13)
and, by definind :
$$\kappa =\frac{k^2}{2\lambda }=\frac{\overline{k}^2}{2m\omega }=\frac{E}{\mathrm{}\omega },$$
(14)
we get
$$y\frac{d^2\mathrm{\Psi }}{dy^2}+\frac{1}{2}\frac{d\mathrm{\Psi }}{dy}+[\frac{\kappa }{2}\frac{1}{4}y]\mathrm{\Psi }=0.$$
(15)
Let us try to solve this equation by first doing its asymptotic analysis in the limit $`y\mathrm{}`$. We first rewrite the previous equation in the form
$$\frac{d^2\mathrm{\Psi }}{dy^2}+\frac{1}{2y}\frac{d\mathrm{\Psi }}{dy}+[\frac{\kappa }{2y}\frac{1}{4}]\mathrm{\Psi }=0.$$
(16)
We notice that in the limit $`y\mathrm{}`$ the equation behaves as follows
$$\frac{d^2\mathrm{\Psi }_{\mathrm{}}}{dy^2}\frac{1}{4}\mathrm{\Psi }_{\mathrm{}}=0.$$
(17)
This equation has as solution
$$\mathrm{\Psi }_{\mathrm{}}(y)=A\mathrm{exp}\frac{y}{2}+B\mathrm{exp}\frac{y}{2}.$$
(18)
Taking $`A=0`$, we eliminate $`\mathrm{exp}\frac{y}{2}`$ since it diverges in the limit $`y\mathrm{}`$, keeping only the attenuated exponential. We can now suggest that $`\mathrm{\Psi }`$ has the following form
$$\mathrm{\Psi }(y)=\mathrm{exp}\frac{y}{2}\psi (y).$$
(19)
Plugging it in the differential equation for $`y`$ ( eq. $`15`$) one gets:
$$y\frac{d^2\psi }{dy^2}+(\frac{1}{2}y)\frac{d\psi }{dy}+(\frac{\kappa }{2}\frac{1}{4})\psi =0.$$
(20)
The latter is a confluent hypergeometric equation <sup>4</sup><sup>4</sup>4It is also known as Kummer’s differential equation. :
$$z\frac{d^2y}{dz^2}+(cz)\frac{dy}{dz}ay=0.$$
(21)
The general solution of this equation is
$$y(z)=A_1F_1(a;c,z)+Bz_1^{1c}F_1(ac+1;2c,z),$$
(22)
where the confluent hypergeometric equation is defined by
$${}_{1}{}^{}F_{1}^{}(a;c,z)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(a)_nx^n}{(c)_nn!}.$$
(23)
Comparing now our equation with the standard confluent hypergeometric equation, one can see that the general solution of the first one is
$$\psi (y)=A_1F_1(a;\frac{1}{2},y)+By_1^{\frac{1}{2}}F_1(a+\frac{1}{2};\frac{3}{2},y),$$
(24)
where
$$a=(\frac{\kappa }{2}\frac{1}{4}).$$
(25)
If we keep these solutions in their present form, the normalization condition is not satisfied for the wavefunction because from the asymptotic behaviour of the confluent hypergeometric function <sup>5</sup><sup>5</sup>5 The asymptotic behavior for $`x\mathrm{}`$ is $`{}_{1}{}^{}F_{1}^{}(a;c,z)\frac{\mathrm{\Gamma }(c)}{\mathrm{\Gamma }(ca)}e^{ia\pi }x^a+\frac{\mathrm{\Gamma }(c)}{\mathrm{\Gamma }(a)}e^xx^{ac}.`$ it follows ( taking into account ony the dominant exponential behavior ) :
$$\mathrm{\Psi }(y)=e^{\frac{y}{2}}\psi (y)const.e^{\frac{y}{2}}y^{a\frac{1}{2}}.$$
(26)
The latter approximation leads to a divergence in the normalization integral, which physically is not acceptable. What one does in this case is to impose the termination condition for the series <sup>6</sup><sup>6</sup>6The truncation condition of the confluent hypergeometric series $`{}_{1}{}^{}F_{1}^{}(a;c,z)`$ is $`a=n`$, where $`n`$ is a nonnegative integer ( i.e., zero included). , that is , the series has only a finite number of terms and therefore it is a polynomial of $`n`$ order.
We thus notice that asking for a finite normalization constant (as already known, a necessary condition for the physical interpretation in terms of probabilities), leads us to the truncation of the series, which simultaneously generates the quantization of energy.
In the following we consider the two possible cases
$`1)a=n`$ and $`B=0`$
$$\frac{\kappa }{2}\frac{1}{4}=n.$$
(27)
The eigenfunctions are given by
$$\mathrm{\Psi }_n(x)=D_n\mathrm{exp}\frac{\lambda x^2}{2}_1F_1(n;\frac{1}{2},\lambda x^2)$$
(28)
and the energy is:
$$E_n=\mathrm{}\omega (2n+\frac{1}{2}).$$
(29)
$`2)a+\frac{1}{2}=n`$ and $`A=0`$
$$\frac{\kappa }{2}\frac{1}{4}=n+\frac{1}{2}.$$
(30)
The eigenfunctions are now
$$\mathrm{\Psi }_n(x)=D_n\mathrm{exp}\frac{\lambda x^2}{2}x_1F_1(n;\frac{3}{2},\lambda x^2),$$
(31)
whereas the stationary energies are
$$E_n=\mathrm{}\omega [(2n+1)+\frac{1}{2}].$$
(32)
The polynomials obtained by this truncation of the confluent hypergeometric series are called Hermite polynomials and in hypergeometric notation they are
$$H_{2n}(\eta )=(1)^n\frac{(2n)!}{n!}_1F_1(n;\frac{1}{2},\eta ^2)$$
(33)
$$H_{2n1}(\eta )=(1)^n\frac{2(2n+1)!}{n!}\eta _1F_1(n;\frac{3}{2},\eta ^2).$$
(34)
We can now combine the obtained results ( because some of them give us the even cases and the others the odd ones ) in a single expression for the eigenvalues and eigenfunctions
$$\mathrm{\Psi }_n(x)=D_n\mathrm{exp}\frac{\lambda x^2}{2}H_n(\sqrt{\lambda }x)$$
(35)
$$E_n=(n+\frac{1}{2})\mathrm{}\omega n=0,1,2\mathrm{}$$
(36)
The HO energy spectrum is equidistant, i.e., there is the same energy difference $`\mathrm{}\omega `$ îbetween any consequitive neighbour levels. Another remark refers to the minimum value of the energy of the oscillator; somewhat surprisingly it is not zero. This is considered by many people to be a pure quantum result because it is zero when $`\mathrm{}0`$. It is known as the zero point energy and the fact that it is different of zero is the main characteristic of all confining potentials.
The normalization constant is easy to calculate
$$D_n=\left[\sqrt{\frac{\lambda }{\pi }}\frac{1}{2^nn!}\right]^{\frac{1}{2}}.$$
(37)
Thus, one gets the following normalized eigenfunctions of the 1D operator
$$\mathrm{\Psi }_n(x)=\left[\sqrt{\frac{\lambda }{\pi }}\frac{1}{2^nn!}\right]^{\frac{1}{2}}\mathrm{exp}(\frac{\lambda x^2}{2})H_n(\sqrt{\lambda }x).$$
(38)
## Creation and anihilation operators: $`\widehat{a}^{}`$ and $`\widehat{a}`$
There is another approach to deal with the HO besides the conventional one of solving the Schrödinger equation. It is the algebraic method, also known as the method of creation and annihilation (ladder) operators. This is a very efficient procedure, which can be successfully applied to many quantum-mechanical problems, especially when dealing with discrete spectra.
Let us define two nonhermitic operators $`a`$ and $`a^{}`$ :
$$a=\sqrt{\frac{m\omega }{2\mathrm{}}}(x+\frac{ip}{m\omega })$$
(39)
$$a^{}=\sqrt{\frac{m\omega }{2\mathrm{}}}(x\frac{ip}{m\omega }).$$
(40)
These operators are known as annihilation operator and creation operator, respectively (the reason of this terminology will be seen in the following, though one can claim that it comes from quantum field theories).
Let us calculate the commutator of these operators
$$[a,a^{}]=\frac{m\omega }{2\mathrm{}}[x+\frac{ip}{m\omega },x\frac{ip}{m\omega }]=\frac{1}{2\mathrm{}}(i[x,p]+i[p,x])=1,$$
(41)
where we have used the commutator
$$[x,p]=i\mathrm{}.$$
(42)
Therefore the annihilation and creation operators do not commute, since we have
$$[a,a^{}]=1.$$
(43)
Let us also introduce the very important number operator $`\widehat{N}`$:
$$\widehat{N}=a^{}a.$$
(44)
This operator is hermitic as one can readily prove using $`(AB)^{}=B^{}A^{}`$ :
$$\widehat{N}^{}=(a^{}a)^{}=a^{}(a^{})^{}=a^{}a=\widehat{N}.$$
(45)
Considering now that
$$a^{}a=\frac{m\omega }{2\mathrm{}}(x^2+\frac{p^2}{m^2\omega ^2})+\frac{i}{2\mathrm{}}[x,p]=\frac{\widehat{H}}{\mathrm{}\omega }\frac{1}{2}$$
(46)
we notice that the Hamiltonian can be written in a quite simple form as a function of the number operator
$$\widehat{H}=\mathrm{}\omega (\widehat{N}+\frac{1}{2}).$$
(47)
The number operator bear this name because its eigenvalues are precisely the subindices of the eigenfunctions on which it acts
$$\widehat{N}n=nn,$$
(48)
where we have used the notation
$$Psi_n=n.$$
(49)
Applying this fact to $`(47)`$, we get
$$\widehat{H}n>=\mathrm{}\omega (n+\frac{1}{2})n>.$$
(50)
On the other hand, from the Schrödinger equation we know that $`\widehat{H}n>=En>`$. In this way, it comes out that the energy eigenvalues are given by
$$E_n=\mathrm{}\omega (n+\frac{1}{2}).$$
(51)
This result is identical (as it should be) to the result $`(36)`$.
We go ahead and show why the operators $`a`$ and $`a^{}`$ bear the names they have. For this, we calculate the commutators
$$[\widehat{N},a]=[a^{}a,a]=a^{}[a,a]+[a^{},a]a=a,$$
(52)
which can be obtained from $`[a,a]=0`$ and $`(43)`$. Similarly, let us calculate
$$[\widehat{N},a^{}]=[a^{}a,a^{}]=a^{}[a,a^{}]+[a^{},a^{}]a=a^{}.$$
(53)
Using these two commutators, we can write
$`\widehat{N}(a^{}n>)`$ $`=`$ $`([\widehat{N},a^{}]+a^{}\widehat{N})n>`$
$`=`$ $`(a^{}+a^{}\widehat{N})n>`$
$`=`$ $`a^{}(1+n)n>=(n+1)a^{}n>.`$
By a similar procedure, one can also obtain
$$\widehat{N}(an>)=([\widehat{N},a]+a\widehat{N})n>=(n1)an>.$$
(55)
The expression $`(54)`$ implies that one can consider the ket $`a^{}n>`$ as an eigenket of that number operator for which the eigenvalue is raised by one unit. In physical terms, this means that an energy quanta has been produced by the action of $`a^{}`$ on the ket. This already expains the name of creation operator. Similar comments with corresponding conclusion can be infered for the operator $`a`$, originating the name of annihilation operator (an energy quanta is eliminated from the system when this operator is put in action).
Moreover, eq. $`(54)`$ implies the proportionality of the kets $`a^{}n>`$ and $`n+1>`$:
$$a^{}n>=cn+1>,$$
(56)
where $`c`$ is a constant that should be determined. Considering in addition
$$(a^{}n>)^{}=<na=c^{}<n+1,$$
(57)
one can perform the following calculation
$$<na(a^{}n>)=c^{}<n+1(cn+1>)$$
(58)
$$<naa^{}n>=c^{}c<n+1n+1>$$
(59)
$$<naa^{}n>=c^2.$$
(60)
But from the commutation relation for the operators $`a`$ and $`a^{}`$
$$[a,a^{}]=aa^{}a^{}a=aa^{}\widehat{N}=1,$$
(61)
we have
$$aa^{}=\widehat{N}+1.$$
(62)
Substituting in $`(60)`$, we get
$$<n\widehat{N}+1n>=<nn>+<n\widehat{N}n>=n+1=c^2.$$
(63)
Asking conventionally for a positive and real $`c`$, the following value is obtained
$$c=\sqrt{n+1}.$$
(64)
Consequently, we have
$$a^{}n>=\sqrt{n+1}n+1>.$$
(65)
For the annihilation operator, following the same procedure one can get the following relation
$$an>=\sqrt{n}n1>.$$
(66)
Let us show now that the values of $`n`$ should be nonnegative integers. For this, we employ the positivity requirement for the norm, applying it to the state vector $`an>`$. The latter condition tells us that the interior product of the vector with its adjunct ($`(an>)^{}=<na^{}`$) should always be nonnegative
$$(<na^{})(an>)0.$$
(67)
This relationship is nothing else but
$$<na^{}an>=<n\widehat{N}n>=n0.$$
(68)
Thus, $`n`$ cannot be negativ. It should be an integer since were it not by applying iteratively the annihilation operator we would be lead to negative values of $`n`$, which would be a contradiction to the previous statement.
It is possible to express the state $`n`$ $`(n)`$ directly as a function of the ground state $`(0)`$ using the creation operator. Let us see how proceeds this important iteration
$`1=a^{}0`$ (69)
$`2=[{\displaystyle \frac{a^{}}{\sqrt{2}}}]1=[{\displaystyle \frac{(a^{})^2}{\sqrt{2!}}}]0`$ (70)
$`3=[{\displaystyle \frac{a^{}}{\sqrt{3}}}]2=[{\displaystyle \frac{(a^{})^3}{\sqrt{3!}}}]0`$ (71)
$`n=[{\displaystyle \frac{(a^{})^n}{\sqrt{n!}}}]0.`$ (72)
One can also apply this method to get the eigenfunctions in the configuration space. To achieve this, we start with the ground state
$$a0=0.$$
(73)
In the $`x`$ representation, we have
$$\widehat{a}\mathrm{\Psi }_0(x)=\sqrt{\frac{m\omega }{2\mathrm{}}}(x+\frac{ip}{m\omega })\mathrm{\Psi }_0(x)=0.$$
(74)
Recalling the form of the momentum operator in the $`x`$ representation, we can obtain a differential equation for the wavefunction of the ground state. Moreover, introducing the definition $`x_0=\sqrt{\frac{\mathrm{}}{m\omega }}`$, we have
$$(x+x_0^2\frac{d}{dx})\mathrm{\Psi }_0=0.$$
(75)
The latter equation can be readily solved, and normalizing (its integral from $`\mathrm{}`$ to $`\mathrm{}`$ should be equal to unity), we obtain the wavefunction of the ground state
$$\mathrm{\Psi }_0(x)=(\frac{1}{\sqrt{\sqrt{\pi }x_0}})e^{\frac{1}{2}(\frac{x}{x_0})^2}.$$
(76)
The rest of the eigenfunctions, which describe the HO excited states, can be obtained employing the creation operator. The procedure is the following
$`\mathrm{\Psi }_1=a^{}\mathrm{\Psi }_0=({\displaystyle \frac{1}{\sqrt{2}x_0}})(xx_0^2{\displaystyle \frac{d}{dx}})\mathrm{\Psi }_0`$ (77)
$`\mathrm{\Psi }_2={\displaystyle \frac{1}{\sqrt{2}}}(a^{})^2\mathrm{\Psi }_0={\displaystyle \frac{1}{\sqrt{2!}}}({\displaystyle \frac{1}{\sqrt{2}x_0}})^2(xx_0^2{\displaystyle \frac{d}{dx}})^2\mathrm{\Psi }_0.`$ (78)
By mathematical induction, one can show that
$$\mathrm{\Psi }_n=\frac{1}{\sqrt{\sqrt{\pi }2^nn!}}\frac{1}{x_0^{n+\frac{1}{2}}}(xx_0^2\frac{d}{dx})^ne^{\frac{1}{2}(\frac{x}{x_0})^2}.$$
(79)
## Time evolution of the oscillator
In this section we shall illustrate on the HO example the way of working with the Heisenberg representation in which the states are fixed in time and only the operators evolve. Thus, we shall consider the operators as functions of time and obtain explicitly the time evolution of the HO position and momentum operators, $`a`$ and $`a^{}`$, respectively. The Heisenberg equations of the motion for $`p`$ and $`x`$ are
$`{\displaystyle \frac{d\widehat{p}}{dt}}`$ $`=`$ $`{\displaystyle \frac{}{\widehat{x}}}V(\widehat{𝐱})`$ (80)
$`{\displaystyle \frac{d\widehat{x}}{dt}}`$ $`=`$ $`{\displaystyle \frac{\widehat{p}}{m}}.`$ (81)
Hence the equations of the motion for $`x`$ and $`p`$ in the HO case are the following
$`{\displaystyle \frac{d\widehat{p}}{dt}}`$ $`=`$ $`m\omega ^2\widehat{x}`$ (82)
$`{\displaystyle \frac{d\widehat{x}}{dt}}`$ $`=`$ $`{\displaystyle \frac{\widehat{p}}{m}}.`$ (83)
These are a pair of coupled equations, which are equivalent to a pair of uncoupled equations for the creation and annihilation operators. Explicitly, we have
$`{\displaystyle \frac{da}{dt}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{m\omega }{2\mathrm{}}}}{\displaystyle \frac{d}{dt}}(\widehat{x}+{\displaystyle \frac{i\widehat{p}}{m\omega }})`$ (84)
$`{\displaystyle \frac{da}{dt}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{m\omega }{2\mathrm{}}}}({\displaystyle \frac{d\widehat{x}}{dt}}+{\displaystyle \frac{i}{m\omega }}{\displaystyle \frac{d\widehat{p}}{dt}}).`$ (85)
Substituting $`(82)`$ and $`(83)`$ in $`(85)`$, we get
$$\frac{da}{dt}=\sqrt{\frac{m\omega }{2\mathrm{}}}(\frac{\widehat{p}}{m}i\omega \widehat{x})=i\omega a.$$
(86)
Similarly, one can obtain a differential equation for the creation operator
$$\frac{da^{}}{dt}=i\omega a^{}.$$
(87)
The differential evolution equations for the creation and annihilation operators can be immediately integrated leading to the explicit evolution of these operators as follows
$`a(t)`$ $`=`$ $`a(0)e^{i\omega t}`$ (88)
$`a^{}(t)`$ $`=`$ $`a^{}(0)e^{i\omega t}.`$ (89)
It is worth noting based on these results and eqs. $`(44)`$ and $`(47)`$ that both the Hamiltonian and the number operator are not time dependent.
Using the latter two results, we can obtain the position and momentum operators as functions of time as far as they are expressed in terms of the creation and annihilation operators
$`\widehat{x}`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mathrm{}}{2m\omega }}}(a+a^{})`$ (90)
$`\widehat{p}`$ $`=`$ $`i\sqrt{{\displaystyle \frac{m\mathrm{}\omega }{2}}}(a^{}a).`$ (91)
Substituting them, one gets
$`\widehat{x}(t)`$ $`=`$ $`\widehat{x}(0)\mathrm{cos}\omega t+{\displaystyle \frac{\widehat{p}(0)}{m\omega }}\mathrm{sin}\omega t`$ (92)
$`\widehat{p}(t)`$ $`=`$ $`m\omega \widehat{x}(0)\mathrm{sin}\omega t+\widehat{p}(0)\mathrm{cos}\omega t.`$ (93)
The time evolution of these operators is the same as for the classical equations of the motion.
Thus, we have shown here the explicit evolution form of the four HO basic operators, and also we illustrated the effective way of working in the Heisenberg representation.
## The 3D HO
We commented on the importance in physics of the HO at the very beginning of our analysis of the quantum HO. If we will consider a 3D analog, we would be led to study a Taylor expansion in three variables<sup>7</sup><sup>7</sup>7It is possible to express the Taylor series in the neighbourhood of $`𝐫_\mathrm{𝟎}`$ as an exponential operator $`e^{[(xx_o)+(yy_o)+(zz_o)](\frac{}{x}+\frac{}{y}+\frac{}{z})}f(𝐫_𝐨).`$ retaining the terms up to the second order, we get a quadratic form in the most general case. The problem at hand in this approximation is not as simple as it might look from the examination of the corresponding potential
$$V(x,y,z)=ax^2+by^2+cz^2+dxy+exz+fyz.$$
(94)
There are however many systems with spherical symmetry or for which this symmetry is sufficiently exact. În acest caz:
$$V(x,y,z)=K(x^2+y^2+z^2),$$
(95)
which is equivalent to saying that the second unmixed partial derivatives have all the same value, denoted by $`K`$ in our case). We can add that this is a good approximation in the case in which the values of the mixed second partial derivatves are small in comparison to the unmixed ones.
When these conditions are satisfied and the potential is given by $`(95)`$, we say that the system is a 3D spherically symmetric HO.
The Hamiltonian in this case is of the form
$$\widehat{H}=\frac{\mathrm{}^2}{2m}^2+\frac{m\omega ^2}{2}r^2,$$
(96)
where the Laplace operator is given in spherical coordinates and $`r`$ is the spherical radial coordinate.
Since the potential is time independent the energy is conserved. In addition, because of the spherical symmetry the orbital momentum is also conserved. having two conserved quantities, we may say that to each of it one can associate a quantum number. Thus, we can assume that the eigenfunctions depend on two quantum numbers (even though for this case we shall see that another one will occur). Taking care of these comments, the equation of interest is
$$\widehat{H}\mathrm{\Psi }_{nl}=E_{nl}\mathrm{\Psi }_{nl}.$$
(97)
The Laplace operator in spherical coordinates reads
$$^2=\frac{^2}{r^2}+\frac{2}{r}\frac{}{r}\frac{\widehat{L}^2}{\mathrm{}^2r^2}$$
(98)
and can be also inferred from the known fact
$$\widehat{L}^2=\mathrm{}^2[\frac{1}{\mathrm{sin}\theta }\frac{}{\theta }(\mathrm{sin}\theta \frac{}{\theta })+\frac{1}{\mathrm{sin}\theta ^2}\frac{^2}{\phi ^2}].$$
(99)
The eigenfunctions of $`\widehat{L}^2`$ are the spherical harmonics, i.e.
$$\widehat{L}^2Y_{lm_l}(\theta ,\phi )=\mathrm{}^2l(l+1)Y_{lm_l}(\theta ,\phi )$$
(100)
The fact that the spherical harmonics ‘wear’ the quantum number $`m_l`$ introduces it in the total wavefunction $`\mathrm{\Psi }_{nlm_l}`$.
In order to achieve the separation of the variables and functions, the following substitution is proposed
$$\mathrm{\Psi }_{nlm_l}(r,\theta ,\phi )=\frac{R_{nl}(r)}{r}Y_{lm_l}(\theta ,\phi ).$$
(101)
Once this is plugged in the Schrödinger equation, the spatial part is separated from the angular one; the latter is identified with an operator that is proportional to the square of the orbital momentum, for which the eigenfunctions are the spherical harmonics, whereas for the spatial part the following equation is obtained
$$R_{nl}^{\prime \prime }+(\frac{2mE_{nl}}{\mathrm{}^2}\frac{m^2\omega ^2}{\mathrm{}^2}r^2\frac{l(l+1)}{r^2})R_{nl}(r)=0.$$
(102)
Using the definitions $`(7)`$ and $`(8)`$, the previous equation is precisely of the form $`(9)`$, unless the angular momentum term, which is commonly known as the unghiular, care în mod comun se cunoaşte ca angular momentum barrier
$$R_{nl}^{\prime \prime }+(k^2\lambda ^2r^2\frac{l(l+1)}{r^2})R_{nl}=0.$$
(103)
To solve this equation, we shall start with its asymptotic analysis. If we shall consider first $`r\mathrm{}`$, we notice that the orbital momentum term is negligible, so that in this limit the asymptotic behavior is similar to that of $`(9)`$, leading to
$$R_{nl}(r)\mathrm{exp}\frac{\lambda r^2}{2}\text{for}limr\mathrm{}.$$
(104)
If now we pass to the behavior close to zero, we can see that the dominant term is that of the orbital momentum, i.e., the differential equation $`(102)`$ in this limit turns into
$$R_{nl}^{\prime \prime }\frac{l(l+1)}{r^2}R_{nl}=0.$$
(105)
This is a differential equation of the Euler type <sup>8</sup><sup>8</sup>8An equation of the Euler type has the form
$$x^ny^{(n)}(x)+x^{n1}y^{(n1)}(x)+\mathrm{}+xy^{}(x)+y(x)=0.$$
Its solutions are of the type $`x^\alpha `$ that are plugged in the equation obtaining a polynomial in $`\alpha `$. , whose two independent solutions are
$$R_{nl}(r)r^{l+1}\text{or}r^l\text{for}limr0.$$
(106)
The previous arguments lead to proposing the substitution
$$R_{nl}(r)=r^{l+1}\mathrm{exp}\frac{\lambda r^2}{2}\varphi (r).$$
(107)
One can also use another substitution
$$R_{nl}(r)=r^l\mathrm{exp}\frac{\lambda r^2}{2}v(r),$$
(108)
which, however, produces the same solutions as $`(107)`$ (showing this is a helpful exercise). Substituing $`(107)`$ in $`(103)`$, the following differential equation for $`\varphi `$ is obtained
$$\varphi ^{\prime \prime }+2(\frac{l+1}{r}\lambda r)\varphi ^{}[\lambda (2l+3)k^2]\varphi =0.$$
(109)
Using now the change of variable $`w=\lambda r^2`$, one gets
$$w\varphi ^{\prime \prime }+(l+\frac{3}{2}w)\varphi ^{}[\frac{1}{2}(l+\frac{3}{2})\frac{\kappa }{2}]\varphi =0,$$
(110)
where $`\kappa =\frac{k^2}{2\lambda }=\frac{E}{\mathrm{}\omega }`$ has been introduced. We see that we found again a differential equation of the confluent hypergeometric type having the solutions (see $`(21)`$ and $`(22)`$)
$$\varphi (r)=A_1F_1[\frac{1}{2}(l+\frac{3}{2}\kappa );l+\frac{3}{2},\lambda r^2]+Br_1^{(2l+1)}F_1[\frac{1}{2}(l+\frac{1}{2}\kappa );l+\frac{1}{2},\lambda r^2].$$
(111)
The second particular solution cannot be normalized because diverges strongly in zero. This forces one to take $`B=0`$, therefore
$$\varphi (r)=A_1F_1[\frac{1}{2}(l+\frac{3}{2}\kappa );l+\frac{3}{2},\lambda r^2].$$
(112)
Using the same arguments as in the 1D HO case, that is, imposing a regular solution at infinity, leads to the truncation of the series, which implies the quantization of the energy. The truncation is explicitly
$$\frac{1}{2}(l+\frac{3}{2}\kappa )=n,$$
(113)
where introducing $`\kappa `$ we get the energy spectrum
$$E_{nl}=\mathrm{}\omega (2n+l+\frac{3}{2}).$$
(114)
One can notice that for the 3D spherically symmetric HO there is a zero point energy $`\frac{3}{2}\mathrm{}\omega `$.
The unnormalized eigenfunctions are
$$\mathrm{\Psi }_{nlm}(r,\theta ,\phi )=r^le_1^{\frac{\lambda r^2}{2}}F_1(n;l+\frac{3}{2},\lambda r^2)Y_{lm}(\theta ,\phi ).$$
(115)
## 5P. Problems
### Problem 5.1
Determine the eigenvalues and eigenfunctions of the HO in the momentum space.
The quantum HO Hamiltonian reads
$$\widehat{H}=\frac{\widehat{p}^2}{2m}+\frac{1}{2}m\omega ^2\widehat{x}^2.$$
In the momentum space, the operators $`\widehat{x}`$ and $`\widehat{p}`$ have the following form
$$\widehat{p}p$$
$$\widehat{x}i\mathrm{}\frac{}{p}.$$
Thus, the HO quantum Hamiltonian in the momentum representation is
$$\widehat{H}=\frac{p^2}{2m}\frac{1}{2}m\omega ^2\mathrm{}^2\frac{d^2}{dp^2}.$$
We have to solve the eigenvalue problem (i.e., to get the eigenfunctions and the eigenvalues) given by $`(5)`$, which, with the previous Hamiltonian, turns into the following differential equation
$$\frac{d^2\mathrm{\Psi }(p)}{dp^2}+(\frac{2E}{m\mathrm{}^2\omega ^2}\frac{p^2}{m^2\mathrm{}^2\omega ^2})\mathrm{\Psi }(p)=0.$$
(116)
One can see that this equation is identical, up to some constants, with the differential equation in the space of configurations (eq. $`(6)`$ ). Just to show another way of solving it, we define two parameters, which are analogous to those in $`(7)`$ and $`(8)`$
$$k^2=\frac{2E}{m\mathrm{}^2\omega ^2}\lambda =\frac{1}{m\mathrm{}\omega }.$$
(117)
With these definitions, we get the differential eq. $`(9)`$ and therefore the solution sought for (after performing the asymptotic analysis) is of the form
$$\mathrm{\Psi }(y)=e^{\frac{1}{2}y}\varphi (y),$$
(118)
where $`y=\lambda p^2`$ and $`\lambda `$ is defined in $`(117)`$. Substitute $`(118)`$ in $`(116)`$ taking care to put $`(118)`$ in the variable $`p`$. One gets a differential equation in the variable $`\varphi `$
$$\frac{d^2\varphi (p)}{dp^2}2\lambda p\frac{d\varphi (p)}{dp}+(k^2\lambda )\varphi (p)=0.$$
(119)
We shall now make the change of variable $`u=\sqrt{\lambda }p`$ that finally leads us to the Hermite equation
$$\frac{d^2\varphi (u)}{du^2}2u\frac{d\varphi (u)}{du}+2n\varphi (u)=0,$$
(120)
where $`n`$ is a nonnegative integer and where we have put
$$\frac{k^2}{\lambda }1=2n.$$
From here and the definitions given in $`(117)`$ one can easily conclude that the eigenvalues are given by
$$E_n=\mathrm{}\omega (n+\frac{1}{2}).$$
The solutions for $`(120)`$ are the Hermite polynomials $`\varphi (u)=H_n(u)`$ and the unnormalized eigenfunctions are
$$\mathrm{\Psi }(p)=Ae^{\frac{\lambda }{2}p^2}H_n(\sqrt{\lambda }p).$$
### Problem 5.2
Prove that the Hermite polynomials can be expressed in the following integral representation
$$H_n(x)=\frac{2^n}{\sqrt{\pi }}_{\mathrm{}}^{\mathrm{}}(x+iy)^ne^{y^2}𝑑y.$$
(121)
This representation of Hermite polynomials is not really usual, though it can prove useful in many cases. In order to accomplish the proof, we shall expand expand the integral and next prove that what we’ve got is identical to the series expansion of the Hermite polynomials that reads
$$\underset{k=0}{\overset{[\frac{n}{2}]}{}}\frac{(1)^kn!}{(n2k)!k!}(2x)^{n2k},$$
(122)
where the symbol $`[c]`$, indicating where the series terminates, denotes the greatest integer less or equal to $`c`$.
The first thing we shall do is to expand the binomial in the integral by using the well-known binomial theorem
$$(x+y)^n=\underset{m=0}{\overset{n}{}}\frac{n!}{(nm)!m!}x^{nm}y^m.$$
Thus
$$(x+iy)^n=\underset{m=0}{\overset{n}{}}\frac{n!}{(nm)!m!}i^mx^{nm}y^m,$$
(123)
which plugged in the integral leads to
$$\frac{2^n}{\sqrt{\pi }}\underset{m=0}{\overset{n}{}}\frac{n!}{(nm)!m!}i^mx^{nm}_{\mathrm{}}^{\mathrm{}}y^me^{y^2}𝑑y.$$
(124)
Inspecting of the integrand we realize that the integral is not zero when $`m`$ is even, whereas it is zero when $`m`$ is odd. Using the even notation $`m=2k`$, we get
$$\frac{2^n}{\sqrt{\pi }}\underset{k=0}{\overset{[\frac{n}{2}]}{}}\frac{n!}{(n2k)!(2k)!}i^{2k}x^{n2k}2_0^{\mathrm{}}y^{2k}e^{y^2}𝑑y.$$
(125)
Under the change of variable $`u=y^2`$, the integral turns into a gamma function
$$\frac{2^n}{\sqrt{\pi }}\underset{k=0}{\overset{[\frac{n}{2}]}{}}\frac{n!}{(n2k)!(2k)!}i^{2k}x^{n2k}_0^{\mathrm{}}u^{k\frac{1}{2}}e^u𝑑u,$$
(126)
more precisely $`\mathrm{\Gamma }(k+\frac{1}{2})`$, which can be expressed in terms of factorials ( of course for $`k`$ an integer)
$$\mathrm{\Gamma }(k+\frac{1}{2})=\frac{(2k)!}{2^{2k}k!}\sqrt{\pi }.$$
Plugging this expression in the sum and using $`i^{2k}=(1)^k`$, one gets
$$\underset{k=0}{\overset{[\frac{n}{2}]}{}}\frac{(1)^kn!}{(n2k)!k!}(2x)^{n2k},$$
(127)
which is identical to $`(122)`$, hence completing the proof.
### Problem 5.3
Show that Heisenberg’s uncertainty relation is satisfied by doing the calculation using the HO eigenfunctions .
We have to show that for any $`\mathrm{\Psi }_n`$, we have
$$<(\mathrm{\Delta }p)^2(\mathrm{\Delta }x)^2>\frac{\mathrm{}^2}{4},$$
(128)
where the notation $`<>`$ means the mean value.
We shall separately calculate $`<(\mathrm{\Delta }p)^2>`$ and $`<(\mathrm{\Delta }x)^2>`$, where each of these expressions is
$$<(\mathrm{\Delta }p)^2>=<(p<p>)^2>=<p^22p<p>+<p>^2>=<p^2><p>^2,$$
$$<(\mathrm{\Delta }x)^2>=<(x<x>)^2>=<x^22x<x>+<x>^2>=<x^2><x>^2.$$
First of all, we shall prove that both the mean of $`x`$ as well as of $`p`$ are zero. For the mean of $`x`$, we have
$$<x>=_{\mathrm{}}^{\mathrm{}}x[\mathrm{\Psi }_n(x)]^2𝑑x.$$
This integral is zero because the integrand is odd. Thus
$$<x>=0.$$
(129)
The same argument holds for the mean of $`p`$, if we do the calculation in the momentum space, employing the functions obtained in problem 1. It is sufficient to notice that the functional form is the same (only the symbol does change). Thus
$$<p>=0.$$
(130)
Let us now calculate the mean of $`x^2`$. We shall use the virial theorem <sup>9</sup><sup>9</sup>9We recall that the virial theorem in quantum mechanics asserts that
$$2<T>=<𝐫V(𝐫)>.$$
For a potential of the form $`V=\lambda x^n`$, the virial theorem gives
$$2<T>=n<V>,$$
where $`T`$ is the kinetic energy and $`V`$ is the potential energy.. We first notice that
$$<V>=\frac{1}{2}m\omega ^2<x^2>.$$
Therefore, it is possible to relate the mean of $`x^2`$ directly to the mean of the potential for this case (implying the usage of the virial theorem).
$$<x^2>=\frac{2}{m\omega ^2}<V>.$$
(131)
We also need the total energy
$$<H>=<T>+<V>,$$
for which again one can make use of the virial theorem (for $`n=2`$)
$$<H>=2<V>.$$
(132)
Thus, we obtain
$$<x^2>=\frac{<H>}{m\omega ^2}=\frac{\mathrm{}\omega (n+\frac{1}{2})}{m\omega ^2}$$
(133)
$$<x^2>=\frac{\mathrm{}}{m\omega }(n+\frac{1}{2}).$$
(134)
Similarly, the mean of $`p^2`$ can be readily calculated
$$<p^2>=2m<\frac{p^2}{2m}>=2m<T>=m<H>=m\mathrm{}\omega (n+\frac{1}{2}).$$
(135)
Employing $`(133)`$ and $`(135)`$, we have
$$<(\mathrm{\Delta }p)^2(\mathrm{\Delta }x)^2>=(n+\frac{1}{2})^2\mathrm{}^2.$$
(136)
Based on this result, we come to the conclusion that in the HO stationary states that actually have not been directly used, Heisenberg’s uncertainty relation is satisfied and it is at the minimum for the ground state, $`n=0`$.
### Problem 5.4
Obtain the matrix elements of the operators $`a`$, $`a^{}`$, $`\widehat{x}`$, and $`\widehat{p}`$.
Let us first find the matrix elements for the creation and annihilation operators, which are very helpful for all the other operators.
We shall use the relatinships $`(65)`$ and $`(66)`$, leading to
$$<man>=\sqrt{n}<mn1>=\sqrt{n}\delta _{m,n1}.$$
(137)
Similarly for the creation operator we have the result
$$<ma^{}n>=\sqrt{n+1}<mn+1>=\sqrt{n+1}\delta _{m,n+1}.$$
(138)
Let us proceed now with the calculation of the matrix elements of the position operator. For this, let us express this operator in terms of creation and annihilation operators. Using the definitions $`(39)`$ and $`(40)`$, one can immediately prove that the position operator is given by
$$\widehat{x}=\sqrt{\frac{\mathrm{}}{2m\omega }}(a+a^{}).$$
(139)
Employing this result, the matrix elements of the operator $`\widehat{x}`$ can be readily calculated
$`<m\widehat{x}n>`$ $`=`$ $`<m\sqrt{{\displaystyle \frac{\mathrm{}}{2m\omega }}}(a+a^{})n>`$ (140)
$`=`$ $`\sqrt{{\displaystyle \frac{\mathrm{}}{2m\omega }}}[\sqrt{n}\delta _{m,n1}+\sqrt{n+1}\delta _{m,n+1}].`$
Following the same procedure we can calculate the matrix elements of the momentum operator, just by taking into account that $`\widehat{p}`$ is given in terms of the creation and annihilation operators as follows
$$\widehat{p}=i\sqrt{\frac{m\mathrm{}\omega }{2}}(a^{}a).$$
(141)
This leads us to
$$<m\widehat{p}n>=i\sqrt{\frac{m\mathrm{}\omega }{2}}[\sqrt{n+1}\delta _{m,n+1}\sqrt{n}\delta _{m,n1}].$$
(142)
One can realize the ease of the calculations when the matrix elements of the creation and annihilation operators are used. Finally, we remark on the nondiagonality of the obtained matrix elements. This is not so much of a surprise because the employed representation is that of the number operator and none of the four operators do not commute with it.
### Problem 5.5
Find the mean values of $`\widehat{x}^2`$ and $`\widehat{p}^2`$ for the1D HO and use them to calculate the mean (expectation) values of the kinetic and potential energies. Compare the result with the virial theorem.
First of all, let us obtain the mean value of $`\widehat{x}^2`$. For this, we use eq. $`(139)`$ that leads us to
$$\widehat{x}^2=\frac{\mathrm{}}{2m\omega }(a^2+(a^{})^2+a^{}a+aa^{}).$$
(143)
Recall that the creation and annihilation operators do not commute. Based on $`(143)`$, we can calculate the mean value of $`\widehat{x}^2`$
$`<\widehat{x}^2>`$ $`=`$ $`<n\widehat{x}^2n>`$ (144)
$`=`$ $`{\displaystyle \frac{\mathrm{}}{2m\omega }}[\sqrt{n(n1)}\delta _{n,n2}+\sqrt{(n+1)(n+2)}\delta _{n,n+2}`$
$`+`$ $`n\delta _{n,n}+(n+1)\delta _{n,n}],`$
which shows that
$$<\widehat{x}^2>=<n\widehat{x}^2n>=\frac{\mathrm{}}{2m\omega }(2n+1).$$
(145)
In order to calculate the mean value of $`\widehat{p}^2`$ we use $`(141)`$ that helps us to express this operator in terms of the creation and annihilation operators
$$\widehat{p}^2=\frac{m\mathrm{}\omega }{2}(a^2+(a^{})^2aa^{}a^{}a).$$
(146)
This leads us to
$$<\widehat{p}^2>=<n\widehat{p}^2n>=\frac{m\mathrm{}\omega }{2}(2n+1).$$
(147)
The latter result practically gives us the mean kinetic energy
$$<\widehat{T}>=<\frac{\widehat{p}^2}{2m}>=\frac{1}{2m}<\widehat{p}^2>=\frac{\mathrm{}\omega }{4}(2n+1).$$
(148)
On the other hand, the mean value of the potential energy
$$<\widehat{V}>=<\frac{1}{2}m\omega ^2\widehat{x}^2>=\frac{1}{2}m\omega ^2<\widehat{x}^2>=\frac{\mathrm{}\omega }{4}(2n+1),$$
(149)
where $`(145)`$ has been used.
We can see that these mean values are equal for any $`n`$, which confirms the quantum virial theorem, telling us that for a quadratic (HO) potential, the mean values of the kinetic and potential energies should be equal and therefore be half of the mean value of the total energy.
6. THE HYDROGEN ATOM
## Introduction
In this chapter we shall study the hydrogen atom by solving the time-independent Schrödinger equation for the potential due to two charged particles, the electron and the proton, and the Laplacian operator in spherical coordinates. From the mathematical viewpoint, the method of separation of variables will be employed, and a physical interpretation of the wavefunction as solution of the Schrödinger equation in this important case will be provided, together with the interpretation of the quantum numbers and of the probability densities.
The very small spatial scale of the hydrogen atom is a clue that the related physical phenomena enter the domain of applicability of the quantum mechanics, for which the atomic processes have been a successful area since the early days of the quantum approaches. Quantum mechanics, as any other theoretical framework, gives relationships between observable quantities. Since the uncertainty principle leads to a substantial change in the understanding of observables at the conceptual level, it is important to have a clear idea on the notion of atomic observable. As a matter of fact, the real quantities on which quantum mechanics offers explicit answers and connections are always probabilites. Instead of saying, for example, that the radius of the electron orbit in the fundamental state of the hydrogen atom is always $`5.3\times 10^{11}`$ m, quantum mechanics asserts that this is a truly mean radius (not in the measurable sense). Thus, if one performs an appropriate experiment, one gets, precisely as in the case of the common arrangement of macroscopic detectors probing macroscopic properties of the matter, random values around the mean value $`5.3\times 10^{11}`$ m. In other words, from the viewpoint of the experimental errors there is no essential difference with regard to the classical physics. The fundamental difference is in the procedure of calculating the mean values within the theoretical framework.
As is known, for performing quantum-mechanical calculations, one needs a corresponding wave function $`\mathrm{\Psi }`$. Although $`\mathrm{\Psi }`$ has no direct physical interpretation, the square modulus $`\mathrm{\Psi }^2`$ calculated at an arbitrary position and given moment is proportional to the probability to find the particle in the infinitesimal neighbourhood of that point at the given time. The purpose of quantum mechanics is to determine $`\mathrm{\Psi }`$ for a specified particle in the prepared experimental conditions.
Before proceeding with the rigorous approaches of getting $`\mathrm{\Psi }`$ for the hydrogen electron, we will argue on several general requirements regarding the wave function. First, the integral of $`\mathrm{\Psi }^2`$ over all space should be finite if we really want to deal with a localizable electron. In addition, if
$$_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }^2𝑑V=0,$$
(1)
then the particle does not exist. $`\mathrm{\Psi }^2`$ cannot be negative or complex because of simple mathematical reasons. In general, it is convenient to identify $`\mathrm{\Psi }^2`$ with the probability P not just the proportionality. In order that $`\mathrm{\Psi }^2`$ be equal to P one imposes
$$_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }^2𝑑V=1,$$
(2)
because
$$_{\mathrm{}}^{\mathrm{}}\mathrm{P}𝑑V=1$$
(3)
is the mathematical way of saying that the particle exists at a point in space at any given moment. A wave function respecting eq. 2 is said to be normalized. Besides this, $`\mathrm{\Psi }`$ should be single valued, because P has a unique value at a given point and given time. Another condition is that $`\mathrm{\Psi }`$ and its partial first derivatives $`\frac{\mathrm{\Psi }}{x}`$, $`\frac{\mathrm{\Psi }}{y}`$, $`\frac{\mathrm{\Psi }}{z}`$ should be continuous at any arbitrary point.
The Schrödinger equation is considered as the fundamental equation of nonrelativistic quantum mechanics in the same sense in which Newton’s force law is the fundamental equation of motion of newtonian mechanics. Notice however that we have now a wave equation for a function $`\mathrm{\Psi }`$ which is not directly measurable.
Once the potential energy is given, one can solve the Schrödinger equation for $`\mathrm{\Psi }`$, implying the knowledge of the probability density $`\mathrm{\Psi }^2`$ as a function of $`x,y,z,t`$. In many cases of interest, the potential energy does not depend on time. Then, the Schrödinger equation simplifies considerably. Notice, for example, that for a 1D free particle the wave function can be written
$`\mathrm{\Psi }(x,t)`$ $`=`$ $`Ae^{(i/\mathrm{})(Etpx)}`$ (4)
$`=`$ $`Ae^{(iE/\mathrm{})t}e^{(ip/\mathrm{})x}`$
$`=`$ $`\psi (x)e^{(iE/\mathrm{})t},`$
i.e., $`\mathrm{\Psi }(x,t)`$ is the product of a time-dependent phase $`e^{(iE/\mathrm{})t}`$ and a stationary wave function $`\psi (x)`$.
In the general case, the stationary Schrödinger equation can be solved, under the aforementioned requirements, only for certain values of the energy E. This is not a mathematical difficulty, but merely a fundamental physical feature. To solve the Schrödinger equation for a given system means to get the wave function $`\psi `$, as a solution for which certain physical boundary condition hold and, in addition, as already mentioned, it is continuous together with its first derivative everywhere in space, is finite, and single valued. Thus, the quantization of energy occurs as a natural theoretical element in wave mechanics, whereas in practice as a universal phenomenon, characteristic for all stable microscopic systems.
## Schrödinger equation for the hydrogen atom
In this section, we shall apply the Schrödinger equation to the hydrogen atom, about which one knows that it is formed of a positive nucleus/proton of charge +$`e`$ and an electron of charge -$`e`$. The latter, being 1836 times smaller in mass than the proton, is by far more dynamic.
If the interaction between two particles is of the type $`u(r)=u(\stackrel{}{r}_1\stackrel{}{r}_2)`$, the problem of the motion is reduced both classically and quantum to the motion of a single particle in a field of spherical symmetry. Indeed, the Lagrangian
$$L=\frac{1}{2}m_1\dot{\stackrel{}{r}_1^2}+\frac{1}{2}m_2\dot{\stackrel{}{r}_2^2}u(\stackrel{}{r}_1\stackrel{}{r}_2)$$
(5)
is transformed, using
$$\stackrel{}{r}=\stackrel{}{r}_1\stackrel{}{r}_2$$
(6)
and
$$\stackrel{}{R}=\frac{m_1\stackrel{}{r}_1+m_2\stackrel{}{r}_2}{m_1+m_2},$$
(7)
in the Lagrangian
$$L=\frac{1}{2}M\dot{\stackrel{}{R}^2}+\frac{1}{2}\mu \dot{\stackrel{}{r}^2}u(r),$$
(8)
where
$$M=m_1+m_2$$
(9)
and
$$\mu =\frac{m_1m_2}{m_1+m_2}.$$
(10)
On the other hand, the momentum is introduced through the Lagrange formula
$$\stackrel{}{P}=\frac{L}{\dot{\stackrel{}{R}}}=M\dot{\stackrel{}{R}}$$
(11)
and
$$\stackrel{}{p}=\frac{L}{\dot{\stackrel{}{r}}}=m\dot{\stackrel{}{r}},$$
(12)
that allows to write the classical Hamilton function in the form
$$H=\frac{P^2}{2M}+\frac{p^2}{2m}+u(r).$$
(13)
Thus, one can obtain the hamiltonian operator for the corresponding quantum problem with commutators of the type
$$[P_i,P_k]=i\mathrm{}\delta _{ik}$$
(14)
and
$$[p_i,p_k]=i\mathrm{}\delta _{ik}.$$
(15)
These commutators implies a Hamiltonian operator of the form
$$\widehat{H}=\frac{\mathrm{}^2}{2M}_R^2\frac{\mathrm{}^2}{2m}_r^2+u(r),$$
(16)
which is fundamental for the study of the hydrogen atom by means of the stationary Schrödinger equation
$$\widehat{H}\psi =E\psi .$$
(17)
This form does not include relativistic effects, i.e., electron velocities close to the velocity of light in vacuum.
The potential energy $`u(r)`$ is the electrostatic one
$$u=\frac{e^2}{4\pi ϵ_0r}$$
(18)
There are two possibilities. The first is to express $`u`$ as a function of the cartesian coordinates $`x,y,z`$, substituing $`r`$ by $`\sqrt{x^2+y^2+z^2}`$. The second is to write the Schrödinger equation in spherical polar coordinates $`r,\theta ,\varphi `$. Because of the obvious spherical symmetry of this case, we shall deal with the latter approach, which leads to considerable mathematical simplifications.
In spherical coordinates, the Schrödinger equation reads
$$\frac{1}{r^2}\frac{}{r}\left(r^2\frac{\psi }{r}\right)+\frac{1}{r^2\mathrm{sin}\theta }\frac{}{\theta }\left(\mathrm{sin}\theta \frac{\psi }{\theta }\right)+\frac{1}{r^2\mathrm{sin}^2\theta }\frac{^2\psi }{\varphi ^2}+\frac{2m}{\mathrm{}^2}(Eu)\psi =0$$
(19)
Substituing (18), and multiplying the whole equation by $`r^2\mathrm{sin}^2\theta `$, one gets
$$\mathrm{sin}^2\theta \frac{}{r}\left(r^2\frac{\psi }{r}\right)+\mathrm{sin}\theta \frac{}{\theta }\left(\mathrm{sin}\theta \frac{\psi }{\theta }\right)+\frac{^2\psi }{\varphi ^2}+\frac{2mr^2\mathrm{sin}^2\theta }{\mathrm{}^2}\left(\frac{e^2}{4\pi ϵ_0r}+E\right)\psi =0.$$
(20)
This equation is a partial differential equation for the electron wavefunction $`\psi (r,\theta ,\varphi )`$ ‘within’ the atomic hydrogen. Together with the various conditions that the wavefunction $`\psi (r,\theta ,\varphi )`$ should fulfill \[for example, $`\psi (r,\theta ,\varphi )`$ should have a unique value at any spatial point ($`r,\theta ,\varphi `$)\], this equation specifies in a complete manner the behavior of the hydrogen electron. To see the explicit behavior, we shall solve eq. 20 for $`\psi (r,\theta ,\varphi )`$ and we shall interpret appropriately the obtained results.
## Separation of variables in spherical coordinates
The real usefulness of writing the hydrogen Schrödinger equation in spherical coordinates consists in the easy way of achieving the separation procedure in three independent equations, each of them being one-dimensional. The separation procedure is to seek the solutions for which the wavefunction $`\psi (r,\theta ,\varphi )`$ has the form of a product of three functions, each of one of the three spherical variables, namely $`R(r)`$, depending only on $`r`$; $`\mathrm{\Theta }(\theta )`$ depending only on $`\theta `$, and $`\mathrm{\Phi }(\varphi )`$ that depends only on $`\varphi `$. This is quite similar to the separation of the Laplace equation. Thus
$$\psi (r,\theta ,\varphi )=R(r)\mathrm{\Theta }(\theta )\mathrm{\Phi }(\varphi ).$$
(21)
The $`R(r)`$ function describes the differential variation of the electron wavefunction $`\psi `$ along the vector radius coming out from the nucleus, with $`\theta `$ and $`\varphi `$ assumed to be constant. The differential variation of $`\psi `$ with the polar angle $`\theta `$ along a meridian of an arbitrary sphere centered in the nucleus is described only by the function $`\mathrm{\Theta }(\theta )`$ for constant $`r`$ and $`\varphi `$. Finally, the function $`\mathrm{\Phi }(\varphi )`$ describes how $`\psi `$ varies with the azimuthal angle $`\varphi `$ along a parallel of an arbitrary sphere centered at the nucleus, under the conditions that $`r`$ and $`\theta `$ are kept constant.
Using $`\psi =R\mathrm{\Theta }\mathrm{\Phi }`$, one can see that
$$\frac{\psi }{r}=\mathrm{\Theta }\mathrm{\Phi }\frac{dR}{dr},$$
(22)
$$\frac{\psi }{\theta }=R\mathrm{\Phi }\frac{d\mathrm{\Theta }}{d\theta },$$
(23)
$$\frac{\psi }{\varphi }=R\mathrm{\Theta }\frac{d\mathrm{\Phi }}{d\varphi }.$$
(24)
Obviously, the same type of formulas are maintained for the unmixed higher-order derivatives. Subtituting them in eq. 20, and after deviding by $`R\mathrm{\Theta }\mathrm{\Phi }`$, we get
$$\frac{\mathrm{sin}^2\theta }{R}\frac{d}{dr}\left(r^2\frac{dR}{dr}\right)+\frac{\mathrm{sin}\theta }{\mathrm{\Theta }}\frac{d}{d\theta }\left(\mathrm{sin}\theta \frac{d\mathrm{\Theta }}{d\theta }\right)+\frac{1}{\mathrm{\Phi }}\frac{d^2\mathrm{\Phi }}{d\varphi ^2}+\frac{2mr^2\mathrm{sin}^2\theta }{\mathrm{}^2}\left(\frac{e^2}{4\pi ϵ_0r}+E\right)=0.$$
(25)
The third term of this equation is a function of the angle $`\varphi `$ only, while the other two terms are functions of $`r`$ and $`\theta `$. We rewrite now the previous equation in the form
$$\frac{\mathrm{sin}^2\theta }{R}\frac{}{r}\left(r^2\frac{R}{r}\right)+\frac{\mathrm{sin}\theta }{\mathrm{\Theta }}\frac{}{\theta }\left(\mathrm{sin}\theta \frac{\mathrm{\Theta }}{\theta }\right)+\frac{2mr^2\mathrm{sin}^2\theta }{\mathrm{}^2}\left(\frac{e^2}{4\pi ϵ_0r}+E\right)=\frac{1}{\mathrm{\Phi }}\frac{^2\mathrm{\Phi }}{\varphi ^2}.$$
(26)
This equation can be correct only if the two sides are equal to the same constant, because they are functions of different variables. It is convenient to denote this (separation) constant by $`m_l^2`$. The differential equation for the $`\mathrm{\Phi }`$ function is
$$\frac{1}{\mathrm{\Phi }}\frac{^2\mathrm{\Phi }}{\varphi ^2}=m_l^2.$$
(27)
If one substitutes $`m_l^2`$ in the right hand side of eq. 26 and devides the resulting equation by $`\mathrm{sin}^2\theta `$, after regrouping the terms, the fllowing result is obtained
$$\frac{1}{R}\frac{d}{dr}\left(r^2\frac{dR}{dr}\right)+\frac{2mr^2}{\mathrm{}^2}\left(\frac{e^2}{4\pi ϵ_0r}+E\right)=\frac{m_l^2}{\mathrm{sin}^2\theta }\frac{1}{\mathrm{\Theta }\mathrm{sin}\theta }\frac{d}{d\theta }\left(\mathrm{sin}\theta \frac{d\mathrm{\Theta }}{d\theta }\right).$$
(28)
Once again, we end up with an equation in which different variables occur in the two sides, thus forcing at equating of both sides to the same constant. For reasons that will become clear later on, we shall denote this constant by $`l(l+1)`$. The equations for the functions $`\mathrm{\Theta }(\theta )`$ and $`R(r)`$ reads
$$\frac{m_l^2}{\mathrm{sin}^2\theta }\frac{1}{\mathrm{\Theta }\mathrm{sin}\theta }\frac{d}{d\theta }\left(sin\theta \frac{d\mathrm{\Theta }}{d\theta }\right)=l(l+1)$$
(29)
and
$$\frac{1}{R}\frac{d}{dr}\left(r^2\frac{dR}{dr}\right)+\frac{2mr^2}{\mathrm{}^2}\left(\frac{e^2}{4\pi ϵ_0r}+E\right)=l(l+1).$$
(30)
The equations 27, 29 and 30 are usually written in the form
$$\frac{d^2\mathrm{\Phi }}{d\varphi ^2}+m_l^2\mathrm{\Phi }=0,$$
(31)
$$\frac{1}{\mathrm{sin}\theta }\frac{d}{d\theta }\left(\mathrm{sin}\theta \frac{d\mathrm{\Theta }}{d\theta }\right)+\left[l(l+1)\frac{m_l^2}{\mathrm{sin}^2\theta }\right]\mathrm{\Theta }=0,$$
(32)
$$\frac{1}{r^2}\frac{d}{dr}\left(r^2\frac{dR}{dr}\right)+\left[\frac{2m}{\mathrm{}^2}\left(\frac{e^2}{4\pi ϵ_0r}+E\right)\frac{l(l+1)}{r^2}\right]R=0.$$
(33)
Each of these equations is an ordinary differential equation for a function of a single variable. In this way, the Schrödinger equation for the hydrogen electron, which initially was a partial differential equation for a function $`\psi `$ of three variables, got a simple form of three 1D ordinary differential equations for unknown functions of one variable.
## Interpreting the separation constants: the quantum numbers
### The solution for the azimuthal part
Eq. 31 is readily solved leading to the following solution
$$\mathrm{\Phi }(\varphi )=A_\varphi e^{im_l\varphi },$$
(34)
where $`A_\varphi `$ is the integration constant. One of the conditions that any wavefunctions should fulfill is to have a unique value for any point in space. This applies to $`\mathrm{\Phi }`$ as a component of the full wavefunction $`\psi `$. One should notice that $`\varphi `$ and $`\varphi +2\pi `$ are identical in the same meridional plane. Therefore, one should have $`\mathrm{\Phi }(\varphi )=\mathrm{\Phi }(\varphi +2\pi )`$, i.e., $`A_\varphi e^{im_l\varphi }=A_\varphi e^{im_l(\varphi +2\pi )}`$. This can be fulfilled only if $`m_l`$ is zero or a positiv or negative integer $`(\pm 1,\pm 2,\pm 3,\mathrm{})`$. $`m_l`$ is known as the magnetic quantum number of the atomic electron and is related to the direction of the projection of the orbital momentum $`L_z`$. It comes into play whenever the effects of axial magnetic fields on the electron may show up. There is also a deep connection between $`m_l`$ and the orbital quantum number $`l`$, which in turn determines the modulus of the orbital momentum of the electron.
The interpretation of the orbital number $`l`$ does not miss some problems. Let us examine eq. 33 that corresponds to the radial wavefunction $`R(r)`$. This equation rules only the radial motion of the electron, i.e., with the relative distance with respect to the nucleus along some guiding ellipses. However, the total energy of the electron $`E`$ is also present. This energy includes the kineticelctron energy in its orbital motion that is not related to the radial motion. This contradiction can be eliminated by the following argument. The kinetic energy $`T`$ has two parts: $`T_{radial}`$ due to the radial oscillatory motion and $`T_{orbital}`$, which is due to the closed orbital motion. The potential energy $`V`$ of the electron is the electrostatic energy. Therefore, its total energy is
$$E=T_{radial}+T_{orbital}\frac{e^2}{4\pi ϵ_0r}.$$
(35)
Substituting this expression of $`E`$ in eq. 33 we get with some regrouping of the terms
$$\frac{1}{r^2}\frac{d}{dr}\left(r^2\frac{dR}{dr}\right)+\frac{2m}{\mathrm{}^2}\left[T_{radial}+T_{orbital}\frac{\mathrm{}^2l(l+1)}{2mr^2}\right]R=0.$$
(36)
If the last two terms in parentheses compansates between themselves, we get a differential equation for the pure radial motion. Thus, we impose the condition
$$T_{orbital}=\frac{\mathrm{}^2l(l+1)}{2mr^2}.$$
(37)
However, the orbital kinetic energy of the electron is
$$T_{orbital}=\frac{1}{2}mv_{orbital}^2$$
(38)
and since the orbital momentum of the electron $`L`$ is
$$L=mv_{orbital}r,$$
(39)
we can express the orbital kinetic energy in the form
$$T_{orbital}=\frac{L^2}{2mr^2}.$$
(40)
Therefore, we have
$$\frac{L^2}{2mr^2}=\frac{\mathrm{}^2l(l+1)}{2mr^2}$$
(41)
and consequently
$$L=\sqrt{l(l+1)}\mathrm{}.$$
(42)
The interpretation of this result is that since the orbital quantum number $`l`$ is constrained to take the values $`l=0,1,2,\mathrm{},(n1)`$, the electron can only have orbital momenta $`L`$ specified by means of eq. 42. As in the case of the total energy $`E`$, the angular momentum is conserved and gets quantized. Its natural unit in quantum mechanics is $`\mathrm{}=h/2\pi =1.054\times 10^{34}`$ J.s.
In the macroscopic planetary motion (putting aside the many-body features), the orbital quantum number is so large that any direct experimental detection is impossible. For example, an electron with $`l=2`$ has an angular momentum $`L=2.6\times 10^{34}`$ J.s., whereas the terrestrial angular momentum is $`2.7\times 10^{40}`$ J.s.!
A common notation for the angular momentum states is by means of the letter $`s`$ for $`l=0`$, $`p`$ for $`l=1`$, $`d`$ for $`l=2`$, and so on. This alphabetic code comes from the empirical spectroscopic classification in terms of the so-called series, which was in use before the advent of quantum mechanics.
The combination of the principal quantum number with the latter corresponding to the angular momentum is another frequently used notation in atomic and molecular physics.. For example, a state for which $`n=2`$ and $`l=0`$ is a state $`2s`$, while a state $`n=4`$ and $`l=2`$ is a state $`4d`$.
On the other hand, for the interpretation of the magnetic quantum number, we shall take into account, as we did for the linear momentum, that the orbital momentum is a vector operator and therefore one has to specify its direction, sense, and modulus. $`L`$, being a vector product, is perpendicular on the plane of rotation. The geometric rules of the vectorial products still hold, in particular the rule of the right hand: its direction and sense are given by the right thumb whenever the other four fingers point at the direction of rotation.
But what significance can be associated to a direction and sense in the limited space of the atomic hydrogen ? The answer may be quick if we think that the rotating electron is nothing but a one-electron loop current that considered as a magnetic dipole has a corresponding magnetic field. Consequently, an atomic electron will always interact with an applied magnetic $`B`$. The magnetic quantum number $`m_l`$ specifies the spatial direction of $`L`$, which is determined by the component of $`L`$ along the direction of the external magnetic field. This effect is commonly known as the quantization of the space in a magnetic field.
If we choose the direction of the magnetic field as the $`z`$ axis, the component of $`L`$ along this direction is
$$L_z=m_l\mathrm{}.$$
(43)
The possible values of $`m_l`$ for a given value of $`l`$, go from $`+l`$ to $`l`$, passing through zero, so that there are $`2l+1`$ possible orientations of the angular momentum $`L`$ in a magnetic field. When $`l=0`$, $`L_z`$ can be only zero; when$`l=1`$, $`L_z`$ can be $`\mathrm{}`$, 0, or $`\mathrm{}`$; when $`l=2`$, $`L_z`$ takes only one of the values $`2\mathrm{}`$, $`\mathrm{}`$, 0, $`\mathrm{}`$, or $`2\mathrm{}`$, and so forth. It is worth mentioning that $`L`$ cannot be put exactly parallel or anti-parallel to $`B`$, because $`L_z`$ is always smaller than the modulus $`\sqrt{l(l+1)}\mathrm{}`$ of the total orbital momentum.
The spatial quantization of the orbital momentum for the hydrogen atom is shown in fig. 6.1 in a particular case.
Fig. 6.1: The spatial quantization of the electron angular momentum for states $`l=2`$, $`L=\sqrt{6}\mathrm{}`$.
One should consider the atom/electron characterized by a given $`m_l`$ as having the orientation of its angular momentum $`L`$ determined relative to the external applied magnetic field.
In the absence of the external magnetic field, the direction of the $`z`$ axis is fully arbitrary. Therefore, the component of $`L`$ in any arbitrary chosen direction is $`m_l\mathrm{}`$; the external magnetic field offers a preferred reference direction from the experimental viewpoint.
Why is quantized only the component $`L_z`$ ? The answer is related to the fact that $`L`$ cannot be put along a direction in an arbitrary way. Its ‘vectorial arrow’ moves always along a cone centered on the quantization axis such that its projection $`L_z`$ is $`m_l\mathrm{}`$. The reason why such a phenomenon occurs is due to the uncertainty principle. If $`L`$ would be fixed in space, in such a way that $`L_x`$, $`L_y`$ and $`L_z`$ would have well-defined values, the electron would have to be confined to a well-defined plane. For example, if $`L`$ would be fixed along the $`z`$ direction, the electron tends to maintain itself in the plane $`xy`$ (fig. 6.2a).
Fig. 6.2: The uncertainty principle forbids a fixed direction in space of the angular momentum.
This can occur only in the case in which the component $`p_z`$ of the electron momentum is ‘infinitely’ uncertain. This is however impossible if the electron is part of the hydrogen atom. But since in reality just the component $`L_z`$ of $`L`$ together with $`L^2`$ have well-defined values and $`L>L_z`$, the electron is not constrained to a single plane (fig. 6.2b). If this would be the case, an uncertainty would exist in the coordinate $`z`$ of the electron. The direction of $`L`$ changes continuously (see fig. 6.3), so that the mean values of $`L_x`$ and $`L_y`$ are zero, although $`L_z`$ keeps all the time its value $`m_l\mathrm{}`$.
Fig. 6.3: The angular momentum displays a constant precession around the $`z`$ axis.
The solution for $`\mathrm{\Phi }`$ should also fulfill the normalization condition given by eq. 2. Thus, we have
$$_0^{2\pi }\mathrm{\Phi }^2𝑑\varphi =1$$
(44)
and substituting $`\mathrm{\Phi }`$, one gets
$$_0^{2\pi }A_\varphi ^2𝑑\varphi =1.$$
(45)
It follows that $`A_\varphi =1/\sqrt{2\pi }`$, and thefore the normalized $`\mathrm{\Phi }`$ is
$$\mathrm{\Phi }(\varphi )=\frac{1}{\sqrt{2\pi }}e^{im_l\varphi }.$$
(46)
### Solution for the polar part
The solution of the $`\mathrm{\Theta }(\theta )`$ equation is more complicated. It is expressed in terms of the associated Legendre polynomials
$$P_l^{m_l}(x)=(1)^{m_l}(1x^2)^{m_l/2}\frac{d^{m_l}}{dx^{m_l}}P_l(x)=(1)^{m_l}\frac{(1x^2)^{m_l/2}}{2^ll!}\frac{d^{m_l+l}}{dx^{m_l+l}}(x^21)^l.$$
(47)
Their orthogonality relationship is
$$_1^1[P_l^{m_l}(cos\theta )]^2𝑑cos\theta =\frac{2}{2l+1}\frac{(l+m_l)!}{(lm_l)!}.$$
(48)
For the case of quantum mechanics, $`\mathrm{\Theta }(\theta )`$ is given by the normalized associated Legendre polynomials. Thus, if
$$\mathrm{\Theta }(\theta )=A_\theta P_l^{m_l}(cos\theta ),$$
(49)
then the normalization condition is
$$_1^1A_\theta ^2[P_l^{m_l}(cos\theta )]^2𝑑cos\theta =1.$$
(50)
Therefore, the normalization constant for the polar part is given by
$$A_\theta =\sqrt{\frac{2l+1}{2}\frac{(lm_l)!}{(l+m_l)!}}$$
(51)
and consequently, the function $`\mathrm{\Theta }(\theta )`$ already normalized reads
$$\mathrm{\Theta }(\theta )=\sqrt{\frac{2l+1}{2}\frac{(lm_l)!}{(l+m_l)!}}P_l^{m_l}(cos\theta ).$$
(52)
For our purposes here, the most important property of these functions is that they exist only when the constant $`l`$ is an integer number greater or at least equal to $`m_l`$, which is the absolute value of $`m_l`$. This condition can be written in the form of the set of values available for $`m_l`$
$$m_l=0,\pm 1,\pm 2,\mathrm{},\pm l.$$
(53)
### Unification of the azimuthal and polar parts: spherical harmonics
The solutions of the azimuthal and polar parts can be unified within spherical harmonics functions that depend on both $`\varphi `$ and $`\theta `$. This simplifies the algebraic manipulations of the full wave functions $`\psi (r,\theta ,\varphi )`$. Spherical harmonics are introduced as follows
$$Y_l^{m_l}(\theta ,\varphi )=(1)^{m_l}\sqrt{\frac{2l+1}{4\pi }\frac{(lm_l)!}{(l+m_l)!}}P_l^{m_l}(cos\theta )e^{im_l\varphi }.$$
(54)
The supplementary factor $`(1)^{m_l}`$ does not produce any problem because the Schrödinger equation is linear and homogeneous. This factor is added for the sake of convenience in angular momentum studies. It is known as the Condon-Shortley phase factor and its effect is to introduce an alternance of the signs $`\pm `$ for the spherical harmonics.
### Solution for the radial part
The solution for the radial part $`R(r)`$ of the wave function $`\psi `$ of the hydrogen atom is somewhat more complicated. It is here where significant differences with respect to the electrostatic Laplace equation do occur. The final result is expressed analytically in terms of the associated Laguerre polynomials (Schrödinger 1926). The radial equation can be solved in exact way only when E is positive or for one of the following negative values $`E_n`$ (in which cases, the electron is in a bound stationary state within atomic hydrogen)
$$E_n=\frac{me^4}{32\pi ^2ϵ_0^2\mathrm{}^2}\left(\frac{1}{n^2}\right),$$
(55)
where $`n`$ is an integer number called the principal quantum number. It gives the quantization of the electron energy in the hydrogen atom. This discrete atomic spectrum has been first obtained in 1913 by Bohr using semi-empirical quantization methods and next by Pauli and Schrödinger almost simultaneously in 1926.
Another condition that should be satisfied to solve the radial equation is that $`n`$ have to be strictly bigger than $`l`$. Its lowest value is $`l+1`$ for a givem $`l`$. Vice versa, the condition on $`l`$ is
$$l=0,1,2,\mathrm{},(n1)$$
(56)
for given $`n`$.
The radial equation can be written in the form
$$r^2\frac{d^2R}{dr^2}+2r\frac{dR}{dr}+\left[\frac{2mE}{\mathrm{}^2}r^2+\frac{2me^2}{4\pi ϵ_0\mathrm{}^2}rl(l+1)\right]R=0,$$
(57)
Dividing by $`r^2`$ and using the substitution $`\chi (r)=rR`$ to eliminate the first derivative $`\frac{dR}{dr}`$, one gets the standard form of the radial Schrödinger equation displaying the effective potential $`U(r)=\mathrm{const}/r+l(l+1)/r^2`$ (actually, electrostatic potential plus quantized centrifugal barrier). These are necessary mathematical steps in order to discuss a new boundary condition, since the spectrum is obtained by means of the $`R`$ equation. The difference between a radial Schrödinger equation and a full-line one is that a supplimentary boundary condition should be imposed at the origin ($`r=0`$). The coulombian potential belongs to a class of potentials that are called weak singular for which $`\mathrm{lim}_{r0}=U(r)r^2=0`$. In these cases, one tries solutions of the type $`\chi r^\nu `$, implying $`\nu (\nu 1)=l(l+1)`$, so that the solutions are $`\nu _1=l+1`$ and $`\nu _2=l`$, just as in electrostatics. The negative solution is eliminated for $`l0`$ because it leads to a divergent normalization constant, nor did it respect the normalization at the delta function for the continuous part of the spectrum. On the other hand, the particular case $`\nu _2=0`$ is elmininated because the mean kinetic energy is not finite. The final conclusion is that $`\chi (0)=0`$ for any $`l`$.
Going back to the analysis of the radial equation for $`R`$, first thing to do is to write it in nondimensional variables. This is performed by noticing that the only space and time scales that one can form on combining the three fundamental constants entering this problem, namely $`e^2`$, $`m_e`$ and $`\mathrm{}`$ are the Bohr radius $`a_0=\mathrm{}^2/me^2=0.52910^8`$ cm. and $`t_0=\mathrm{}^3/me^4=0.24210^{16}`$ sec., usually known as atomic units. Employing these units, one gets
$$\frac{d^2R}{dr^2}+\frac{2}{r}\frac{dR}{dr}+\left[2E+\frac{2}{r}\frac{l(l+1)}{r^2}\right]R=0,$$
(58)
where we are especially interested in the discrete part of the spectrum ($`E<0`$). The notations $`n=1/\sqrt{E}`$ and $`\rho =2r/n`$ leads us to
$$\frac{d^2R}{d\rho ^2}+\frac{2}{\rho }\frac{dR}{d\rho }+\left[\frac{n}{\rho }\frac{1}{4}\frac{l(l+1)}{\rho ^2}\right]R=0.$$
(59)
For $`\rho \mathrm{}`$, this equation reduces to $`\frac{d^2R}{d\rho ^2}=\frac{R}{4}`$, having solutions $`Re^{\pm \rho /2}`$. Because of the normalization condition only the decaying exponential is acceptable. On the other hand, the asymptotics at zero, as we already commented on, should be $`R\rho ^l`$. Therefore, we can write $`R`$ as a product of three radial functions $`R=\rho ^le^{\rho /2}F(\rho )`$, of which the first two give the asymptotic behaviors, whereas the third is the radial function in the intermediate region. The latter function is of most interest because its features determine the energy spectrum. The equation for $`F`$ is
$$\rho \frac{d^2F}{d\rho ^2}+(2l+2\rho )\frac{dF}{d\rho }+(nl1)F=0.$$
(60)
This is a particular case of confluent hypergeometric equation for which the two ‘hyper’geometric parameters depend on the pair of quantum numbers $`n,l`$. It can be identified as the equation for the associated Laguerre polynomials $`L_{n+l}^{2l+1}(\rho )`$. Thus, the normalized form of $`R`$ is
$$R_{nl}(r)=\frac{2}{n^2}\sqrt{\frac{(nl1)!}{2n[(n+l)!]^3}}e^{\rho /2}\rho ^lL_{n+l}^{2l+1}(\rho ),$$
(61)
where the following Laguerre normalization condition has been used
$$_0^{\mathrm{}}e^\rho \rho ^{2l}[L_{n+l}^{2l+1}(\rho )]^2\rho ^2𝑑\rho =\frac{2n[(n+l)!]^3}{(nl1)!}.$$
(62)
We have now the solutions of all the equations depending on a single variable and therefore we can build the wave function for any electronic state of the hydrogen atom. The full wave function reads
$$\psi (r,\theta ,\varphi )=𝒩_H(\alpha r)^le^{\alpha r/2}L_{n+l}^{2l+1}(\alpha r)P_l^{m_l}(cos\theta )e^{im_l\varphi },$$
(63)
where $`𝒩_H=\frac{2}{n^2}\sqrt{\frac{2l+1}{4\pi }\frac{(lm_l)!}{(l+m_l)!}\frac{(nl1)!}{[(n+l)!]^3}}`$ and $`\alpha =2/na_0`$.
Using the spherical harmonics, the solution is written as follows
$$\psi (r,\theta ,\varphi )=\frac{2}{n^2}\sqrt{\frac{(nl1)!}{[(n+l)!]^3}}(\alpha r)^le^{\alpha r/2}L_{n+l}^{2l+1}(\alpha r)Y_l^{m_l}(\theta ,\varphi ).$$
(64)
The latter formula may be considered as the final result for the Schrödinger solution of the hydrogen atom for any stationary electron state. Indeed, one can see explicitly both the asmptotic dependence and the two orthogonal and complete sets of functions, i.e., the associated Laguerre polynomials and the spherical harmonics that correspond to this particular case of linear partial second-order differential equation. The parabolic coordinates \[$`\xi =r(1\mathrm{cos}\theta )`$, $`\eta =r(1+\mathrm{cos}\theta )`$, $`\varphi =\varphi `$\], are another coordinate system in which the Schrödinger hydrogen equation is separable (E. Schrödinger, Ann. Physik 80, 437, 1926; P.S. Epstein, Phys. Rev. 28, 695, 1926; I. Waller, Zf. Physik 38, 635, 1926). The final solution in this case is expressed as the product of factors of asymptotic nature, azimuthal harmonics, and two sets of associate Laguerre polynomials in the variables $`\xi `$ and $`\eta `$, respectively. The energy spectrum ($`1/n^2`$) and the degeneracy ($`n^2`$) of course do not depend on the coordinate system.
## Electronic probability density
In the Bohr model of the hydrogen atom, the electron rotates around the nucleus on circular or elliptic trajectories. It is possible to think of appropriate experiments allowing to “see” that the electron moves within experimental errors at the predicted radii $`r=n^2a_0`$ (where $`n`$ is the principal quantum number labeling the orbit and $`a_0=0.53`$ $`\AA `$ is the Bohr radius) in the equatorial plane $`\theta =90^o`$, whereas the azimuthal angle may vary according to the specific experimental conditions.
The more rigorous quantum theory changes the conclusions of the Bohr model in at least two important aspects. First, one cannot speak about exact values of $`r,\theta ,\varphi `$, but only of relative probabilities to find the electron within an infinitesimal given region of space. This feature is a consequence of the wave nature of the electron. Secondly, the electron does not move around the nucleus in the classical conventional way because the probability density $`\psi ^2`$ does not depend on time but can vary substantially as a function of the relative position of the infinitesimal region.
The hydrogenic electron wave function $`\psi `$ is $`\psi =R\mathrm{\Theta }\mathrm{\Phi }`$, where $`R=R_{nl}(r)`$ describes the way $`\psi `$ changes with $`r`$ when the principal and orbital quantum numbers have the values $`n`$ and $`l`$, respectively. $`\mathrm{\Theta }=\mathrm{\Theta }_{lm_l}(\theta )`$ describes in turn how $`\psi `$ varies with $`\theta `$ when the orbital and magnetic quantum numbers have the values $`l`$ and $`m_l`$, respectively. Finally, $`\mathrm{\Phi }=\mathrm{\Phi }_{m_l}(\varphi )`$ gives the change of $`\psi `$ with $`\varphi `$ when the magnetic quantum number has the value $`m_l`$. The probability density $`\psi ^2`$ can be written
$$\psi ^2=R^2\mathrm{\Theta }^2\mathrm{\Phi }^2.$$
(65)
Notice that the probability density $`\mathrm{\Phi }^2`$, which measures the possibility to find the electron at a given azimuthal angle $`\varphi `$, is a constant (does not depend on $`\varphi `$). Therefore, the electronic probability density is symmetric with respect to the $`z`$ axis and independent on the magnetic substates (at least until an external magnetic field is applied). Consequently, the electron has an equal probability to be found in any azimuthal direction. The radial part $`R`$ of the wave function, contrary to $`\mathrm{\Phi }`$, not only varies with $`r`$, but it does it differently for any different combination of quantum numbers $`n`$ and $`l`$. Fig. 6.4 shows plots of $`R`$ as a function of $`r`$ for the states $`1s`$, $`2s`$, and $`2p`$. $`R`$ is maximum at the center of the nucleus ($`r=0`$) for all the $`s`$ states, whereas it is zero at $`r=0`$ for all the states of nonzero angular momentum.
Fig. 6.4: Approximate plots of the radial functions $`R_{1s}`$, $`R_{2s}`$, $`R_{2p}`$; ($`a_0=0.53`$ Å).
Fig. 6.5: Probability density of finding the hydrogenic electron between $`r`$ and $`r+dr`$ with respect to the nucleus for the states $`1s`$, $`2s`$, $`2p`$.
The electronic probability density at the point $`r,\theta ,\varphi `$ is proportional to $`\psi ^2`$, but the real probability in the infinitesimal volume element $`dV`$ is $`\psi ^2dV`$. In spherical coordinates
$$dV=r^2\mathrm{sin}\theta drd\theta d\varphi ,$$
(66)
and since $`\mathrm{\Theta }`$ and $`\mathrm{\Phi }`$ are normalized functions, the real numerical probability $`P(r)dr`$ to find the electron at a relative distance with respect to the nucleus between $`r`$ and $`r+dr`$ is
$`P(r)dr`$ $`=`$ $`r^2R^2dr{\displaystyle _0^\pi }\mathrm{\Theta }^2\mathrm{sin}\theta d\theta {\displaystyle _0^{2\pi }}\mathrm{\Phi }^2𝑑\varphi `$ (67)
$`=`$ $`r^2R^2dr`$
$`P(r)`$ is displayed in fig. 6.5 for the same states for which the radial functions $`R`$ appear in fig. 6.4. In principle, the curves are quite different. We immediately see that $`P(r)`$ is not maximal in the nucleus for the states $`s`$, as happens for $`R`$. Instead, their maxima are encountered at a finite distance from the nucleus. The most probable value of $`r`$ for a $`1s`$ electron is exactly $`a_0`$, the Bohr radius. However, the mean value of $`r`$ for a $`1s`$ electron is $`1.5a_0`$. At first sight this might look strange, because the energy levels are the same both in quantum mechanics and in Bohr’s model. This apparent unmatching is eliminitated if one takes into account that the electron energy depends on $`1/r`$ and not on $`r`$, and the mean value of $`1/r`$ for a $`1s`$ electron is exactly $`1/a_0`$.
The function $`\mathrm{\Theta }`$ varies with the polar angle $`\theta `$ for all the quantum numbers $`l`$ and $`m_l`$, unless $`l=m_l=0`$, which are the $`s`$ states. The probability density $`\mathrm{\Theta }^2`$ for a $`s`$ state is a constant (1/2). This means that since $`\mathrm{\Phi }^2`$ is also a constant, the electronic probability density $`\psi ^2`$ has the same value for a given $`r`$ value, not depending on the direction. In other states, the electrons present an angular behavior that in many cases may be quite complicated. This can be seen in fig.6.5, where the electronic probability densities for different atomic states are displayed as a function of $`r`$ and $`\theta `$. (The plotted term is $`\psi ^2`$ and not $`\psi ^2dV`$). Because $`\psi ^2`$ is independent of $`\varphi `$, a three-dimensional representation of $`\psi ^2`$ can be obtained by rotating a particular representation around a vertical axis. This can prove that the probability densities for the $`s`$ states have spherical symmetry, while all the other states do not possess it. In this way, one can get more or less pronounced lobes of characteristic forms depending on state. These lobes are quite important in chemistry for specifying the atomic interaction in the molecular bulk.
6N. Note:
1. In 1933, E. Schrödinger has been awarded the Nobel Prize in Physics (together with Dirac) for the “discovery of new productive forms of atomic theory”. Schrödinger wrote a remarkable series of four papers “Quantisierung als Eigenwertproblem” \[“Quantization as an eigenvalue problem”\] (I-IV, received by Annalen der Physik on 27 January, 23 February, 10 May and 21 June 1926, respectively).
## 6P. Problems
Problem 6.1 \- Obtain the formulas for the stable orbits and the energy levels of the electron in the atomic hydrogen using only arguments based on the de Broglie wavelength associated to the electron and the empirical value $`5.310^{11}`$ m for the Bohr radius.
Solution: The electron wavelength is given by $`\lambda =\frac{h}{mv}`$, whereas if we equate the electric force and the centripetal force $`\frac{mv^2}{r}=\frac{1}{4\pi ϵ_0}\frac{e^2}{r^2}`$ we obtain the electron ‘velocity’ $`v=\frac{e}{\sqrt{4\pi ϵ_0mr}}.`$Thus, the wavelength of the electron is $`\lambda =\frac{h}{e}\sqrt{\frac{4\pi ϵ_0r}{m}}`$. If we now use the value $`5.3\times 10^{11}`$m for the radius $`r`$ of the electron orbit, we can see that the wavelength of the electron is $`\lambda =33\times 10^{11}`$ m. But this is exactly the same value as of the circumference of the orbit, $`2\pi r=33\times 10^{11}`$ m. One may say that the electron orbit in the atomic hydrogen corresponds to a wave “closing into itself” (i.e., stationary). This fact can be compared to the vibrations of a metallic ring. If the wavelengths are multiples of the circumference, the ring goes on with its vibrations for a long time with very small dissipation If, on the other hand, the number of wavelengths making a circumference is not an integer, the interference of the waves is negative and they dissapear in a short period of time. One may say that the electron will rotate around the nucleus without radiating its energy for an infinite time as far as its orbit contains an integer number of de Broglie wavelengths. Thus, the stability/stationary condition is
$`n\lambda =2\pi r_n,`$
where $`r_n`$ is the radius of the electron orbit containing $`n`$ wavelengths. Substituting $`\lambda `$, we have
$`{\displaystyle \frac{nh}{e}}\sqrt{{\displaystyle \frac{4\pi ϵ_0r_n}{m}}}=2\pi r_n,`$
and therefore the stationary electron orbits are
$`r_n={\displaystyle \frac{n^2\mathrm{}^2ϵ_0}{\pi me^2}}.`$
To get the energy levels, we use $`E=T+V`$ and substituting the kinetic and potential energies leads to
$`E={\displaystyle \frac{1}{2}}mv^2{\displaystyle \frac{e^2}{4\pi ϵ_0r}},`$
or equivalently
$`E_n={\displaystyle \frac{e^2}{8\pi ϵ_0r_n}}.`$
Plugging the value of $`r_n`$ into the latter equation, we get
$`E_n={\displaystyle \frac{me^4}{8ϵ_0^2\mathrm{}^2}}\left({\displaystyle \frac{1}{n^2}}\right).`$
Problem 6.2 \- Unsöld’s theorem tells that for any value of the orbital number $`l`$, the probability densities, summed over all possible substates, from $`m_l=l`$ to $`m_l=+l`$ give a constant that is independent of the angles $`\theta `$ and $`\varphi `$, i.e.
$`{\displaystyle \underset{m_l=l}{\overset{+l}{}}}\mathrm{\Theta }_{lm_l}^2\mathrm{\Phi }_{m_l}^2=ct.`$
This theorem shows that any atom or ion with closed (occupied) sublevels has a spherically-symmetric charge distribution. Check Unsöld’s theorem for $`l=0`$, $`l=1`$, and $`l=2`$.
Solution: For $`l=0`$, $`\mathrm{\Theta }_{00}=1/\sqrt{2}`$ and $`\mathrm{\Phi }_0=1/\sqrt{2\pi }`$, so that
$`\mathrm{\Theta }_{0,0}^2\mathrm{\Phi }_0^2={\displaystyle \frac{1}{4\pi }}.`$
For $`l=1`$, we have
$`{\displaystyle \underset{m_l=1}{\overset{+1}{}}}\mathrm{\Theta }_{lm_l}^2\mathrm{\Phi }_{m_l}^2=\mathrm{\Theta }_{1,1}^2\mathrm{\Phi }_1^2+\mathrm{\Theta }_{1,0}^2\mathrm{\Phi }_0^2+\mathrm{\Theta }_{1,1}^2\mathrm{\Phi }_1^2.`$
On the other hand, the wave functions are given by $`\mathrm{\Theta }_{1,1}=(\sqrt{3}/2)sin\theta `$, $`\mathrm{\Phi }_1=(1/\sqrt{2\pi })e^{i\varphi }`$, $`\mathrm{\Theta }_{1,0}=(\sqrt{6}/2)cos\theta `$, $`\mathrm{\Phi }_0=1/\sqrt{2\pi }`$, $`\mathrm{\Theta }_{1,1}=(\sqrt{3}/2)sin\theta `$, $`\mathrm{\Phi }_1=(1/\sqrt{2\pi })e^{i\varphi }`$ , which plugged into the previous equation give
$`{\displaystyle \underset{m_l=1}{\overset{+1}{}}}\mathrm{\Theta }_{lm_l}^2\mathrm{\Phi }_{m_l}^2={\displaystyle \frac{3}{8\pi }}sin^2\theta +{\displaystyle \frac{3}{4\pi }}cos^2\theta +{\displaystyle \frac{3}{8\pi }}sin^2\theta ={\displaystyle \frac{3}{4\pi }}`$
and again we’ve got a constant.
For $`l=2`$, we have
$$\underset{m_l=2}{\overset{+2}{}}\mathrm{\Theta }_{lm_l}^2\mathrm{\Phi }_{m_l}^2=\mathrm{\Theta }_{2,2}^2\mathrm{\Phi }_2^2\mathrm{\Theta }_{2,1}^2\mathrm{\Phi }_1^2$$
$$+\mathrm{\Theta }_{2,0}^2\mathrm{\Phi }_0^2+\mathrm{\Theta }_{2,1}^2\mathrm{\Phi }_1^2+\mathrm{\Theta }_{2,2}^2\mathrm{\Phi }_2^2,$$
and the wave functions are $`\mathrm{\Theta }_{2,2}=(\sqrt{15}/4)sin^2\theta `$, $`\mathrm{\Phi }_2=(1/\sqrt{2\pi })e^{2i\varphi }`$, $`\mathrm{\Theta }_{2,1}=(\sqrt{15}/2)sin\theta cos\theta `$, $`\mathrm{\Phi }_1=(1/\sqrt{2\pi })e^{i\varphi }`$, $`\mathrm{\Theta }_{2,0}=(\sqrt{10}/4)(3cos^2\theta 1)`$, $`\mathrm{\Phi }_0=1/\sqrt{2\pi }`$, $`\mathrm{\Theta }_{2,1}=(\sqrt{15}/2)sin\theta cos\theta `$, $`\mathrm{\Phi }_1=(1/\sqrt{2\pi })e^{i\varphi }`$, $`\mathrm{\Theta }_{2,2}=(\sqrt{15}/4)sin^2\theta `$, $`\mathrm{\Phi }_2=(1/\sqrt{2\pi })e^{2i\varphi }`$, Plugging them into the previous equation give
$`{\displaystyle \underset{m_l=2}{\overset{+2}{}}}\mathrm{\Theta }_{lm_l}^2\mathrm{\Phi }_{m_l}^2={\displaystyle \frac{5}{4\pi }},`$
which again fulfills Unsöld’s theorem.
Problem 6.3 \- The probability to find an atomic electron whose radial wave functions is that of the ground state $`R_{10}(r)`$ outside a sphere of Bohr radius $`a_0`$ centered on the nucleus is
$`{\displaystyle _{a_0}^{\mathrm{}}}R_{10}(r)^2r^2𝑑r.`$
Obtain the probability to find the electron in the ground state at a distance from the nucleus bigger than $`a_0`$.
Solution: The radial wave function corresponding to the ground state is
$`R_{10}(r)={\displaystyle \frac{2}{a_0^{3/2}}}e^{r/a_0}.`$
Substituting it in the integral, we get $`_{a_0}^{\mathrm{}}R(r)^2r^2𝑑r=\frac{4}{a_0^3}_{a_0}^{\mathrm{}}r^2e^{2r/a_0}𝑑r,`$or
$`{\displaystyle _{a_0}^{\mathrm{}}}R(r)^2r^2𝑑r={\displaystyle \frac{4}{a_0^3}}\left[{\displaystyle \frac{a_0}{2}}r^2e^{2r/a_0}{\displaystyle \frac{a_0^2}{2}}re^{2r/a_0}{\displaystyle \frac{a_0^3}{4}}e^{2r/a_0}\right]_{a_0}^{\mathrm{}}.`$
This leads us to
$`{\displaystyle _{a_0}^{\mathrm{}}}R(r)^2r^2𝑑r={\displaystyle \frac{5}{e^2}}68\%!!,`$
which is the result asked for in this problem.
7. QUANTUM SCATTERING
## Introduction
One usually begins the quantum theory of scattering by referring to results already known from the classical scattering in central fields with some simplifying assumptions helping to avoid unnecessary calculations in getting basic results. It is generally known that studying scatterings in the laboratory provides information on the distribution of matter in the target and other details of the interaction between the incident beam and the target. The hypotheses that we shall assume correct in the following are
i) The particles are spinless. This, of course, does not mean that spin effects are not important in quantum scatterings.
ii) We shall study only elastic scattering for which the internal structure of the particles is not taken into account.
iii) The target is sufficiently thin to neglect multiple scatterings.
iv) The interactions are described by a potential that depends only on the relative distance between the particles (central potential).
These hypotheses eliminate some quantum effects that are merely details. They also represent conditions for getting the quantum analogs of basic classical results. We now define
$$\frac{d\sigma }{d\mathrm{\Omega }}\frac{I(\theta ,\phi )}{I_0},$$
(1)
where $`d\mathrm{\Omega }`$ is the solid angle infinitesimal element, $`I_0`$ is the number of incident particles per unit transverse area, and $`Id\mathrm{\Omega }`$ is the number of scattered particles in the solid angle element.
Employing these well-known concepts, together with the asymptotic notion of impact parameter $`b`$ associated to each classical incident particle, one gets in classical mechanics the following important formula
$$\frac{d\sigma }{d\mathrm{\Omega }}=\frac{b}{\mathrm{sin}\theta }|\frac{db}{d\theta }|.$$
(2)
If one wants to study the scattering phenomenology in quantum terminology, one should investigate the time evolution of a ‘scattering’ wave packet. Let $`F_i`$ be the flux of incident particles, i.e., the number of particles per unit of time passing through the unit of transverse surface onto the propagation axis. An appropriate detector configuration is usually placed far away from the effective interaction region, ‘seeing’ a solid angle $`d\mathrm{\Omega }`$ of that region. In general, the number of particles $`dn/dt`$ scattered per unit of time in $`d\mathrm{\Omega }`$ in the direction $`(\theta ,\phi )`$ is detected.
$`dn/dt`$ is proportional to $`d\mathrm{\Omega }`$ and $`F_i`$. Let us call $`\sigma (\theta ,\phi )`$ the coefficient of proportionality between $`dn`$ and $`F_id\mathrm{\Omega }`$:
$$dn=\sigma (\theta ,\phi )F_id\mathrm{\Omega },$$
(3)
which is by definition the differential cross section.
The number of particles per unit of time reaching the detector is equal to the number of particles crossing the surface $`\sigma (\theta ,\phi )d\mathrm{\Omega }`$, which is perpendicular to the beam axis. The total section is by definition
$$\sigma =\sigma (\theta ,\phi )𝑑\mathrm{\Omega }.$$
(4)
To further simplify the calculation, we choose the z axis along the incident beam direction.
On the negative side of the axis, for large negative $`t`$, the particle is practically free: it is not affected by $`V(𝐫)`$ and its state can be represented by plane waves. Therefore, the wave function contains terms of the form $`e^{ikz}`$, where $`k`$ is the constant ocurring in the Helmholtz equation. By analogy with optics, the form of the scattered wave is
$$f(r)=\frac{e^{ikr}}{r}.$$
(5)
Indeed
$$(^2+k^2)e^{ikr}0$$
(6)
and
$$(^2+k^2)\frac{e^{ikr}}{r}=0$$
(7)
for $`r>r_0`$, where $`r_0`$ is any positive number.
We assume that the motion of the particle is described by the Hamiltonian
$$H=\frac{𝐩^\mathrm{𝟐}}{2\mu }+V=H_0+V.$$
(8)
$`V`$ is different of zero only in a small neighbourhood close to the origin. A wave packet at $`t=0`$ can be written
$$\psi (𝐫,0)=\frac{1}{(2\pi )^{\frac{3}{2}}}\phi (𝐤)\mathrm{exp}[i𝐤(𝐫𝐫_\mathrm{𝟎})]𝐝^\mathrm{𝟑}𝐤,$$
(9)
where $`\psi `$ is a function that is nonzero in a ‘width’ $`\mathrm{\Delta }𝐤`$ centered on $`𝐤_\mathrm{𝟎}`$. We also assume that $`𝐤_\mathrm{𝟎}`$ is antiparallel to $`𝐫_\mathrm{𝟎}`$. In order to see quantitatively what happens to the wave packet when scatters the target, one can use the expansion of $`\psi (𝐫,0)`$ in the eigenfunctions $`\psi _n(𝐫)`$ of $`H`$, i.e., $`\psi (𝐫,0)=_nc_n\psi _n(𝐫)`$. Thus, the wave packet at time $`t`$ is
$$\psi (𝐫,t)=\underset{n}{}c_n\phi _n(𝐫)\mathrm{exp}(\frac{i}{\mathrm{}}E_nt).$$
(10)
This is an eigenfunction of the operator $`H_0`$, not of $`H`$, but we can substitute these eigenfunctions by eigenfunctions of $`H`$, which we denote by $`\psi _k^{(+)}(𝐫)`$. The asymptotic form of the latter is
$$\psi _k^{(+)}(𝐫)𝐞^{\mathrm{𝐢𝐤}𝐫}+𝐟(𝐫)\frac{𝐞^{\mathrm{𝐢𝐤𝐫}}}{|𝐫|},$$
(11)
where, as usually $`𝐩=\mathrm{}𝐤`$ and $`E=\frac{\mathrm{}^2k^2}{2m}`$.
This corresponds to a plane wave of the incident beam type and a divergent spherical wave resulting from the interaction between the incident beam and the target. One can expand $`\psi (𝐫,0)`$ in plane waves and $`\psi _k(𝐫)`$
$$\psi (𝐫,0)=\phi (𝐤)\mathrm{exp}(i𝐤𝐫_\mathrm{𝟎})\psi _𝐤(𝐫)d^3k,$$
(12)
where $`\mathrm{}\omega =\frac{\mathrm{}^2k^2}{2m}`$. The divergent spherical wave does not contribute to the initial wave packet because it is an additive part.
## Scattering of a wave packet
Any wave is dispersed during its propagation. This is why one cannot ignore the effect of the divergent wave from this viewpoint. One can make use of the following trick
$$\omega =\frac{\mathrm{}}{2m}k^2=\frac{\mathrm{}}{2m}[𝐤_\mathrm{𝟎}+(𝐤𝐤_\mathrm{𝟎})]^2=\frac{\mathrm{}}{2m}[2𝐤_\mathrm{𝟎}𝐤𝐤_\mathrm{𝟎}^\mathrm{𝟐}+(𝐤𝐤_\mathrm{𝟎})^\mathrm{𝟐}],$$
(13)
pentru a neglija ultimul termen în paranteze. Substituting $`\omega `$ in $`\psi `$, we ask that $`\frac{\mathrm{}}{2m}(𝐤𝐤_\mathrm{𝟎})^2T1`$, where $`T\frac{2mr_0}{\mathrm{}k_0}`$. Therefore
$$\frac{(\mathrm{\Delta }k)^2r_0}{k_0}1.$$
(14)
This condition tells us that the wave packet does not disperse significantly even when it moves over amacroscopic distance $`r_0`$.
Choosing the direction of the vector $`𝐤`$ of the incident wave along one of the three cartesian directions (we use the $`z`$ one), we can write in spherical coordinates the following important formula
$`\psi _k(r,\theta ,\phi )e^{ikz}+\frac{f(k,\theta ,\phi )e^{ikr}}{r}.`$
Since the Hamiltonian $`H`$, up to now not considered as an operator (the class of the results presented are the same both at the classical and quantum level), is invariant under $`z`$ rotations, we can choose boundary conditions of spherical symmetry too. Thus
$`\psi _k(r,\theta ,\phi )e^{ikz}+\frac{f(\theta )e^{ikr}}{r}.`$
This type of functions are known as scattering wave functions. The coefficient $`f(\theta )`$ of the spherical wave is known as the scattering amplitude. It is a basic concept in the formal theory of quantum scatterings.
## Probability amplitude in scattering
We write the Schrödinger equation as follows
$$i\mathrm{}\frac{\psi }{t}=\frac{\mathrm{}^2}{2m}^2\psi +V(𝐫,t)\psi .$$
(15)
Recall that the expression
$$P(𝐫,t)=\psi ^{}(𝐫,t)\psi (𝐫,t)=|\psi (𝐫,t)|^2$$
(16)
can be interpreted, cf. Max Born, as a probability density under normalization conditions of the type
$$|\psi (𝐫,t)|^2d^3r=1.$$
(17)
This normalization integral should be time independent. This can be noted by writing
$$I=\frac{}{t}_\mathrm{\Omega }P(𝐫,t)d^3r=_\mathrm{\Omega }(\psi ^{}\frac{\psi }{t}+\frac{\psi ^{}}{t}\psi )d^3r,$$
(18)
and from Schrödinger’s equation
$$\frac{\psi }{t}=\frac{i\mathrm{}}{2m}^2\psi \frac{i}{\mathrm{}}V(𝐫,t)\psi $$
(19)
one gets
$$I=\frac{i\mathrm{}}{2m}_\mathrm{\Omega }[\psi ^{}^2(^2\psi ^{})\psi ]d^3r=\frac{i\mathrm{}}{2m}_\mathrm{\Omega }[\psi ^{}\psi (\psi ^{})\psi ]d^3r=$$
$$=\frac{i\mathrm{}}{2m}_A[\psi ^{}\psi (\psi ^{})\psi ]_n𝑑A,$$
(20)
where the Green theorem has been used to evaluate the volume integral. $`dA`$ is the infinitesimal surface element on the boundary of the integration region and $`[]_n`$ denotes the component along the normal direction to the surface element $`dA`$.
Defining
$$𝐒(𝐫,t)=\frac{\mathrm{}}{2im}[\psi ^{}\psi (\psi ^{})\psi ],$$
(21)
we get
$$I=\frac{}{t}_\mathrm{\Omega }P(𝐫,t)d^3r=_\mathrm{\Omega }𝐒d^3r=_AS_n𝑑A,$$
(22)
for well-bahaved wave packets (not funny asymptotically) so that the normalization integral converges. The surface integral is zero when $`\mathrm{\Omega }`$ covers the whole space. One can prove (see P. Dennery & A. Krzywicki, Mathematical methods for physicists) that the surface integral is zero. Therefore, the normalization integral is constant in time and the initial condition holds. From the same equation for $`𝐒`$, we get
$$\frac{P(𝐫,t)}{t}+𝐒(𝐫,t)=0,$$
(23)
which is the continuity equation for the density flux $`P`$ and the current density $`𝐒`$ in the absence of any type of sources or sinks. If we interpret $`\frac{\mathrm{}}{im}`$ as a sort of velocity ‘operator’ (as for time, it is difficult to speak rigorously about a velocity operator in quantum mechanics!), then
$$𝐒(𝐫,t)=Re(\psi ^{}\frac{\mathrm{}}{im}\psi ).$$
(24)
To calculate the quantum current density for a scattering wave function is a tricky and inspiring (not illustrative) exercise! The final result is $`j_r=\frac{\mathrm{}k}{mr^2}|f(\theta )|^2`$, where the direction $`\theta =0`$ should not be included.
## Green’s function in scattering theory
Another way of writing the Schrödinger equation at hand is $`(\frac{\mathrm{}^2}{2m}^2+V)\psi =E\psi `$, or $`(^2+k^2)\psi =U\psi `$, where $`k^2=\frac{2mE}{\mathrm{}^2}`$, şi $`U=\frac{2mV}{\mathrm{}^2}`$.
It follows that it is more convenient to put this equation in an integral form. This can be done if we consider $`U\psi `$ in the right hand side of the equation as a inhomogeneity. This allows to build the solution by means of Green’s function (integral kernel), which, by definition, is the solution of
$$(^2+k^2)G(𝐫,𝐫^{})=\delta (𝐫𝐫^{}).$$
(25)
One can write now the Schrödinger solution as the sum of the homogeneous equation and the inhomogeneous one of Green’s type
$$\psi (𝐫)=\lambda (𝐫)𝐆(𝐫,𝐫^{})𝐔(𝐫^{})\psi (𝐫^{})𝐝^\mathrm{𝟑}𝐫^{}.$$
(26)
We seek now a $`G`$ function in the form of a product of linear independent functions, for example, plane waves
$$G(𝐫,𝐫^{}=A(𝐪)e^{i𝐪(𝐫𝐫^{})}dq.$$
(27)
Using eq. 25, we have
$$A(𝐪)(k^2q^2)e^{i𝐪(𝐫𝐫^{})}𝑑q=\delta (𝐫𝐫^{}),$$
(28)
which turns in an identity if
$$A(𝐪)=(2\pi )^3(k^2q^2)^1.$$
(29)
Thus
$$G(𝐫,𝐫^{})=\frac{1}{(2\pi )^3}\frac{e^{iqR}}{k^2q^2}d^3q,$$
(30)
where $`R=|𝐫𝐫^{}|`$. After performing a calculation of complex variable <sup>10</sup><sup>10</sup>10See problem 7.1., we get
$$G(r)=\frac{1}{4\pi }\frac{e^{ikr}}{r}.$$
(31)
This function is not determined univoquely since the Green function can be any solution of the eq. 25. The right particular solution is chosen by imposing boundary conditions on the eigenfunctions $`\psi _k(𝐫)`$.
The Green function obtained in this way is
$$G(𝐫,𝐫^{})=\left(\frac{e^{ik|𝐫𝐫^{}|}}{4\pi |𝐫𝐫^{}|}\right).$$
(32)
Thus, we finally get the integral equation for the scattering wave function
$$\psi (k,𝐫)=\phi (k,𝐫)\frac{m}{2\pi \mathrm{}^2}\frac{e^{ik|𝐫𝐫^{}|}}{𝐫𝐫^{}}U(𝐫^{})\psi (k,𝐫)𝑑𝐫,$$
(33)
where $`\phi `$ is a solution of the Helmholtz equation. Noticing that $`|𝐫𝐫^{}|=R`$, then
$$(^2+k^2)\psi =(^2+k^2)[\phi +G(𝐫,𝐫^{})U(𝐫^{})\psi (𝐫^{})d^3r^{}]$$
(34)
and assuming that we can change the order of operations and put the $``$ operator inside the integral, we get
$$(^2+k^2)\psi =(^2+k^2)G(𝐫,𝐫^{})U(𝐫^{})\psi (𝐫^{})d^3r^{}=U(𝐫)\psi (𝐫),$$
(35)
which shows us that $`G(R)=\frac{1}{4\pi }\frac{e^{ikR}}{R}`$ is indeed a solution.
## Optical theorem
The total cross section is given by
$$\sigma _{tot}(k)=\frac{d\sigma }{d\mathrm{\Omega }}𝑑\mathrm{\Omega }.$$
(36)
Let us express now $`f(\theta )`$ as a function of the phase shift $`S_l(k)=e^{2i\delta _l(k)}`$ in the form
$$f(\theta )=\frac{1}{k}\underset{l=0}{\overset{\mathrm{}}{}}(2l+1)e^{i\delta _i(k)}\mathrm{sin}\delta _l(k)P_l(\mathrm{cos}\theta ).$$
(37)
Then
$$\sigma _{tot}=[\frac{1}{k}\underset{l=0}{\overset{\mathrm{}}{}}(2l+1)e^{i\delta _l(k)}\mathrm{sin}\delta _l(k)P_l(\mathrm{cos}\theta )]$$
$$[[\frac{1}{k}\underset{l^{}=0}{\overset{\mathrm{}}{}}(2l^{}+1)e^{i\delta _l^{}(k)}\mathrm{sin}\delta _l^{}(k)P_l^{}(\mathrm{cos}\theta )].$$
(38)
Using now $`P_l(\mathrm{cos}\theta )P_l^{}(\mathrm{cos}\theta )=\frac{4\pi }{2l+1}\delta _{ll^{}}`$, we get
$$\sigma _{tot}=\frac{4\pi }{k^2}\underset{l=0}{\overset{\mathrm{}}{}}(2l+1)\mathrm{sin}\delta _l(k)^2.$$
(39)
Of interest is the relationship
$$\mathrm{Im}f(0)=\frac{1}{k}\underset{l=0}{\overset{\mathrm{}}{}}(2l+1)\mathrm{Im}[e^{i\delta _l(k)}\mathrm{sin}\delta _l(k)]P_l(1)=\frac{1}{k}\underset{l=0}{\overset{\mathrm{}}{}}(2l+1)\mathrm{sin}\delta _l(k)^2=$$
$$\frac{k}{4\pi }\sigma _{tot},$$
(40)
which is known as the optical theorem. Its physical significance is related to the fact that the interference of the incident wave with the dispersed wave at zero/forward angle produces the “getting out” of the particle from the incident wave, allowing in this way the conservation of the probability.
## Born approximation
Let us consider the situation of Fig. 7.2:
The observation point M is far away from P, which is in the range of the potential $`U`$. The geometrical conditions are $`rL`$, $`r^{}l`$. The segment MP that corresponds to $`|𝐫𝐫^{}|`$ is in the aforementioned geometrical conditions approxiamtely equal to the projection of MP onto MO
$$|𝐫𝐫^{}|r𝐮𝐫^{},$$
(41)
where $`𝐮`$ is a unit vector (versor) in the $`𝐫`$ direction. Then, for large $`r`$
$$G=\frac{1}{4\pi }\frac{e^{ik|𝐫𝐫^{}|}}{|𝐫𝐫^{}|}_r\mathrm{}\frac{1}{4\pi }\frac{e^{ikr}}{r}e^{ik𝐮𝐫}.$$
(42)
We now substitute $`G`$ in the integral expression for the scattering wave function
$$\psi (𝐫)=e^{ikz}\frac{1}{4\pi }\frac{e^{ikr}}{r}e^{ik𝐮𝐫}U(𝐫^{})\psi (𝐫^{})d^3r^{}.$$
(43)
The latter is already not a function of the distance $`r=OM`$, but only of $`\theta `$ and $`\psi `$. Thus
$$f(\theta ,\psi )=\frac{1}{4\pi }e^{ik𝐮𝐫}U(𝐫^{})\psi (𝐫^{})d^3r^{}.$$
(44)
We define now the incident wave vector $`𝐤_𝐢`$ as a vector of modulus $`k`$ directed along the polar axis of the beam. Then $`e^{ikz}=e^{i𝐤_𝐢𝐫}`$. Similarly, $`𝐤_𝐝`$, of modulus $`k`$ and of direction fixed by $`\theta `$ and $`\phi `$, is called the shifted wave vector in the direction $`(\theta ,\phi )`$: $`𝐤_𝐝=k𝐮`$.
The momentum transfer in the direction $`(\theta ,\phi )`$ is introduced as the vectorial difference $`𝐊=𝐤_𝐝𝐤_𝐢`$.
Hence we can write the integral equation in the form
$$\psi (𝐫)=e^{i𝐤_𝐢𝐫}+G(𝐫,𝐫^{})U(𝐫^{})\psi (𝐫^{})d^3r^{}$$
(45)
One can try to solve this equation iteratively. Putting $`𝐫𝐫^{};𝐫^{}𝐫^{\prime \prime }`$, we can write
$$\psi (𝐫^{})=e^{i𝐤_𝐢𝐫^{}}+G(𝐫^{},𝐫^{\prime \prime })U(𝐫^{\prime \prime })\psi (𝐫^{\prime \prime })d^3r^{\prime \prime }.$$
(46)
Substituting in 45, we get
$$\psi (𝐫)=e^{i𝐤_ir}+G(𝐫,𝐫^{})U(𝐫^{})e^{i𝐤_𝐢𝐫^{}}d^3r^{}+$$
$$G(𝐫,𝐫^{})U(𝐫^{})G(𝐫^{},𝐫^{\prime \prime })U(𝐫^{\prime \prime })\psi (𝐫^{\prime \prime })d^3r^{\prime \prime }d^3r^{}.$$
(47)
The first two terms in the right hand side are known and it is only the third one that includes the unknown function $`\psi (𝐫)`$. We can repeat the procedure: substituting $`𝐫`$ by $`𝐫^{\prime \prime }`$, and $`𝐫^{}`$ by $`𝐫^{\prime \prime \prime }`$, we get $`\psi (𝐫^{\prime \prime })`$ , that we can reintroduce in the eq. 47
$$\psi (𝐫)=e^{i𝐤_𝐢𝐫}+G(𝐫,𝐫^{})U(𝐫^{})e^{i𝐤_𝐢𝐫^{}}+$$
$$G(𝐫,𝐫^{})U(𝐫^{})G(𝐫^{},𝐫^{\prime \prime })U(𝐫^{\prime \prime })e^{i𝐤_𝐢𝐫^{\prime \prime }}d^3r^{}d^3r^{\prime \prime }+$$
$$G(𝐫,𝐫^{})U(𝐫^{})G(𝐫^{},𝐫^{\prime \prime })U(𝐫^{\prime \prime })e^{i𝐤_𝐢𝐫^{\prime \prime }}G(𝐫^{\prime \prime },𝐫^{\prime \prime \prime })U(𝐫^{\prime \prime \prime })\psi (𝐫^{\prime \prime \prime }).$$
(48)
The first three terms are now known and the unknown function $`\psi (𝐫)`$ has been sent to the fourth term. In this way, by succesive iterations we can build the stationary dispersed wave function. Notice that each term of the series expansion has one more power in the potential with respect to the previous one. We can go on until we get a negligible expression in the right hand side, obtaining $`\psi (𝐫)`$ as a function of only known quantities.
Substituting the expression of $`\psi (𝐫)`$ in $`f(\theta ,\phi )`$, we get the expansion in Born series of the scattering amplitude. In first order in $`U`$, one should replace $`\psi (𝐫^{})`$ by $`e^{i𝐤_𝐢𝐫^{}}`$ in the right hand side to get
$$f^{(B)}(\theta ,\phi )=\frac{1}{4\pi }e^{i𝐤_𝐢𝐫^{}}U(𝐫^{})e^{ik𝐮𝐫^{}}d^3r^{}=\frac{1}{4\pi }e^{i(𝐤_𝐝𝐤_𝐢)𝐫^{}}U(𝐫^{})d^3r^{}=$$
$$\frac{1}{4\pi }e^{i𝐊𝐫^{}}U(𝐫^{})d^3r^{}$$
(49)
$`𝐊`$ is the momentum transfer vector. Thus, the differential cross section is simply related to the potential, $`V(𝐫)=\frac{\mathrm{}^2}{2m}U(𝐫)`$. Since $`\sigma (\theta ,\phi )=|f(\theta ,\phi )|^2`$, the result is
$$\sigma ^{(B)}(\theta ,\phi )=\frac{m^2}{4\pi ^2\mathrm{}^4}|e^{i𝐊𝐫}V(𝐫)d^3r|^2$$
(50)
The direction and modulus of $`𝐊`$ depends on the modulus $`k`$ of $`𝐤_𝐢`$ and $`𝐤_𝐝`$ as well as on the scattering direction $`(\theta ,\phi )`$. For given $`\theta `$ and $`\phi `$, it is a function of $`k`$, the energy of the incident beam. Analogously, for a given energy, $`\sigma ^{(B)}`$ is a function of $`\theta `$ and $`\phi `$. Born’s approximation allows one to get information on the potential $`V(𝐫)`$ from the dependence of the differential cross section on the scattering direction and the incident energy.
7N. Note \- The following paper of Born was practically the first dealing with quantum scattering:
M. Born, “Quantenmechanik der Stossvorgänge” \[“Quantum mechanics of scattering processes ”\], Zf. f. Physik 37, 863-867 (1926)
## 7P. Problems
Problem 7.1
Calculus of complex variable for the scattering Green function
We recall that we already obtained the result
$`G(𝐫,𝐫^{})=\frac{1}{(2\pi )^3}\frac{e^{iqR}}{k^2q^2}d^3q,`$ cu $`R=|𝐫𝐫^{}|`$. Since $`d^3q=q^2\mathrm{sin}\theta dqd\theta d\varphi `$, we get after integrating in angular variables
$`G(𝐫,𝐫^{})=\frac{i}{4\pi ^2R}_{\mathrm{}}^{\mathrm{}}\frac{(e^{iqR}e^{iqR})}{k^2q^2}q𝑑q.`$
Putting $`C=\frac{i}{4\pi ^2R}`$, we separate the integral in two parts
$`C(_{\mathrm{}}^{\mathrm{}}\frac{e^{iqR}}{k^2q^2}q𝑑q_{\mathrm{}}^{\mathrm{}}\frac{e^{iqR}}{k^2q^2}q𝑑q).`$
Let us make now $`qq`$ in the first integral
$`_{\mathrm{}}^{\mathrm{}}\frac{e^{i(q)R}}{k^2(q)^2}(q)d(q)=_{\mathrm{}}^{\mathrm{}}\frac{e^{iqR}}{k^2q^2}q𝑑q=_{\mathrm{}}^{\mathrm{}}\frac{e^{iqR}}{k^2q^2}q𝑑q,`$
so that
$`G(𝐫,𝐫^{})=2C(_{\mathrm{}}^{\mathrm{}}\frac{qe^{iqR}}{k^2q^2}𝑑q).`$
Substituting $`C`$, leads to
$`G(𝐫,𝐫^{})=\frac{i}{2\pi ^2R}_{\mathrm{}}^{\mathrm{}}\frac{qe^{iqR}}{k^2q^2}𝑑q`$
In this form, the integral can be calculated by means of the theorem of residues of its poles. Notice the presence of simple poles at $`q=_{}^+k`$.
Fig. 7.4: Contour rules around the poles for $`G_+`$ and $`G_{}`$
We use the contour of fig. 7.4 encircling the poles as shown, because in this way we get the physically correct effect from the theorem of residues
$`G(r)=\frac{1}{4\pi }\frac{e^{ikr}}{r}(\mathrm{Im}k>0)`$ ,
$`G(r)=\frac{1}{4\pi }\frac{e^{ikr}}{r}(\mathrm{Im}k<0).`$
The solution of interest is the first one, because it provides divergent waves, whereas the latter solution holds for convergent waves (propagating towards the target). Moreover, the linear combination
$`\frac{1}{2}lim_{ϵ0}[G_{k+iϵ}+G_{kiϵ}]=\frac{1}{4\pi }\frac{\mathrm{cos}kr}{r}`$
corresponds to stationary waves.
The formal calculation of the integral can be performed by taking $`k^2q^2k^2+iϵq^2`$, so that: $`_{\mathrm{}}^{\mathrm{}}\frac{qe^{iqR}}{k^2q^2}𝑑q_{\mathrm{}}^{\mathrm{}}\frac{qe^{iqR}}{(k^2+iϵ)q^2}𝑑q.`$ This is possible for $`R>0`$. This is why the contour for the calculation will be placed in the upper half plane. Thus, the poles of the integrand are located at $`q=\pm \sqrt{k^2+iϵ}\pm (k+\frac{iϵ}{2k})`$. The procedure of taking the limit $`ϵ0`$ should be applied after calculating the integral.
Problem 7.2
Asymptotic form of the radial function
As we have already seen in the chapter Hydrogen atom the radial part of the Schrödinger equation can be written
$`(\frac{d^2}{dr^2}+\frac{2}{r}\frac{d}{dr})R_{nlm}(r)\frac{2m}{\mathrm{}^2}[V(r)+\frac{l(l+1)\mathrm{}^2}{2mr^2}]R_{nlm}(r)+\frac{2mE}{\mathrm{}^2}R_{nlm}(r)=0.`$
$`n,l,m`$ are the spherical quantum numbers. For the sake of convenience of writing we shall discard them hereafter. $`R`$ is the radial wave function (i.e., depends only on $`r`$). We assume that the potential goes to zero stronger than $`1/r`$, and that $`lim_{r0}r^2V(r)=0`$.
Using $`u(r)=rR`$, since $`(\frac{d^2}{dr^2}+\frac{2}{r}\frac{d}{dr})\frac{u}{r}=\frac{1}{r}\frac{d^2}{dr^2}u`$, we have
$`\frac{d^2}{dr^2}u+\frac{2m}{\mathrm{}^2}[EV(r)\frac{l(l+1)\mathrm{}^2}{2mr^2}]u=0.`$
Notice that the potential displays a supplementary term
$`V(r)V(r)+\frac{l(l+1)\mathrm{}^2}{2mr^2},`$
which corresponds to a repulsive centrifugal barrier. For a free particle $`V(r)=0`$, and the equation becomes
$`[\frac{d^2}{dr^2}+\frac{2}{r}\frac{d}{dr})\frac{l(l+1)}{r^2}]R+k^2R=0.`$
Introducing the variable $`\rho =kr`$, we get
$`\frac{d^2R}{d\rho ^2}+\frac{2}{\rho }\frac{dR}{d\rho }\frac{l(l+1)}{\rho ^2}R+R=0.`$
The solutions are the so-called spherical Bessel functions. The regular solution is
$`j_l(\rho )=(\rho )^l(\frac{1}{\rho }\frac{d}{d\rho })^l(\frac{\mathrm{sin}\rho }{\rho }),`$
while the irregular one
$`n_l(\rho )=(\rho )^l(\frac{1}{\rho }\frac{d}{d\rho })^l(\frac{\mathrm{cos}\rho }{\rho }).`$
For large $`\rho `$, the functions of interest are the spherical Hankel functions
$`h_l^{(1)}(\rho )=j_l(\rho )+in_l(\rho )`$ şi $`h_l^{(2)}(\rho )=[h_l^{(1)}(\rho )]^{}.`$
The behaviour for $`\rho l`$ is of special interest
$$j_l(\rho )\frac{1}{\rho }\mathrm{sin}(\rho \frac{l\pi }{2})$$
(51)
$$n_l(\rho )\frac{1}{\rho }\mathrm{cos}(\rho \frac{l\pi }{2}).$$
(52)
Then
$`h_l^{(1)}\frac{i}{\rho }e^{i(\rho l\pi /2)}.`$
The solution regular at the origin is $`R_l(r)=j_l(kr).`$
The asymptotic form is (using eq. 51)
$`R_l(r)\frac{1}{2ikr}[e^{ikrl\pi /2}e^{ikrl\pi /2}].`$
Problem 7.3
Born approximation for Yukawa potentials
Let us consider the potential of the form
$$V(𝐫)=V_0\frac{e^{\alpha r}}{r},$$
(53)
where $`V_0`$ and $`\alpha `$ are real constants and $`\alpha `$ is positive. The potential is either attractive or repulsive depending on the sign of $`V_0`$; the larger $`|V_0|`$, the stronger the potential. We assume that $`|V_0|`$ is sufficiently small that Born’s approximation holds. According to a previous formula, the scattering amplitude is given by
$`f^{(B)}(\theta ,\phi )=\frac{1}{4\pi }\frac{2mV_0}{\mathrm{}^2}e^{i𝐊𝐫}\frac{e^{\alpha r}}{r}d^3r.`$
Since this potential depends only on $`r`$, the angular integrals are trivial leading to the form
$`f^{(B)}(\theta ,\phi )=\frac{1}{4\pi }\frac{2mV_0}{\mathrm{}^2}\frac{4\pi }{|𝐊|}_0^{\mathrm{}}\mathrm{sin}|𝐊|r\frac{e^{\alpha r}}{r}rdr.`$
Thus, we obtain
$`f^{(B)}(\theta ,\phi )=\frac{2mV_0}{\mathrm{}^2}\frac{1}{\alpha ^2+|𝐊|^2}.`$
From the figure we can notice that $`|𝐊|=2k\mathrm{sin}\frac{\theta }{2}`$. Therefore
$`\sigma ^{(B)}(\theta )=\frac{4m^2V_0^2}{\mathrm{}^4}\frac{1}{[\alpha ^2+4k^2\mathrm{sin}\frac{\theta }{2}^2]^2}.`$
The total cross section is obtained by integrating
$`\sigma ^{(B)}=\sigma ^{(B)}(\theta )𝑑\mathrm{\Omega }=\frac{4m^2V_0^2}{\mathrm{}^4}\frac{4\pi }{\alpha ^2(\alpha ^2+4k^2)}.`$
8. PARTIAL WAVES
## Introduction
The partial waves method is quite general and applies to particles interacting in very small spatial regions with another one, which is usually known as scattering center because of its physical characteristics. (for example, because it can be considered as fixed). Beyond the interaction region, the interaction between the two particles is usually negligible. Under this circumstances, it is possible to describe the scattered particle by means of the Hamiltonian
$$H=H_0+V,$$
(1)
where $`H_0`$ corresponds to the free particle Hamiltonian. Our problem is to solve the equation
$$(H_0+V)\psi =E\psi .$$
(2)
Obviously, the spectrum will be continuous since we study the case of elastic scattering. The solution will be
$$\psi =\frac{1}{EH_0}V\psi +\varphi .$$
(3)
It is easy to see that for $`V=0`$ one can obtain the solution $`\varphi `$, i.e., the solution corresponding to the free particle. It is worth noting that in a certain sense the operator $`\frac{1}{EH_0}`$ is anomalous, because it has a continuum of poles on the real axis at positions coinciding with the eigenvalues of $`H_0`$. To get out of this trouble, it is common to produce a small shift in the imaginary direction ($`\pm iϵ`$) of the cut on the real axis
$$\psi ^\pm =\frac{1}{EH_0\pm i\epsilon }V\psi ^\pm +\varphi $$
(4)
This equation is known as the Lippmann-Schwinger equation. Finally, the shift of the poles is performed in the positive sense of the imaginary axis because in this case the causality principle holds (cf. Feynman). Let us consider the x representation
$$𝐱\psi ^\pm =𝐱\varphi +d^3x^{^{}}𝐱|\frac{1}{EH_0\pm i\epsilon }|𝐱^{^{}}𝐱^{^{}}V\psi ^\pm .$$
(5)
The first term on the right hand side corresponds to a free particle, while the second one is interpreted as a spherical wave getting out from the scattering center. The kernel of the previous integral can be considered as a Green function (also called propagator in quantum mechanics). It is a simple matter to calculate it
$$G_\pm (𝐱,𝐱^{^{}})=\frac{\mathrm{}^2}{2m}𝐱|\frac{1}{EH_0\pm i\epsilon }|𝐱^{^{}}=\frac{1}{4\pi }\frac{e^{\pm ik𝐱𝐱^{^{}}}}{𝐱𝐱^{^{}}},$$
(6)
where $`E=\mathrm{}^2k^2/2m`$. Writing the wave function as a plane wave plus a divergent spherical one (up to a constant factor),
$$𝐱\psi ^+=e^{𝐤𝐱}+\frac{e^{ikr}}{r}f(𝐤,𝐤^{^{}}).$$
(7)
the quantity $`f(𝐤,𝐤^{^{}})`$ is known as the scattering amplitude and is explicitly
$$f(𝐤,𝐤^{^{}})=\frac{1}{4\pi }(2\pi )^3\frac{2m}{\mathrm{}^2}𝐤^{^{}}V\psi ^+.$$
(8)
Let us now define an operator $`T`$ such that
$$T\varphi =V\psi ^+$$
(9)
If we multiply the Lippmann-Schwinger equation by $`V`$ and make use of the previous definition, we get
$$T\varphi =V\varphi +V\frac{1}{EH_0+i\epsilon }T\varphi .$$
(10)
Iterating this equation (as in perturbation theory) we can get the Born approximation and its higher-order corrections.
## Partial waves method
Let us now consider the case of a central potential. In this case, using the definition (9), it is found that the operator $`T`$ commutes with $`\stackrel{}{L}^2`$ and $`\stackrel{}{L}`$; it is said that $`T`$ is a scalar operator. To simplify the calculations it is convenient to use spherical coordinates, because of the symmetry of the problem that turns the $`T`$ operator diagonal. Let us see now a more explicit form of the scattering amplitude
$$f(𝐤,𝐤^{^{}})=\mathrm{const}.\underset{lml^{^{}}m^{^{}}}{}𝑑E𝑑E^{^{}}𝐤^{^{}}E^{^{}}l^{^{}}m^{^{}}E^{^{}}l^{^{}}m^{^{}}TElmElm𝐤,$$
(11)
where $`\mathrm{const}.=\frac{1}{4\pi }\frac{2m}{\mathrm{}^2}(2\pi )^3`$. After some calculation, one gets
$$f(𝐤,𝐤^{^{}})=\frac{4\pi ^2}{k}\underset{l}{}\underset{m}{}T_l(E)Y_l^m(𝐤^{^{}})Y_l^m^{}(𝐤).$$
(12)
Choosing the coordinate system such that the vector $`𝐤`$ have the same direction with the z axis, one infers that only the spherical harmonics of $`m=0`$ will contribute to the scattering amplitude. If we define by $`\theta `$ the angle between $`𝐤`$ and $`𝐤^{^{}}`$, we will get
$$Y_l^0(𝐤^{^{}})=\sqrt{\frac{2l+1}{4\pi }}P_l(cos\theta ).$$
(13)
Employing the following definition
$$f_l(k)\frac{\pi T_l(E)}{k},$$
(14)
eq. (12) can be written as follows
$$f(𝐤,𝐤^{^{}})=f(\theta )=\underset{l=0}{\overset{\mathrm{}}{}}(2l+1)f_l(k)P_l(cos\theta ).$$
(15)
For $`f_l(k)`$ a simple interpretation can be provided, which is based on the expansion of a plane wave in spherical waves. Thus, we can write the function $`𝐱\psi ^+`$ for large values of $`r`$ in the following form
$$𝐱\psi ^+=\frac{1}{(2\pi )^{3/2}}\left[e^{ikz}+f(\theta )\frac{e^{ikr}}{r}\right]=$$
$$\frac{1}{(2\pi )^{3/2}}\left[\underset{l}{}(2l+1)P_l(\mathrm{cos}\theta )\left(\frac{e^{ikr}e^{i(krl\pi )}}{2ikr}\right)+\underset{l}{}(2l+1)f_l(k)P_l(\mathrm{cos}\theta )\frac{e^{ikr}}{r}\right]$$
$$=\frac{1}{(2\pi )^{3/2}}\underset{l}{}(2l+1)\frac{P_l(\mathrm{cos}\theta )}{2ik}\left[\left[1+2ikf_l(k)\right]\frac{e^{ikr}}{r}\frac{e^{i(krl\pi )}}{r}\right].$$
(16)
This expression can be interpreted as follows. The two exponential terms correspond to spherical waves: the first to a divergent wave, and the latter to a convergent one. Moreover, the scattering effect is conveniently displayed in the coefficient of the divergent wave, which is unity when there are no scattering centers.
## Phase shifts
We consider now a surface enclosing the scattering center. Assuming that there is no creation and annihilation of particles, one has
$$𝐣𝑑𝐒=0,$$
(17)
where the integration region is the aforementioned surface, and $`𝐣`$ is the probability current density. Moreover, because of the conservation of the orbital momentum, the latter equation should hold for each partial wave. The theoretical formulation of the problem does not change if one assumes the wave packet as a flux of noninteracting particles propagating through a region of central potential for which the angular momentum of each particle is conserved, so that the ‘particle’ content of the wave packet really does not change. Thus, one may think even intuitevely that only phase factor effects can be introduced under these circumstances. Thus, if one defines
$$S_l(k)1+2ikf_l(k)$$
(18)
we should have
$$S_l(k)=1.$$
(19)
These results can be interpreted using the conservation of probabilities. They are natural and expected because we assumed that there is no creation and annihilation of particles. Therefore, the effects of the scattering center is reduced to adding a phase factor in the components of the divergent wave. Taking into account the unitarity of the phase factor, we can write it in the form
$$S_l=e^{2i\delta _l},$$
(20)
where $`\delta _l`$ is a real function of $`k`$. Taking into account the definition (18), we can write
$$f_l=\frac{e^{2i\delta _l}1}{2ik}=\frac{e^{i\delta _l}\mathrm{sin}(\delta _l)}{k}=\frac{1}{k\mathrm{cot}(\delta _l)ik}.$$
(21)
The total cross section has the following form
$$\sigma _{total}=f(\theta )^2𝑑\mathrm{\Omega }=$$
$$\frac{1}{k^2}_0^{2\pi }𝑑\varphi _1^1d(\mathrm{cos}(\theta ))\underset{l}{}\underset{l^{^{}}}{}(2l+1)(2l^{^{}}+1)e^{i\delta _l}\mathrm{sin}(\delta _l)e^{i\delta _l^{^{}}}\mathrm{sin}(\delta _l^{^{}})P_lP_l^{^{}}$$
$$=\frac{4\pi }{k^2}\underset{l}{}(2l+1)\mathrm{sin}{}_{}{}^{2}(\delta _l^{^{}}).$$
(22)
## Getting the phase shifts
Let us consider now a potential V that is zero for $`r>R`$, where the parameter $`R`$ is known as the range of the potential. Thus, the region $`r>R`$ corresponds to a spherical unperturbed/free wave. On the other hand, the general form of the expansion of a plane wave in spherical ones is
$$𝐱\psi ^+=\frac{1}{(2\pi )^{3/2}}\underset{l}{}i^l(2l+1)A_l(r)P_l(\mathrm{cos}\theta )(r>R),$$
(23)
where the coefficient $`A_l`$ is by definition
$$A_l=c_l^{(1)}h_l^{(1)}(kr)+c_l^{(2)}h_l^{(2)}(kr).$$
(24)
$`h_l^{(1)}`$ and $`h_l^{(2)}`$ are the spherical Hankel functions whose asymptotic forms are the following
$$h_l^{(1)}\frac{e^{i(krl\pi /2)}}{ikr}$$
$$h_l^{(2)}\frac{e^{i(krl\pi /2)}}{ikr}.$$
Inspecting the following asymptotic form of the expression (23)
$$\frac{1}{(2\pi )^{3/2}}\underset{l}{}(2l+1)P_l\left[\frac{e^{ikr}}{2ikr}\frac{e^{i(krl\pi )}}{2ikr}\right],$$
(25)
one can see that
$$c_l^{(1)}=\frac{1}{2}e^{2i\delta _l}c_l^{(2)}=\frac{1}{2}.$$
(26)
This allows to write the radial wave function for $`r>R`$ in the form
$$A_l=e^{2i\delta _l}\left[\mathrm{cos}\delta _lj_l(kr)\mathrm{sin}\delta _ln_l(kr)\right].$$
(27)
Using the latter equation, we can get the logarithmic derivative in $`r=R`$, i.e., at the boundary of the potential range
$$\beta _l\left(\frac{r}{A_l}\frac{dA_l}{dr}\right)_{r=R}=kR\left[\frac{j_l^{^{}}\mathrm{cos}\delta _ln_l^{^{}}(kR)\mathrm{sin}\delta _l}{j_l\mathrm{cos}\delta _ln_l(kR)\mathrm{sin}\delta _l}\right].$$
(28)
$`j_l^{^{}}`$ is the derivative of $`j_l`$ with respect to $`r`$ evaluated at $`r=R`$. Another important result that can be obtained from the knowledge of the previous one is the phase shift
$$\mathrm{tan}\delta _l=\frac{kRj_l^{^{}}(kR)\beta _lj_l(kR)}{kRn_l^{^{}}(kR)\beta _ln_l(kR)}.$$
(29)
To get the complete solution of the problem in this case, it is necessary to make the calculations for $`r<R`$, i.e., within the range of the potential. For a central potential, the 3D Schrödinger equation reads
$$\frac{d^2u_l}{dr^2}+\left(k^2\frac{2m}{\mathrm{}^2}V\frac{l(l+1)}{r^2}\right)u_l=0,$$
(30)
where $`u_l=rA_l(r)`$ is constrained by the boundary condition $`u_l_{r=0}=0`$. Thus, one can calculate the logarithmic derivative, which, taking into account the continuity of the log-derivative (equivalent to the continuity condition of the derivative at a discontinuity point) leads to
$$\beta _l_{in}=\beta _l_{out}.$$
(31)
## An example: scattering on a hard sphere
Let us now consider an important illustrative case, that of the hard sphere potential
$$V=\{\begin{array}{cc}\mathrm{}\hfill & \text{ }r<R\hfill \\ 0\hfill & r>R.\hfill \end{array}$$
(32)
It is known that a particle cannot penetrate into a region where the potential is infinite. Therefore, the wave function should be zero at $`r=R`$. Since we deal with an impenetrable sphere we also have
$$A_l(r)_{r=R}=0.$$
(33)
Thus, from eq. (27), we get
$$\mathrm{tan}\delta _l=\frac{j_l(kR)}{n_l(kR)}.$$
(34)
One can see that the phase shift calculation is an easy one for any $`l`$. In the $`l=0`$ case (s wave scattering), we have
$$\delta _l=kR$$
and from eq. (27)
$$A_{l=0}(r)\frac{\mathrm{sin}kr}{kr}\mathrm{cos}\delta _0+\frac{\mathrm{cos}kr}{kr}\mathrm{sin}\delta _0=\frac{1}{kr}\mathrm{sin}(kr+\delta _0).$$
(35)
We immediately see that there is an additional phase contribution with regard to the motion of the free particle. It is also clear that in more general cases the various waves will have different phase shifts leading to a transient distortion of the scattered wave packet. At small energies, i.e., $`kR<<1`$, the spherical Bessel functions (entering the formulas for the spherical Hankel functions) are the following
$$j_l(kr)\frac{(kr)^l}{(2l+1)!!}$$
(36)
$$n_l(kr)\frac{(2l1)!!}{(kr)^{l+1}},$$
(37)
leading to
$$\mathrm{tan}\delta _l=\frac{(kR)^{2l+1}}{(2l+1)[(2l1)!!]^2}.$$
(38)
From this formula, one can see that a substantial contribution to the phase shift is given by the $`l=0`$ waves. Moreover, since $`\delta _0=kR`$ the cross section is obtained as follows
$$\sigma _{total}=\frac{d\sigma }{d\mathrm{\Omega }}𝑑\mathrm{\Omega }=4\pi R^2.$$
(39)
One can see that the total scattering cross section is four times bigger than the classical one and coincides with the total area of the impenetrable sphere. For large values of the incident energy, one can work in the hypothesis that all values of $`l`$ up to a maximum value $`l_{max}kR`$ contribute to the total cross section
$$\sigma _{total}=\frac{4\pi }{k^2}\underset{l=0}{\overset{lkR}{}}(2l+1)\mathrm{sin}^2\delta _l.$$
(40)
In this way, from eq. (34), we have
$$\mathrm{sin}^2\delta _l=\frac{\mathrm{tan}^2\delta _l}{1+\mathrm{tan}^2\delta _l}=\frac{[j_l(kR)]^2}{[j_l(kR)]^2+[n_l(kR)]^2}\mathrm{sin}^2\left(kR\frac{l\pi }{2}\right),$$
(41)
where the expressions
$$j_l(kr)\frac{1}{kr}\mathrm{sin}\left(kr\frac{l\pi }{2}\right)$$
$$n_l(kr)\frac{1}{kr}\mathrm{cos}\left(kr\frac{l\pi }{2}\right).$$
have been used.
Inspection of $`\delta _l`$ shows a negative jump of $`\frac{\pi }{2}`$ whenever $`l`$ is augmented by a unity. Thus, it is clear that $`\mathrm{sin}^2\delta _l+\mathrm{sin}^2\delta _{l+1}=1`$ holds. Approximating $`\mathrm{sin}^2\delta _l`$ by its mean value $`\frac{1}{2}`$ over a period and using the sum of odd numbers, one gets
$$\sigma _{total}=\frac{4\pi }{k^2}(kR)^2\frac{1}{2}=2\pi R^2.$$
(42)
Once again the quantum-mechanical result, although quite similar to the corresponding classical result is nevertheless different. What might be the origin of the factor of two that makes the difference ? To get an explanation, we first separate eq. (15) in two parts
$$f(\theta )=\frac{1}{2ik}\underset{l=0}{\overset{l=kR}{}}(2l+1)e^{2i\delta _l}P_l\mathrm{cos}(\theta )+\frac{i}{2k}\underset{l=0}{\overset{l=kR}{}}(2l+1)P_l\mathrm{cos}(\theta )=f_{\text{refl}}+f_{\text{shadow}}.$$
(43)
Calculation of $`|f_{\text{ refl}}|^2𝑑\mathrm{\Omega }`$ gives
$$|f_{\text{ refl}}|^2𝑑\mathrm{\Omega }=\frac{2\pi }{4k^2}\underset{l=0}{\overset{l_{max}}{}}_{1}^{}{}_{}{}^{1}(2l+1)^2[P_l\mathrm{cos}(\theta )]^2d(\mathrm{cos}\theta )=\frac{\pi l_{max}^2}{k^2}=\pi R^2.$$
(44)
Analysing now $`f_{\text{shadow}}`$ at small angles, we get
$$f_{\text{shadow}}\frac{i}{2k}(2l+1)J_0(l\theta )ik_0^RbJ_0(kb\theta )𝑑b=\frac{iRJ_1(kR\theta )}{\theta }.$$
(45)
This formula is rather well known in optics. It corresponds to the Fraunhofer diffraction. Employing the change of variable $`z=kR\theta `$ one can calculate the integral $`|f_{\text{ shadow }}|^2𝑑\mathrm{\Omega }`$
$$|f_{\text{shadow}}|^2𝑑\mathrm{\Omega }2\pi R^2_0^{\mathrm{}}\frac{[J_1(z)]^2}{z}𝑑z\pi R^2.$$
(46)
Finally, neglecting the interference between $`f_{\text{refl}}`$ and $`f_{\text{ shadow }}`$ (since the phase oscillates between $`2\delta _{l+1}=2\delta _l\pi `$), one gets the result (42). The label ‘shadow’ for one of the terms is easily explained if one thinks of the wavy behaviour of the scattered particle (from the physical viewpoint there is no difference between a wave packet and a particle in this case). Its origin can be traced back to the backward-scattered components of the wave packet leading to a phase shift with respect to the incident waves and destructive interference.
## Coulomb scattering
In this section we briefly consider the Coulomb scattering in the quantum-mechanical approach. For this case, the Schrödinger equation is
$$\left(\frac{\mathrm{}^2}{2m}^2\frac{Z_1Z_2e^2}{r}\right)\psi (𝐫)=E\psi (𝐫),E>0,$$
(47)
where $`m`$ is the reduced mass of the system, $`E>0`$ since we deal with the simple scattering case where no kind of bound states are allowed to form. The previous equation is equivalent to the following expression (for adequate values of the constants $`k`$ and $`\gamma `$)
$$\left(^2+k^2+\frac{2\gamma k}{r}\right)\psi (𝐫)=0.$$
(48)
If we do not consider the centrifugal barrier, i.e., we look only to the $`s`$ waves, we really deal with a pure coulombian interaction, for which one can propose a solution of the following form
$$\psi (𝐫)=e^{i𝐤𝐫}\chi (u),$$
(49)
where
$$u=ikr(1\mathrm{cos}\theta )=ik(rz)=ikw,$$
$$𝐤𝐫=kz.$$
$`\psi (𝐫)`$ is the complete solution of the Schrödinger equation with an asymptotic ‘physical’ behaviour to which a plane wave $`e^{i𝐤𝐫}`$ and a spherical wave are expected to contribute $`r^1e^{ikr}`$ are expected to contribute. Defining new variables
$$z=zw=rz\lambda =\varphi ,$$
and by employing of previous relationships, eq. (48) takes the form
$$\left[u\frac{d^2}{du^2}+(1u)\frac{d}{du}i\gamma \right]\chi (u)=0.$$
(50)
To solve this equation, one should first study its asymptotic behaviour. Since we have already tackled this issue, we merely present the asymptotic normalized wave function that is the final result of all previous calculations
$$\psi _𝐤(𝐫)=\frac{1}{(2\pi )^{3/2}}\left(e^{i[𝐤𝐫\gamma ln(kr𝐤𝐫)]}+\frac{f_c(k,\theta )e^{i[kr+\gamma ln2kr]}}{r}\right).$$
(51)
As one can see, this wave function displays terms that turns it quite different from the form in eq. (7). This is due to the fact that the Coulomb potential is of infinite range. Performing the exact calculation for the Coulomb scattering amplitude is not an easy matter. Here we give only the final result for the ‘normalized’ wave function
$$\psi _𝐤(𝐫)=\frac{1}{(2\pi )^{3/2}}\left(e^{i[𝐤𝐫\gamma ln(kr𝐤𝐫)]}+\frac{g_1^{}(\gamma )}{g_1(\gamma )}\frac{\gamma }{2k\mathrm{sin}(\theta /2)^2}\frac{e^{i[kr+\gamma ln2kr]}}{r}\right),$$
(52)
where $`g_1(\gamma )=\frac{1}{\mathrm{\Gamma }(1i\gamma )}`$.
In addition, we reduce the partial wave analysis to a clear cut presentation of the results, of which some have already been mentioned. First of all, we write the wave function $`\psi (𝐫)`$ in (49) as follows
$$\psi (𝐫)=e^{i𝐤𝐫}\chi (u)=Ae^{i𝐤𝐫}_Ce^{ut}t^{i\gamma 1}(1t)^{i\gamma }𝑑t,$$
(53)
where $`A`$ is a ‘normalization’ constant, while all the integral part is the inverse Laplace transform of the direct transform of eq. (50). A convenient form of the latter equation is
$$\psi (𝐫)=A_Ce^{i𝐤𝐫}(1t)e^{ikrt}(1t)d(t,\gamma )𝑑t$$
(54)
where
$$d(t,\gamma )=t^{i\gamma 1}(1t)^{i\gamma 1}.$$
(55)
Within the partial wave analysis we proceed by writing
$$\psi (𝐫)=\underset{l=0}{\overset{\mathrm{}}{}}(2l+1)i^lP_l(\mathrm{cos}\theta )A_l(kr),$$
(56)
where
$$A_l(kr)=A_Ce^{ikrt}j_l[kr(1t)](1t)d(t,\gamma ).$$
(57)
Applying the relationships between the spherical Bessel functions and the Hankel functions, we get
$$A_l(kr)=A_l^{(1)}(kr)+A_l^{(2)}(kr).$$
(58)
We shall not sketch here how these coefficients are obtained (this is quite messy). They are
$$A_l^{(1)}(kr)=0$$
(59)
$$A_l^{(2)}(kr)\frac{Ae^{\pi \gamma /2}}{2ikr}[2\pi ig_1(\gamma )]\left(e^{i[kr(l\pi /2)+\gamma \mathrm{ln}2kr]}e^{2i\eta _l(k)}e^{i[kr(l\pi /2)+\gamma \mathrm{ln}2kr]}\right)$$
(60)
where
$$e^{2i\eta _l(k)}=\frac{\mathrm{\Gamma }(1+li\gamma )}{\mathrm{\Gamma }(1+l+i\gamma )}.$$
(61)
## Calculation of the Coulomb scattering amplitude
If we perform the Laplace transform of eq. (50), we get
$$\chi (u)=A_Ce^{ut}t^{i\gamma 1}(1t)^{i\gamma }𝑑t.$$
(62)
The contour $`C`$ goes from $`\mathrm{}`$ to $`\mathrm{}`$ on the real axis and closes through the upper half-plane. There are two poles in this case at $`t=0`$ and $`t=1`$. By the change of variable $`s=ut`$, we get
$$\chi (u)=A_{C_1}e^ss^{i\gamma 1}(us)^{i\gamma }.$$
(63)
$`\chi (u)`$ should be regular in zero. Indeed, we get
$$\chi (0)=(1)^{i\gamma }A_{C_1}\frac{e^s}{s}ds.=(1)^{i\gamma }A2\pi i$$
(64)
Performing now the limit $`u\mathrm{}`$, let’s do an infinitesimal shift to avoid the location of the poles on the contour. Moreover, by the change of variable $`\frac{s}{u}=\frac{(s_0\pm i\epsilon )}{i\kappa }`$, we see that this expression goes to zero when $`u\mathrm{}`$. Thus, we can expand $`(us)`$ in power series of $`\frac{s}{u}`$ for the pole with $`s=0`$. This expansion is not the right one in $`s=1`$, because in this case $`s=s_0+i(\kappa \pm \epsilon )`$. It comes out that $`\frac{s}{u}=1\frac{(s_0\pm i\epsilon )}{\kappa }`$ tends to $`1`$ when $`\kappa \mathrm{}`$. If instead we do the change of variable $`s^{^{}}=su`$, we get rid of this difficulty
$$\chi (u)=A_{\mathrm{C}_2}\left([e^ss^{i\gamma 1}(us)^{i\gamma }]ds+[e^{s^{^{}}+u}(s^{^{}})^{i\gamma }(u+s^{^{}})^{i\gamma 1}]ds^{^{}}\right).$$
(65)
Expanding the power series, it is easy to calculate the previous integrals, but one should take the limit $`\frac{s}{u}0`$ in the result in order to get the correct asymptotic forms for the Coulomb scattering
$$\chi (u)2\pi iA\left[u^{i\gamma }g_1(\gamma )(u)^{i\gamma 1}e^ug_2(\gamma )\right]$$
$$2\pi g_1(\gamma )=i_{\mathrm{C}_2}e^ss^{i\gamma 1}𝑑s$$
$$2\pi g_2(\gamma )=i_{\mathrm{C}_2}e^ss^{i\gamma }𝑑s.$$
(66)
After all this chain of variable changes, we get back to the original $`s`$ one to obtain
$$(u^{})^{i\gamma }=(i)^{i\gamma }[k(rz)]^{i\gamma }=e^{\gamma \pi /2}e^{i\gamma \mathrm{ln}k(rz)}$$
$$(u)^{i\gamma }=(i)^{i\gamma }[k(rz)]^{i\gamma }=e^{\gamma \pi /2}e^{i\gamma \mathrm{ln}k(rz)}.$$
(67)
The calculation of $`\chi `$, once effected, is equivalent with having $`\psi _𝐤(𝐫)`$ starting from (49).
## Eikonal approximation
We shall briefly expound on the eikonal approximation whose philosophy is the same to that used when one wants to pass from the wave optics to the geometrical optics. Therefore, it is the right approximation when the potential varies slowly over distances comparable to to the wavelength of the scattered wave packet, i.e., for the case $`E>>|V|`$. Thus, this approximation may be considered as a quasiclassical one. First, we propose that the quasiclassical wave function has the known form
$$\psi e^{iS(𝐫)/\mathrm{}},$$
(68)
where $`S`$ satisfies the Hamilton-Jacobi equation, having the solution
$$\frac{S}{\mathrm{}}=_{\mathrm{}}^z\left[k^2\frac{2m}{\mathrm{}^2}V\left(\sqrt{b^2+z^2}\right)\right]^{1/2}𝑑z^{}+\text{ constant}.$$
(69)
The additive constant is chosen in such a way to fulfill
$$\frac{S}{\mathrm{}}kz\mathrm{for}V0.$$
(70)
The term multiplying the potential can be interpreted as a change of phase of of the wave packet, having the following explicit form
$$\mathrm{\Delta }(b)\frac{m}{2k\mathrm{}^2}_{\mathrm{}}^{\mathrm{}}V\left(\sqrt{b^2+z^2}\right)𝑑z.$$
(71)
Within the method of partial waves, the eikonal approximation has the following application. We know it is correct at high energies, where many partial waves do contribute to the scattering. Thus, we can consider $`l`$ as a continuous variable and by analogy to classical mechanics we let $`l=bk`$. Moreover, as we already mentioned $`l_{max}=kR`$, which plugged into eq. (15) leads to
$$f(\theta )=ikbJ_0(kb\theta )[e^{2i\mathrm{\Delta }(b)}1]𝑑b.$$
(72)
## 8P. Problems
Problem 8.1
Obtain the phase shift and the differential cross section at small angles for a scattering centre of potential $`U(r)=\frac{\alpha }{r^2}`$. It should be taken into account that for low-angle scattering the main contribution is given by the partial waves of large $`l`$.
Solution:
Solving the equation
$$R_l^{^{\prime \prime }}+\left[k^2\frac{l(l+1)}{r^2}\frac{2m\alpha }{\mathrm{}^2r^2}\right]=0$$
with the boundary conditions $`R_l(0)=0`$, $`R_l(\mathrm{})=N`$, where $`N`$ is a finite number, we get
$$R_l(r)=A\sqrt{r}I_\lambda (kr),$$
where $`\lambda =\left[(l+\frac{1}{2})^2+\frac{2m\alpha }{\mathrm{}^2}\right]^{1/2}`$ and $`I`$ is the first modified Bessel function.
To determine $`\delta _l`$, one should use the asymptotic expression of $`I_\lambda `$:
$$I_\lambda (kr)\left(\frac{2}{\pi kr}\right)^{1/2}\mathrm{sin}(kr\frac{\lambda \pi }{2}+\frac{\pi }{4}).$$
Therefore
$$\delta _l=\frac{\pi }{2}\left(\lambda l\frac{1}{2}\right)=\frac{\pi }{2}\left(\left[(l+\frac{1}{2})^2+\frac{2m\alpha }{\mathrm{}^2}\right]^{1/2}\left(l+\frac{1}{2}\right)\right).$$
The condition of large $`l`$ leads us to
$$\delta _l=\frac{\pi m\alpha }{(2l+1)\mathrm{}^2},$$
whence one can see that $`|\delta _l|1`$ for large $`l`$.
From the general expression of the scattering amplitude
$$f(\theta )=\frac{1}{2ik}\underset{l=0}{\overset{\mathrm{}}{}}(2l+1)P_l(\mathrm{cos}\theta )(e^{2i\delta _l}1),$$
at small angles one gets $`e^{2i\delta _l}1+2i\delta _l`$, so that
$$\underset{l=0}{\overset{\mathrm{}}{}}P_l(\mathrm{cos}\theta )=\frac{1}{2\mathrm{sin}\frac{\theta }{2}}.$$
Thus
$$f(\theta )=\frac{\pi \alpha m}{k\mathrm{}^2}\frac{1}{2\mathrm{sin}\frac{\theta }{2}}.$$
The final result is
$$\frac{d\sigma }{d\theta }=\frac{\pi ^3\alpha ^2m}{2\mathrm{}^2E}\text{c}tg\frac{\theta }{2}.$$
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# Transport and Elastic Properties of Fractal Media
## 1 Introduction
In the past decade a number of important natural and manufactured materials have been shown to exhibit microstructures with fractal characteristics. For example, it has been shown that sedimentary rocks have internal surfaces that exhibit fractal roughness and in some cases the fractal regime may include the entire pore volume . This geometric complexity has a number of important implications for transport properties in porous rocks. In another example it was shown that fractured surfaces of metals are fractal in nature and proposed that the fractal dimensions of the surfaces could be correlated to the impact energy of the sample . Finally low-weight polymeric gel materials (aerogels) give a unique combination of properties; low-density and insulative, yet mechanically strong . Aerogels are good examples of fractal materials, made up of clusters generated from the aggregation of primary particles .
While the concept of a fractal geometry provides a description of disordered material structure, it is not clear if the properties of these materials are affected by fractal characteristics nor how one may design and optimise fractal structure for a given application. Microstructural characteristics other than fractal dimension (e.g., connectivity, coordination) may be more important when deriving Structure-Property relationships. For example, in aerogels, the macroscopic connectivity at low density is the key microstructural characteristic associated with the mechanical strength of the material.
To date there has been no explicit investigation of the influence of fractal structure on material properties. Theoretical studies have relied on simplistic models of fractal microstructure. In general the relationship of such models to real materials is unclear, and the predictive ability of the models is qualitative in nature. In this paper we investigate the influence of fractal structure on material properties. We find that a fractally rough interface has a strong influence on material properties. A fractal volume, in contrast, has little effect on properties.
## 2 Evaluation of Bounds
For general random composites there is no exact method of predicting properties. However a number of rigorous bounds, which depend on the microstructure of the composites, have been derived. These include bounds on the conductivity and the bulk and shear moduli of composite materials (reviewed in Ref. ). The bounds are expressed in terms of the volume fractions and properties of each of the phases and two microstructure parameters:
$`\zeta _1`$ $`=`$ $`{\displaystyle \frac{9}{2pq}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dr}{r}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}{\displaystyle _1^1}𝑑uP_2(u)\left(p_3(r,s,t){\displaystyle \frac{p_2(r)p_2(s)}{p}}\right)`$ (1)
$`\eta _1`$ $`=`$ $`{\displaystyle \frac{5\zeta _1}{21}}+{\displaystyle \frac{150}{7pq}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dr}{r}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}{\displaystyle _1^1}𝑑uP_4(u)\left(p_3(r,s,t){\displaystyle \frac{p_2(r)p_2(s)}{p}}\right)`$ (2)
where $`p_3(r,s,t)`$ is the three point correlation function , $`t^2=r^2+s^22rsu`$ and $`P_n(u)`$ denotes the Legendre polynomial of order $`n`$. Until recently the use of the bounds has been restricted to ‘particulate’ media, such as uncorrelated overlapping spheres . Such models have limited utility in describing the microstructure of composite materials. We have shown that microstructure generated from level cuts of a random standing wave mimics the microstructure of a wide range of real composite materials including polymer blends , porous rocks and foamed solids . We have subsequently evaluated the bounds for a variety of level-cut microstructures and shown that it is possible to correlate quantitatively the effective physical properties of materials to their microstructure . Berk has shown how the level-cut scheme may be used to model fractal interfaces, and the extension to volume fractals is simple. Here we investigate the properties of these surface- and volume-fractal materials using rigorous bounds. This allows us to investigate the influence of fractal structure on material properties.
A key quantity in the characterization of two phase materials is the two-point correlation function, $`p_2(r)=p(1p)\gamma (r)+p^2`$, where $`p`$ is the volume fraction of the first phase and $`\gamma (r)`$ is the normalized correlation function with $`\gamma (0)=1`$, $`\gamma (\mathrm{})=0`$. Physically $`p_2`$ gives the probability that two points a distance $`r`$ apart will lie in the first phase. Two principle fractal types are evident in composites: surface and volume (or mass) fractals. Each can be characterized by the behaviour of the normalized correlation function $`\gamma (r)`$. Surface fractal behaviour occurs when the interface between the phases is rough at all scales. In this case Bale and Schmidt have shown that $`1\gamma (r)r^{3D_s}`$ as $`r0`$ where $`2D_s<3`$ is the surface-fractal dimension. Experimental techniques have determined that a variety of materials exhibit fractal surfaces with $`D_s>2`$. These include crushed glass, zeolites , silica gels , lignite coals , porous rocks , soils and composite steels (see Korvin for a review). On the other hand, volume (or mass) fractals occur when voids (or inclusions) are present at all scales. In this case $`\gamma (r)r^{(3D_v)}`$ as $`r\mathrm{}`$ where $`D_v<3`$ is the volume-fractal dimension . Fractal behavior of this type has been observed in silica gels and soils .
## 3 Model morphologies: The level cut Gaussian random field
Random composites with a wide variety of morphological properties can be generated by taking level-cuts of a Gaussian random field (GRF). A simple definition of a GRF is
$$y(𝒓)=\sqrt{\frac{2}{N}}\underset{i=1}{\overset{N}{}}\mathrm{cos}(k_i\widehat{𝒌}𝒓+\varphi _i)$$
(3)
where $`\varphi _i`$ is a uniform deviate on $`[0,2\pi )`$ and $`\widehat{𝒌}_i`$ is uniformly distributed on a unit sphere. The magnitude of the wave vectors $`k_i`$ are distributed on $`[0,\mathrm{})`$ with a probability (spectral) density $`P(k)`$ ($`P(k)𝑑k=1`$). A composite material can then be defined by taking the region in space where $`\alpha y(𝒓)\beta `$ as the first phase, while the two regions contiguous to this $`y(𝒓)<\alpha ;y(𝒓)>\beta `$ define a complementary second phase. In the case $`\beta =\mathrm{}`$ a 1-cut material results; in the case $`\beta =\alpha `$ a 2-cut material.
The statistical correlation functions of these materials can be calculated . The volume fraction is given by $`p=(2\pi )^{\frac{1}{2}}_\alpha ^\beta e^{\frac{1}{2}t^2}𝑑t`$ and the two point correlation function is
$`p_2(r)=`$ $`p^2+{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{g(r)}}{\displaystyle \frac{dt}{\sqrt{1t^2}}}\times [\mathrm{exp}({\displaystyle \frac{\alpha ^2}{1+t}})`$ (4)
$`\mathrm{exp}({\displaystyle \frac{1}{2}}{\displaystyle \frac{\alpha ^22\alpha \beta t+\beta ^2}{(1t^2)}})+\mathrm{exp}({\displaystyle \frac{\beta ^2}{1+t}})].`$
Here $`g(r)=y(0)y(𝒓)`$ is the field-field correlation function and is related to the spectral density of the field
$$g(r)=_0^{\mathrm{}}P(k)\frac{\mathrm{sin}kr}{kr}𝑑k.$$
(5)
In terms of the normalized 2-point function $`\gamma (r)=(p_2(r)p^2)/(pp^2)`$ it can be shown that
$$1\gamma (r)(1g(r))^{\frac{1}{2}}r0\&\gamma (r)g(r)r\mathrm{}.$$
(6)
Higher order correlation functions can also be evaluated for level-cut GRF media . The freedom in specifying $`\alpha `$, $`\beta `$ and $`P(k)`$ (or $`g(r)`$) allows a wide variety of materials to be modelled.
## 4 Surface fractals
Recent microstructural studies of composite steels have revealed a fractal boundary between a ferrite and pearlite phase. Using a scaling relationship between the area and perimeter of the pearlite grains it was found that $`d_f1.5`$. In three-dimensions this implies a surface fractal dimension of $`D_s=1+d_f2.5`$. Small angle scattering studies have shown that lignite coal has $`D_s2.5`$ and sedimentary sandstones have $`2.25D_s2.96`$ with many exhibiting $`D_s=2.5\pm 0.1`$ . Materials with such a fractal surface can be generated in the GRF scheme by using, for example, $`g(r)=e^{r/l}`$ where $`l`$ is a pore/grain length scale which we normalize to unity. In this case $`1\gamma (r)r^{\frac{1}{2}}`$ as $`r0`$ so that $`D_s=2.5`$. The corresponding spectral density is
$$P(k)=\frac{4k^2}{\pi (1+k^2)^2}.$$
(7)
To account for the fact that physical surfaces only exhibit fractal scaling down to some finite length scale (e.g. the lattice constant) a cut-off length can be introduced in the model by setting $`P(k)=0`$ for $`k>K`$ (and scaling so $`P(k)𝑑k=1`$). This leaves the large scale microstructure unchanged and filters out the high frequency (rougher) sinusoids in Eqn. (3) above a length scale $`2\pi /K`$. The magnitude of $`K`$ is estimated from experimental studies. In ferritic-pearlite steels the minimum scale reported is $`.1\mu m`$ and the grain size is $`O(50\mu m)`$ so $`K2\pi 50/.1=3.1\times 10^3`$. In sandstones the minimum scale measured is $`.5\times 10^3\mu m`$ and the pore-size is around $`10\mu m`$ so $`K2\pi 10/.0005=1.3\times 10^6`$. We assume $`K=\mathrm{}`$ provides a reasonable model of fractal steel and sandstone interfaces. To clearly demonstrate the effect of a fractal surface on material properties we also consider materials with $`K=8`$ (providing a smooth interface with $`D_s=2`$).
Cross sections of the surface fractal material are shown in Figs. 1 (a) & (b). The morphology of the 2-cut field is very similar to that observed in steel composites . The roughness of the interface for the 1-cut case is clearly visible in Fig. 2. The spectral density and associated 2-point correlation function of each of the materials are shown in Fig. 3. Ten two-dimensional realizations of each model are calculated from Eqn. (3) for the case $`\alpha =0`$ & $`\beta =\mathrm{}`$ ($`p=\frac{1}{2}`$). The average value of $`\gamma (r)`$ is represented by symbols in this figure: the fractal scaling of $`\gamma (r)`$ is evident.
To determine the effect of the fractal microstructure on the conductivity and elastic moduli of each material we evaluate the microstructure parameters for both models ($`K=\mathrm{},8`$). The results are given in Tables 1 & 2 and bounds on the shear moduli are shown in Fig. 4 for the case where the the shear and bulk moduli are both equal to 1 in phase 1 and 0 in phase 2. Note that the lower bound vanishes at this contrast. Curves for the bulk moduli and conductivity are qualitatively similar. The shear moduli of the fractal and smooth 1-cut models are quite similar. In contrast the fractal 2-cut model has a significantly reduced resistance to shear when compared to the smooth 2-cut model: a rough interface decreases the elastic strength of a material if the inclusion phase is stronger than the bulk phase. This result is similar to a trend observed in composite steels: an increase in $`D_s`$ leads to a decrease in impact toughness.
## 5 Volume fractals
Aerogels provide good examples of volume fractal materials with $`D_v2`$ . The formation process has been modelled by the cluster-cluster aggregation model which gives $`D_v1.71.8`$. To see how the level-cut scheme can model volume-fractal materials consider $`\gamma (r)C(r/l)^A`$ for $`r\mathrm{}`$. The volume of solid within a radius $`R`$ is
$`V(R)`$ $`=`$ $`{\displaystyle _0^R}𝑑vp_2(r)/p={\displaystyle _0^R}𝑑r4\pi r^2\left((1p)\gamma (r)+p\right)`$ (8)
$``$ $`{\displaystyle \frac{(1p)4\pi Cl^A}{3A}}R^{3A}+p{\displaystyle \frac{4\pi }{3}}R^3.`$
Now the the former term dominates for $`R<R_s=\left(\frac{3(1p)C}{p(3A)}\right)^{1/A}lp^{1/A}l`$ where $`R_s`$ is a saturation scale and $`V(R)R^{3A}=R^{D_v}`$. Thus the volume fractal dimension is $`D_v=3A`$. A simple field-field correlation function which gives rise to this behaviour in the level-cut GRF is
$$g(r)=\left(1+\frac{2}{A}\frac{r^2}{l^2}\right)^{A/2}.$$
(9)
This field-field correlation function actually leads to a spectral density
$$P(k)=\frac{1}{(2\pi )^3}_0^{\mathrm{}}4\pi r^2g(r)\frac{\mathrm{sin}kr}{kr}𝑑r$$
(10)
which is not strictly positive. Therefore we define a new $`P(k)`$ which is set to zero for all $`k`$ beyond the first point $`K`$ at which $`P(k)<0`$ (and re-scale to ensure $`P(k)𝑑k=1`$). This modification does not change the behaviour of the model at large $`r`$. In aerogels $`l`$ (which we normalize to unity) is related to the length scale of the monomers. In experiments the saturation scale is $`O(10`$nm) and in a recent model of a silica gel based on the cluster-cluster aggregation algorithm the saturation scale is around six times the particle diameter at concentrations of $`c=0.05`$ . This is consistent with Eqn. (8): $`R_sp^{1/A}=0.05^{1/1.7}=5.8`$. Although the level-cut materials generated using Eqn. (9) has fractal scaling over the correct range of $`r`$, it does not necessarily follow that this is an appropriate model for aerogels. It does, however, allow the dependence of properties on $`D_v`$ to be gauged within the level-cut scheme.
We choose a fractal $`A=1.5`$ ($`D_v=1.5`$) and non-fractal case $`A=3`$ (so the second term in Eqn. (8) dominates and $`D_v=3`$) to determine the effect of volume fractal behaviour on material properties. The functions $`P(k)`$ are shown in Fig. 5. We have generated realizations of the materials using Eqn. (3). Cross-sections of the model for the case $`D_v=1.5`$ are shown for the 1-cut and 2-cut cases in Figs. 1 (c) & (d). A three dimensional representation of the 2-cut interface is shown in Fig. 2. The 2-point function of each model has been measured for the case $`p=\frac{1}{2}`$ and averaged over 50 cross-section realizations of each model. The results, plotted in Fig. 5, show that the level-cut GRF model has volume-fractal scaling.
The microstructure parameters for both $`D_v=1.5`$ and $`D_v=3.0`$ are given in Tables 1 & 2 and bounds on the conductivity are shown in Fig. 6 for the case $`\sigma _{1,2}=1,0`$ (the lower bounds vanish at this contrast). As we would expect the sheet-like nature of the structures in the 2-cut media are much better conductors than the node/bond-like structures present in the 1-cut media. We find virtually no difference in the fractal and non-fractal materials in both the 1-cut and 2-cut cases. This surprising result suggests that the properties of volume-fractal composites (such as aerogels) are not explicitly dependent on the fractal dimension.
## 6 Conclusion
We investigate the influence of fractal structure on material properties. We have calculated rigorous bounds on the conductivity and elasticity of fractal media generated using the level-cut random field model. The behaviour of the bounds indicates that a fractal interface plays a minimal role in the properties of 1-cut media. For the two-cut model, which mimics the microstructure of both foams and porous rocks a much stronger influence is observed. In contrast, varying the volume-fractal dimension of both 1-cut and 2-cut media has little effect on the property bounds. The latter result indicates that the remarkable properties of aerogels are influenced more by the fact that they contain very well connected structures at high porosities, rather than their fractal characteristics . In the future we shall utilise a range of microstructural studies to develop a more appropriate model of aerogel structure. This will allow a more rigorous comparison between model and experimental properties.
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# Magnetization and Level Statistics at Quantum Hall Liquid-Insulator Transition in the Lattice Model
## Abstract
Statistics of level spacing and magnetization are studied for the phase diagram of the integer quantum Hall effect in a 2D finite lattice model with Anderson disorder.
The way in which the increasing disorder induces the insulating state when starting from the integer quantum Hall (IQH) state is a topic of controversy between the continuum and lattice models of the 2D electronic gas in strong magnetic field. The continuum approach predicts the crossover between the adjacent quantum Hall plateaus, ending up with the insulating state when the degree of disorder increases or, equivalently, the limit of small magnetic field is considered. This is due to the so-called ’floating up’ of the critical energies $`E_c`$ which occurs with increasing disorder. ($`E_c`$ is the energy where the localization-delocalization transition takes place in the thermodynamic limit). In the critical region, i.e., when the Fermi energy $`E_f`$ crosses an extended state energy $`E_c`$, the transverse (Hall) conductivity is $`\sigma _{xy}^c=\nu 1/2`$ (at the transition between the plateaus $`\nu `$ and $`\nu 1`$). This means that at large disorder (or low field), the cascade of transitions must end with $`\nu =1\nu =0`$ (insulator). The experiments give controversial information in what concerns the possibility to observe this last transition (see Ref.2 and the references therein, Ref.3) The sensible conclusion can be found in Ref.2 and Ref.4 suggesting that the theoretical results of the scaling theory, which are obtained for zero temperature and infinite systems, cannot be checked easily by experiments which are done for finite samples and at low (but nevertheless finite) temperature. The evaluation of the critical value of the longitudinal conductivity is also a difficult task. Lee, Kivelson and Zhang show in the frame of corresponding states law that $`\sigma _{xx}^c=1/2`$ for any $`\nu `$; approaching the question in the opposite way, Zirnbauer assumes $`\sigma _{xx}^c=1/2`$ and finds agreement with the numerical simulations. The numerical calculation performed by Huo, Hetzel and Bhatt for the lowest Landau level produces also $`0.5`$ ; Huckestein and Backhaus show that this value is correct even in the presence of the electron-electron interaction .
More recently, the same problem has been approached also in lattice models. The results are again controversial, since Yang and Bhatt , in a discussion based on the calculation of the Chern numbers find a tiny floating-up of the extended states, while Xie at al do not find any changing in the position of extended states and affirm that the floating picture is not valid in this model. The type of boundary conditions used in the lattice model is important; the above mentioned authors use periodic conditions. The Dirichlet conditions around the plaquette, which represent a more physical situation, are used by Sheng and Weng ; their approach, consisting in numerical calculation of the conductance performed over ensembles of disordered finite plaquettes, indicates that the QH-insulator transition may occur from any IQH plateau of index $`\nu `$ directly to the insulator. In the middle of the transition region the conductances satisfy the relation $`<\sigma _{xx}>=<\sigma _{xy}>=\nu /2`$ (here $`<\mathrm{}>`$ means ensemble average) and the localization length diverges; this region is called ’metallic’.
In this context we study some new relevant features of the disordered lattice model in magnetic field with vanishing boundary conditions, the attention being paid especially to magnetization, level spacing distribution and critical conductances in the metallic region.
The discussion is based on the spin-less one-electron Hamiltonian in perpendicular magnetic field defined on a 2D square lattice with N sites in one direction and M sites along the other one, which reads as follows:
$$H=\underset{n=1}{\overset{N}{}}\underset{m=1}{\overset{M}{}}[ϵ_{nm}|n,mn,m|+^{i2\pi \varphi m}|n,mn+1,m|$$
$$+|n,mn,m+1|+h.c.](1)$$
where $`|n,m`$ is a set of orthonormal states, localized at the sites (n,m), and $`\varphi `$ is the magnetic flux through the unit cell measured in quantum flux units. In Eq.(1) the hopping integral at $`\varphi =0`$ is taken unity serving as the energy unit and the diagonal energy $`ϵ_{nm}`$ is a random variable distributed according to the probability density:
$$P(ϵ)=\{\begin{array}{cc}1/W,W/2<ϵ<W/2\hfill & \\ 0,otherwise,(2)\hfill & \end{array}$$
(having zero mean $`\overline{ϵ}=0`$).
The averaged spectrum of Hamiltonian (1) is depicted in Fig.1 for $`\varphi =1/10`$ and the disorder amplitude in the range $`W[0,10]`$. The lines represent the mean eigenvalues $`<E_n>`$ as function of disorder amplitude $`W`$.
In order to study the phase diagram we calculate the longitudinal and Hall conductances of this system at constant magnetic flux and varying disorder. The different phases: quantum Hall, metallic, insulating are characterized not only by conductance but also by the specific distribution of the level spacing and by the current density on the plaquette, described by the operator:
$$𝐉_{nm}^{n^{}m^{}}=it_{nm}^{n^{}m^{}}(𝐫_{nm}𝐫_{n^{}m^{}})|nmn^{}m^{}|+h.c(4)$$
(here $`t_{nm}^{n^{}m^{}}`$ is the hopping integral between the sites $`𝐫_{nm}`$ and $`𝐫_{n^{}m^{}}`$).
It is opportune to remind previously that for a clean system ($`ϵ_{nm}=0`$) with cyclic boundary conditions (i.e for a torus) and commensurate values of the magnetic flux through the unit cell, the spectrum consists of degenerate bands separated by gaps (the well-known Hofstadter butterfly). However, when vanishing boundary conditions are imposed (i.e., for plaquette or cylinder geometry) the gaps get filled with ’edge states’, localized close to the edges of the sample. The other states, the ’bulk’ ones, remain grouped in bands on the energy scale, while geometrically are concentrated in the middle of the plaquette. The two types of states differ also by their chirality, i.e. by the sign of the derivative $`dE_n/d\varphi `$.
The effect on the orbital magnetization of each state is immediate: the expectation values of the operator $`𝐌=(𝐫\times 𝐣(𝐫))𝑑S`$ calculated on the eigenstates of the Hamiltonian (1) have different signs depending on whether the state is bulk- or edge-type.
Fig.2 shows that the magnetization of the edge eigenstate No.11 is positive $`M_{11}>0`$, but the bulk eigenstate No.12 has $`M_{12}<0`$ ; the local currents corresponding to the two states are also shown in insets. In the same figure one anticipates that the increasing disorder produces a monotonic decrease of the magnetization. The magnetization of all states in the spectrum is shown in Fig.3a and b, for the clean and disordered system, respectively (the well-known electron-hole symmetry of the Hofstadter spectrum is evident also in the aspect of the magnetization). The disorder effect consisting in the broadening of the bands and narrowing of the gaps is obvious in the second figure. More notable is the magnetization of the ground state $`M_g`$ which can be compared successfully with experimental results. Assuming that the spectrum is filled up to the Fermi energy $`E_f`$, due to the alternating sign of $`M_n`$ in different regions of the spectrum, the quantity
$$M_g=\underset{(E_n<E_f)}{}M_n(5)$$
shows, as function of the number of occupied states, a sawtooth aspect (see Fig.4) which is the same as in the experiments by Wiegers at al , including the fact that the jumps of $`M_g`$ occur at the center of the gaps. In our model, the number of teeth depends on the number of gaps that can be resolved. For modulated quantum wires such a sawtooth aspect was obtained also for the thermodynamic magnetization .
When the Anderson potential (Eq.2) is switched on and $`W`$ is increased continuously, the bands become broader and broader, the disorder spreads the states over the whole plaquette, giving rise to extended disordered states, and finally produces a quasi-continuum of localized states; even the edge states, which are more robust, disappear gradually into the quasi-continuum. The nature of the states can be checked by calculating the distribution of level spacing for various degrees of disorder, in different domains of the spectrum. Let $`s_n`$ be the level spacing between two consecutive eigenvalues $`E_n`$ and $`E_{n+1}`$ and define $`t_n=s_n/s_n`$, where $`s_n`$ is the mean level spacing.
In Fig.5 we have three typical distribution functions $`P(t)`$ relevant for different values of disorder and energy interval. The Wigner-Dyson (WD) surmise with $`\beta =2`$ (unitary) indicates the presence of extended states (Fig.5a) and the Poisson-like distribution (Fig.5c) shows the existence of uncorrelated localized states for strong disorder. The case of edge states is also studied, in which situation the distribution can be fitted well with the Gaussian function (Fig.5b). Every inset shows a representative distribution of local currents for each case. Very illustrative is Fig.6 which shows the variance $`\delta t_n`$ for all level spacing of the Hamiltonian (1). One may learn that: a) at $`W<4`$, for most of the states, $`\delta t_n0.42`$, which is the typical value for WD distribution with $`\beta =2`$. b) for larger $`W`$, the variance increases towards $`\delta t_n=1.0`$ specific to the Poisson distribution. This value cannot practically be reached because of finite dimension of the plaquette. c) inbetween, at $`W6`$, the variance equals 0.52, and the probability distribution is well-fitted by the WD function with $`\beta =1`$ (so, here the influence of the magnetic field is lost and the system behaves ’orthogonal’). d) the lowest states, originating from the first band ( n=1,.,6 ) get localized faster than the others, while the states from the first gap (n=7,.,10) are very robust against the localization process.
At last we discuss the transport properties; for this purpose the Landauer-Büttiker formalism and the technics from Ref.9 are used. The conductance in the three regimes (IQH, metal and insulator) may be correlated with the spectrum characteristics of the isolated system discussed above. Since the edge states are responsible for IQH effect and they are robust against disorder, this regime survives also in the presence of disorder as long as $`W`$ is not too high. In what concerns the transmittances between different leads, it is well-known that the only non-vanishing ones are $`<T_{\alpha \alpha +1}>`$, which connect consecutive leads ($`\alpha `$=lead index) and equal an integer value.
The metallic regime occurs when the transport process takes place on the states which are extended over the whole plaquette. In this case also the local current is distributed on the whole area allowing for a non-zero transmission probability between any pair of leads. This behaviour corroborates with the multi-fractal properties of the local density of states . The metallic region can be crossed in different ways. Let assume a constant disorder (say $`W`$ =1.0) and change the Fermi level ; in this case, metallic regions are intercalated between the QH regions, excepting the lowest one which ensures the transition between QH1 and the insulating phase. The consequence is that, when crossing the metallic region, the Hall conductance $`\sigma _{xy}`$ will have a smooth decay between consecutive plateaus, while $`\sigma _{xx}`$ will differ from zero and have a maximum in the middle of the transition. In Fig.4, one has to notice that $`R_{xx}`$ equals $`R_{xy}`$ in two places: 1) at the transition between QH1 and insulator, where one has also $`\sigma _{xx}1/2`$ in agreement with Huo at al , and 2) near the centre of the spectrum where the system behaves already classically.
Another way to look at the metal-insulator transition consists in crossing the metallic zone by increasing the amplitude of disorder $`W`$, while keeping constant the number of electrons ($`N_e`$) or the Fermi level. Now the transition is of the type $`\nu 0`$ and exhibits maxima of the longitudinal conductance at a critical disorder $`W_c`$ where, as in Ref.8 , the condition $`\sigma _{xx}=\sigma _{xy}`$ is fulfilled. Such kind of transitions are shown in the insets of Fig.6 for $`N_e`$=8 (i.e. in the first gap) and for $`N_e`$=20 (i.e. in the second gap). One may ask what is going on when $`N_e`$ corresponds to a band ; this is the situation for the third inset $`N_e`$=14. It can be observed that $`\sigma _{xx}`$ is different from zero even at small $`W`$, meaning that the bulk states are the first ones which become extended under the influence of disorder. The dotted line which crosses transversally the spectrum in the metallic region of Fig.1 represents the critical disorder $`W_c`$ corresponding to each $`N_e`$.
In the insulating phase, when the Poisson distribution of the level spacing is installed, the Hall and longitudinal conductances tend to zero and the longitudinal resistance increases exponentially. The transport of the electron through the plaquette is performed by tunneling on localized states.
In conclusion, the metallic regime is characterized by a Wigner-Dyson distribution of the level spacing with $`\beta =2`$ (unitary ensemble). As the system evolves towards insulator, the orthogonal WD distribution ($`\beta =1`$) that showes up at a given higher disorder indicates the loss of influence of the magnetic field. Simultaneously, the magnetization decays to zero. The QH phase is characterized by a Gaussian distribution of level spacing. Due to the different chirality of the edge and extended states, when crossing the metallic zone (at relatively small disorder) the magnetization of the ground state shows a toothsaw behavior as function of the filling factor similar to experimental results.
Acknowledgments. A.A. is very grateful to Professor Johannes Zittartz for his hospitality at Institut für Theoretische Physik der Universität Köln where this work was partially performed under SFB 341. Illuminating discussions with J.Hajdu, B.Huckestein, M.Zirnbauer, K.Maschke and A.Manolescu are thankfully acknowledged. We thank Romanian Academy for the support under the grant No.69/1999.
Permanent address: National Institute of Materials Physics POBox MG7, Bucharest-Magurele, Romania.
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# Static and dynamic image potential for tunneling into a Luttinger liquid
\[
## Abstract
We study electron tunneling from a tip or a lead into an interacting quantum wire described by Luttinger liquid theory. Within a WKB-type approach, the Coulomb interaction between the wire and the tunneling electrons, as well as the finite traversal time are taken into account. Although the static image potential is only logarithmically suppressed against the bare Coulomb interaction, the dynamic image potential is not strong enough to alter power-law exponents entering the tunneling density of states.
\]
One-dimensional (1D) quantum wires (QWs) are at the focal point of current activities in condensed-matter physics. Fabrication advances in semiconductor heterostructures and carbon nanotubes allow the systematic study of phenomena arising only in one dimension. In particular, the physics of 1D nanowires is intimately connected with the concept of a Luttinger liquid (LL), exhibiting spin-charge separation, suppression of the tunneling density of states (TDOS), and interaction dependent power-laws in transport properties. For nanotubes, the theoretically predicted LL behavior has been convincingly established in several recent experiments. Of particular importance was the measurement of the TDOS power-law exponent $`\alpha `$, which provides information about the dimensionless interaction strength parameter $`g`$ of the LL (where $`0<g1`$). The TDOS can be obtained, e.g. from the tunneling current measured via weakly attached metallic leads. Another possibility is to use a scanning tunneling microscopy (STM) tip. Apart from the resolution of atomic and electronic properties of individual nanotubes successfully achieved in recent STM experiments, a detailed analysis of nonlinear current-voltage curves in the spectroscopy (STS) mode would allow to extract the TDOS exponent.
So far in all theoretical studies of tunneling into a LL, it has tacitly been assumed that the electrons tunneling into the correlated fermion liquid do not modify the effective TDOS. Conventional treatments employ the tunneling Hamiltonian and thereby assume the traversal time of tunneling to be infinitely short. In addition, the interaction between the tunneling electron and the LL electrons is neglected. Given these two assumptions, the TDOS is indeed an intrinsic property of the LL with a power-law dependence on energy, $`\rho (ϵ)ϵ^\alpha ,`$ for $`ϵ0`$, so that $`\alpha `$ depends only on $`g`$. (For clarity, we focus on the case of bulk tunneling at zero temperature throughout this paper.) For an STS experiment, such a calculation predicts that the measured differential conductance is proportional to the TDOS, and hence the exponent $`\alpha `$ can be extracted from experimental data. A similar reasoning has been employed to understand the data in Ref. , where electron tunneling between metallic leads and a nanotube was important.
Here we address the question of whether this measured exponent really characterizes the unperturbed LL, or whether it is affected by the dynamics and charge of the tunneling particle. To that end, we take into account correlations both within the QW and between the QW and the tunneling electron, as well as a finite traversal time. Our main findings are as follows: Although the static image potential experienced by the tunneling electron is very strong, dynamical effects turn out to be of crucial importance. Within the framework of a semiclassical theory related to Nazarov’s tunnel junction theory, the low-energy power-law exponent of the TDOS is governed solely by intrinsic LL properties. Thereby we provide the a posteriori justification for the (previously assumed) connection linking the (observable) value of $`\alpha `$ and the interaction parameter $`g`$. Since spin and charge are decoupled in a LL, we study only the spinless single-channel case, with the same conclusions applying to spin-$`\frac{1}{2}`$ electrons or nanotubes.
For the case of tunneling from an STM tip into the QW, the relevant geometry is depicted in Figure 1. We consider a clean and very long $`(L\mathrm{})`$ QW, and mainly focus on an interaction potential of the form
$$V(x,y)=\frac{1}{\sqrt{x^2+(Dy)^2+a^2}},$$
(1)
where $`a`$ is the lattice constant of the QW. For $`y=D`$, this describes the intra-wire interaction responsible for the LL state, while for $`0<y<D`$, it gives the interaction between the wire electrons and the tunneling electron at $`x=0`$ (see below). The one-sided Fourier transform of Eq. (1) is
$$\stackrel{~}{V}(q,y)=2K_0\left(|q|\sqrt{(Dy)^2+a^2}\right),$$
(2)
with the modified Bessel function $`K_0`$. When using Eq. (1), one neglects or crudely approximates several important effects, e.g. screening due to the tip, or the orbital structure of both tip and quantum wire (see, e.g. Ref. ). However, our qualitative findings concerning the connection between $`\alpha `$ and $`g`$ are expected to hold quite generally. In addition to the model potential (1), these findings can be made rigorous for a general class of separable potentials.
Tunneling out of the tip (or lead) proceeds from the ground state with energy $`E_0^{}`$ into the LL ground state with energy $`E_0`$, and is usually described in terms of a simple tunneling Hamiltonian,
$$H_T=T\mathrm{\Psi }^{}(0)c_0+\mathrm{h}.\mathrm{c}.,$$
(3)
where $`c_0`$ annihilates an electron in a state at the center of the tip and $`\mathrm{\Psi }^{}(0)`$ creates an electron at $`x=0`$ in the LL. If the tunneling matrix element $`T`$ is small, $`H_T`$ can be treated as a perturbation to $`H^{}+H`$, where $`H^{}`$ describes the electrons in the tip and $`H`$ denotes the Hamiltonian of the LL. More specifically, the perturbative treatment is appropriate when the tunneling resistance is large compared to $`h/e^2`$. Then the tunneling rate $`\mathrm{\Gamma }`$ can be calculated with the aid of Fermi’s golden rule. Labeling states in the tip by $`|\nu `$ and states in the LL by $`|n`$, the tunneling rate is (we put $`\mathrm{}=1`$)
$`\mathrm{\Gamma }=2\pi T^2{\displaystyle \underset{\nu ,n}{}}|\nu ,n|\mathrm{\Psi }^{}(0)c_0|0,0|^2\delta (E_0^{}+E_0E_\nu ^{}E_n),`$
where $`|0,0`$ denotes the ground state in the absence of tunneling. Using the identity
$`\delta (E_0^{}+E_0E_\nu ^{}E_n)={\displaystyle 𝑑ϵ\delta (E_0+ϵE_n)\delta (E_0^{}ϵE_\nu ^{})},`$
the rate may be written as
$$\mathrm{\Gamma }=2\pi T^2_0^{eV}𝑑ϵ\rho (ϵ)\rho ^{}(eVϵ).$$
(4)
Here the TDOS for adding an electron with energy $`\mu +ϵ`$ to the LL is
$$\rho (ϵ)=\underset{n}{}|n|\mathrm{\Psi }^{}(0)|0|^2\delta (E_0+\mu +ϵE_n),$$
(5)
where the electrochemical potential $`\mu `$ is the minimal energy required to add an electron. Furthermore,
$$\rho ^{}(ϵ^{})=\underset{\nu }{}|\nu |c_0|0|^2\delta (E_0^{}\mu ^{}+ϵ^{}E_\nu ^{})$$
(6)
is the DOS for removing an electron from the tip/lead with energy $`\mu ^{}ϵ^{}`$. Now $`\rho (ϵ)`$ and $`\rho ^{}(ϵ)`$ are nonvanishing only for $`ϵ>0`$, and $`\mu ^{}\mu =eV`$ determines the applied voltage. With these definitions we readily obtain Eq. (4). In a metallic lead, the DOS $`\rho ^{}(ϵ)`$ is essentially constant, while for an STM tip, we expect pronounced peaks reflecting the discrete level structure. In the latter case, Eq. (4) reproduces the rate from an STM tip into a metal first obtained by Tersoff and Hamann.
To determine $`\rho (ϵ)`$ explicitly, we write
$`\rho (ϵ)`$ $`=`$ $`{\displaystyle \frac{\mathrm{Re}}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle \underset{n}{}}0|\mathrm{\Psi }(0)|nn|\mathrm{\Psi }^{}(0)|0e^{i(E_0+\mu +ϵE_n)t}`$ (7)
$`=`$ $`{\displaystyle \frac{\mathrm{Re}}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑tG_0(t)e^{i(\mu +ϵ)t},`$ (8)
with the single-electron Greens function ($`t>0`$) for an electron at position $`x=0`$ in the ground state, $`G_0(t)=\mathrm{\Psi }(0,t)\mathrm{\Psi }^{}(0,0)`$. To evaluate $`G_0(t)`$, standard bosonization methods can be applied. The kinetic part of the Hamiltonian is
$`H_0={\displaystyle \frac{v_F}{2}}{\displaystyle 𝑑x\left[(_x\vartheta )^2+(_x\phi )^2\right]}+\mu \widehat{N}.`$
Here $`v_F`$ is the Fermi velocity, $`\widehat{N}`$ the particle number operator, and $`\vartheta (x)`$ is conjugate to the phase field $`\phi (x)`$ describing the plasmon excitations in the wire. The intra-wire interaction part is $`H_V=\frac{1}{2}𝑑x𝑑x^{}n(x)V(xx^{},D)n(x^{})`$ with the electron charge density $`n(x)=e\pi ^{1/2}_x\phi `$. (The LL model appropriate for the low energy sector follows by effectively using a local interaction, $`V(x,D)=V_0\delta (x)`$.) By virtue of a Bogoliubov transformation, $`H=H_0+H_V`$ can easily be diagonalized. With bosonic operators $`b_q^{()}`$, the phase field $`\phi (x)`$ reads
$$\phi (x)=i\underset{q0}{}\left(\frac{g(q)}{2L|q|}\right)^{1/2}\mathrm{exp}(iqx)\mathrm{sgn}(q)[b_q^{}+b_q],$$
(9)
with the $`q`$-dependent interaction parameter \[where $`g=g(q=2\pi /L)`$\]
$$g(q)=[1+e^2\stackrel{~}{V}(q,D)/\pi v_F]^{1/2}=g(q).$$
(10)
We then arrive at
$$H=\underset{q0}{}\omega _qb_q^{}b_q+\mu \widehat{N},$$
(11)
with the plasmon dispersion relation $`\omega _q=v_F|q|/g(q)`$. In terms of the chiral (right- or left-moving) phase fields ($`p=R/L=\pm `$),
$$\varphi _p(x)=[p\phi (x)+\vartheta (x)]/\sqrt{4\pi },$$
(12)
the bosonized electron operator at $`x=0`$ is $`\mathrm{\Psi }(0,t)_{p=\pm }\mathrm{exp}[2\pi i\varphi _p(0,t)]`$, implying
$`G_0(t)t^{(g+g^1)/2}e^{i\mu t}`$
at long times. Hence one obtains the well-known exponent $`\alpha =(g+g^12)/2`$ governing the bulk TDOS.
Let us now look at the static image potential experienced by an electron with charge $`e`$ held fixed at position $`x=0`$ and $`0<y<D`$ due to its interaction with the QW electrons, $`H_I=e𝑑xn(x)V(x,y)=e\varphi (y)`$. Using Eq. (9), we get the fluctuating field
$$\varphi (y,t)=\underset{q0}{}\lambda _q(y)\left(b_q^{}(t)+b_q(t)\right)$$
(13)
with couplings $`\lambda _q(y)=e[g(q)|q|/2\pi L]^{1/2}\stackrel{~}{V}(q,y)`$. Next we shift the bosonic operators, $`B_q=b_qe\lambda _q(y)/\omega _q`$, whence
$$H=\underset{q0}{}\omega _qB_q^{}B_q+\mu \widehat{N}+V_{im}(y).$$
(14)
Here the static image potential
$$V_{im}(y)=e^2\underset{q0}{}\lambda _q^2(y)/\omega _q$$
(15)
describes the energy gained by the plasmons relaxing to their equilibrium state in the presence of the additional electron. For the unscreened interaction (1), one has $`g(q)=[1+\xi K_0(|q|a)]^{1/2}`$ with the dimensionless parameter $`\xi =2e^2/\pi v_F`$. Then the image potential (15) for $`(Dy)a`$ reads
$`V_{im}(y)={\displaystyle \frac{2e^2/\pi }{(Dy)\{\mathrm{ln}[(Dy)/a]+1/\xi \}}}.`$
Therefore the static image potential is only logarithmically suppressed against the bare Coulomb interaction, and hence is very strong.
Next we turn to dynamical effects due to tunneling. We envision the latter as penetration through a rectangular barrier of width $`D`$. If the barrier is sufficiently thick and its transparency low, the main contribution to the tunnel current comes from electrons with momenta perpendicular to the QW. Therefore we effectively obtain a 1D Schrödinger equation for the underbarrier motion $`\psi (y,t)`$ of the tunneling electron ($`0<y<D`$),
$$i_t\psi (y,t)=[(2m)^1_y^2+\mu ^{}+Ue\varphi (y,t)]\psi (y,t),$$
(16)
where $`U`$ is the work function of the tip/lead. In the absence of $`\varphi (y,t)`$, the solution for an electron at energy $`\mu ^{}`$ reads $`\psi (y,t)e^{i\mu ^{}tmvy}`$, where $`v=\sqrt{2U/m}`$ is an effective velocity related to the traversal time $`D/v`$. Under the WKB approximation, the dominant effect of the potential $`\varphi `$ can be incorporated as additional phase factor,
$$\psi (y,t)e^{i\mu ^{}tmvyi\theta (y,t)}.$$
(17)
Linearizing the resulting WKB equation gives for $`|\varphi (y,t)|U`$:
$$_t\theta +iv_y\theta =e\varphi (y,t),$$
(18)
supplemented by the boundary condition $`\theta (0,t)=0`$. This equation can be solved separately for each bosonic mode using the ansatz
$`\theta (y,t)={\displaystyle \underset{q0}{}}\left[w(y,\omega _q)b_q(t)+\stackrel{~}{w}(y,\omega _q)b_q^{}(t)\right].`$
From Eq. (18) and $`b_q(t)=b_qe^{i\omega _qt}`$, we obtain $`\stackrel{~}{w}(y,\omega _q)=w(y,\omega _q)`$ and
$`(v_y\omega _q)w(y,\omega _q)=ie\lambda _q(y),`$
which can easily be solved. For $`y=D`$, we finally get
$$\theta (t)=\underset{q0}{}\left[w(\omega _q)b_q^{}(t)+w(\omega _q)b_q(t)\right],$$
(19)
with
$$w(\omega _q)=(ie/v)_0^D𝑑y\lambda _q(y)\mathrm{exp}[\omega _q(Dy)/v].$$
(20)
Because of the associated dynamic image potential, the electron wave function acquires the phase factor $`\mathrm{exp}[i\theta (t)]`$ during the tunneling process.
This effect can be properly incorporated by a modification of the tunneling Hamiltonian,
$$\stackrel{~}{H}_T=T\mathrm{\Psi }^{}(0)e^{i\theta }c_0+\mathrm{h}.\mathrm{c}..$$
(21)
The Greens function determining the effective TDOS is now given by
$$G(t)=\mathrm{exp}[i\theta ^{}(t)]\mathrm{\Psi }(0,t)\mathrm{\Psi }^{}(0,0)\mathrm{exp}[i\theta (0)].$$
(22)
Putting $`G(t)=G_0(t)K(t)`$, the contribution of the dynamic image potential then gives rise to the factor
$$K(t)=\mathrm{exp}[C_1(t)+C_2(t)],$$
(23)
where we introduce the functions
$`C_1(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\theta ^2(t)+\theta ^2(0)2\theta ^{}(t)\theta (0),`$
$`C_2(t)`$ $`=`$ $`2\pi [\varphi _p(t)\theta (0)+\theta ^{}(t)\varphi _p(0)`$
$`\theta ^{}(t)\varphi _p(t)\varphi _p(0)\theta (0)],`$
where $`C_2`$ is independent of $`p=\pm `$. Doing the Gaussian averages yields
$`C_1(t)`$ $`=`$ $`{\displaystyle \underset{q0}{}}[w(\omega _q)w(\omega _q)+w^2(\omega _q)\mathrm{exp}(i\omega _qt)],`$ (24)
$`C_2(t)`$ $`=`$ $`i(2\pi /L)^{1/2}{\displaystyle \underset{q0}{}}{\displaystyle \frac{w(\omega _q)}{\sqrt{g(q)|q|}}}[1\mathrm{exp}(i\omega _qt)].`$ (25)
Since we are interested in the power-law exponent governing the TDOS, we focus on the time-dependent parts of Eqs. (24) and (25), and do not explicitly compute the prefactor. For the interaction (1), numerical calculation of $`K(t)`$ gives the result shown in Fig. 2 which is well approximated by $`|K(t)|=1+A(t)\mathrm{cos}(\mathrm{\Omega }t)`$ for long times, with oscillation frequency $`\mathrm{\Omega }`$. The amplitude decays according to $`A(t)t^\beta `$ with $`\beta 1.15`$. This result is insensitive to the precise parameter values taken for $`\xi ,D/a`$, and $`v/v_F`$. As a consequence, the power-law exponent $`\alpha `$ of the TDOS for small energy $`ϵ`$ remains unchanged by the dynamic image potential. Hence one can indeed obtain the LL parameter $`g`$ from a measurement of $`\alpha `$.
This finding can be inferred analytically for a class of separable interaction potentials of the form
$$V(x,y)=V_0\delta (x)f(y),$$
(26)
where $`f(y)`$ is an arbitrary function with $`f(D)=1`$. In this case, we get $`\stackrel{~}{V}(q,y)=V_0f(y)`$, leading to $`g(q)=g`$. The time-dependent parts of Eqs. (24) and (25) read
$`C_1(t)`$ $``$ $`{\displaystyle _0^{\mathrm{}}}𝑑qe^{iv_Fqt/g}q\left[{\displaystyle _0^D}𝑑yf(y)e^{v_F(Dy)q/gv}\right]^2,`$
$`C_2(t)`$ $``$ $`{\displaystyle _0^{\mathrm{}}}𝑑qe^{iv_Fqt/g}\left[{\displaystyle _0^D}𝑑yf(y)e^{v_F(Dy)q/gv}\right].`$
The asymptotic long-time behavior of $`C_{1,2}(t)`$ can be accurately calculated in stationary-phase approximation. We find that $`C_1(t)`$ decays faster than $`1/t`$, while $`C_2(t)1/t`$. Therefore, from Eq. (23), the TDOS exponent for small $`ϵ`$ remains unchanged. We expect this result to be correct and generic for arbitrary physically relevant interaction potentials.
We conclude by summarizing our results. We have presented a simple theory of electron tunneling from a tip or a lead into a strongly correlated 1D metal, explicitly incorporating the finite traversal time and the dynamic response of the correlated metal to the incoming electron. We have solved this problem within a WKB-type approximation for different interaction potentials. Despite the presence of a strong static image potential, the power-law exponent entering the tunneling density of states is not affected by these effects, but completely determined by the correlation strength in the 1D metal.
We acknowledge financial support by the Deutsche Forschungsgemeinschaft under the Gerhard-Hess program.
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# Neutrino Oscillation and Charged Lepton-Flavor Violation in the Supersymmetric Standard Models
## 1 Introduction
Result of the atmospheric neutrino experiment by the superKamiokande detector indicates that the neutrinos have finite masses and the lepton-flavor symmetry of muon is violating in nature $`^\mathrm{?}`$. This is the first signature of the physics beyond the standard model (SM), and this discovery will be confirmed by further experiments, such as the long base-line experiments. Also, the solar neutrino experimental data suggest that the lepton-flavor symmetry of electron is violating $`^\mathrm{?}`$.
Processes, such as $`\mu ^+e^+\gamma `$ and $`\tau ^\pm \mu ^\pm \gamma `$, are also lepton-flavor violating (LFV) processes. Unfortunately, the event rates are too small to be observed in near future experiments even if the neutrino masses are introduced into the standard model. The event rates are suppressed by the fourth order of the ratio of the tiny neutrino mass to the $`W`$ boson mass due to the GIM suppression. However, if the standard model is supersymmetrized, the processes may be accessible in near future experiment and we may study the origin of the neutrino masses.
The supersymmetric standard model (SUSY SM) is a solution of the naturalness problem, and is one of the most promising model beyond the standard model. In this model, introduction of the SUSY breaking terms allows the lepton-flavor symmetries to be violating in the slepton masses $`^\mathrm{?}`$. Then, the orders of magnitude of the event rates for the LFV processes depend on the origin of the SUSY breaking in the SUSY SM and physics beyond the SUSY SM. One of the successful candidates for the origin of the SUSY breaking is the minimal supergravity, where the SUSY breaking scalar masses are generated universally in the flavor space at the tree level. In this scenario, the source of the LFV processes comes from the LFV radiative correction to the SUSY breaking masses for the sleptons by the LFV interaction in physics beyond the SUSY SM $`^\mathrm{?}`$. Then, we have a chance to study the origin of the neutrino masses through the LFV processes in the supersymmetric models.
The see-saw mechanism $`^\mathrm{?}`$ is the simplest model to generate the tiny neutrino masses. In this mechanism the Yukawa interactions are lepton-flavor violating due to introduction of the right-handed neutrinos, similar to the quark sector. Then, in the minimal supergravity scenario, if the lepton-flavor violation in the Yukawa coupling constants is strong enough, the radiative correction generates sizable LFV masses for the sleptons $`^\mathrm{?}`$$`^\mathrm{?}`$. Moreover, the large mixing angles observed on the solar and atmospheric neutrino observations may enhance the event rates for the LFV processes $`^\mathrm{?}`$$`^\mathrm{?}`$$`^\mathrm{?}`$.
In this article, we study the charged lepton-flavor violating processes, $`\mu ^+e^+\gamma `$ and $`\tau ^\pm \mu ^\pm \gamma `$, using the current neutrino experimental data, in the SUSY SM with the right-handed neutrinos.<sup>a</sup><sup>a</sup>a After sleptons are discovered in the future large colliders, the slepton oscillation will be a powerful tool to study the lepton-flavor violation $`^\mathrm{?}`$. We assume the minimal supergravity scenario. The large mixing angle in the atmospheric neutrino result may enhance $`\tau ^\pm \mu ^\pm \gamma `$, and the large mixing angles in the MSW and the vacuum oscillation solutions may lead to a large event rate of $`\mu ^+e^+\gamma `$.
This article is organized as follows. In the next section we review the radiative generation of the LFV masses for slepton in the SUSY SM with the right-handed neutrinos. In Section 3 we show the branching rate for $`\tau ^\pm \mu ^\pm \gamma `$, using the atmospheric neutrino result. In Section 4 discuss the branching rate for $`\mu ^+e^+\gamma `$, using the solar neutrino result. The other processes are also discussed. Section 5 is devoted to Conclusion.
## 2 The SUSY SM with the right-handed neutrinos
We review the radiative generation of the LFV masses for sleptons in the SUSY SM with the right-handed neutrinos. We adopt the minimal supergravity scenario as the origin of the SUSY breaking in the SUSY SM. The superpotential of the lepton sector in the SUSY SM with right-handed neutrinos is given as
$`W`$ $`=`$ $`f_{\nu _{ij}}H_2\overline{N}_iL_j+f_{e_{ij}}H_1\overline{E}_iL_j+{\displaystyle \frac{1}{2}}M_{\nu _i\nu _j}\overline{N}_i\overline{N}_j,`$ (1)
where $`L`$ is a chiral superfield for the left-handed lepton, and $`\overline{N}`$ and $`\overline{E}`$ are for the right-handed neutrino and the charged lepton. $`H_1`$ and $`H_2`$ are for the Higgs doublets in the SUSY SM. Here, $`i`$ and $`j`$ are generation indices. After redefinition of the fields, the Yukawa coupling constants and the Majorana masses can be taken as
$`f_{\nu _{ij}}`$ $`=`$ $`f_{\nu _i}V_{Dij},`$
$`f_{e_{ij}}`$ $`=`$ $`f_{e_i}\delta _{ij},`$
$`M_{\nu _i\nu _j}`$ $`=`$ $`U_{ik}^{}M_{\nu _k}U_{kj}^{},`$ (2)
where $`V_D`$ and $`U`$ are unitary matrices. In this model the mass matrix for the left-handed neutrinos $`(m_\nu )`$ becomes
$`(m_\nu )_{ij}`$ $`=`$ $`V_{Dik}^{}(\overline{m}_\nu )_{kl}V_{Dlj},`$ (3)
where
$`(\overline{m}_\nu )_{ij}`$ $`=`$ $`m_{\nu _iD}\left[M^1\right]_{ij}m_{\nu _jD}`$ (4)
$``$ $`V_{Mik}^{}m_{\nu _k}V_{Mkj}.`$
Here, $`m_{\nu _iD}=f_{\nu _i}v\mathrm{sin}\beta /\sqrt{2}`$ and $`V_M`$ is a unitary matrix.<sup>b</sup><sup>b</sup>b $`h_1=(v\mathrm{cos}\beta /\sqrt{2},0)^{}`$ and $`h_2=(0,v\mathrm{sin}\beta /\sqrt{2})^{}`$ with $`v246`$GeV. The observed mixing angles on the atmospheric and solar neutrino experiments are $`(V_MV_D)_{\tau \mu }`$ and $`(V_MV_D)_{\mu e}`$, respectively, if they come from the oscillations of $`\nu _\mu \nu _\tau `$ and $`\nu _e\nu _\mu `$.
The SUSY breaking terms for the lepton sector in the SUSY SM with the right-handed neutrinos are in general given as
$`_{\mathrm{SUSY}\mathrm{breaking}}`$ $`=`$ $`(m_{\stackrel{~}{L}}^2)_{ij}\stackrel{~}{l}_{Li}^{}\stackrel{~}{l}_{Lj}+(m_{\stackrel{~}{e}}^2)_{ij}\stackrel{~}{e}_{Ri}^{}\stackrel{~}{e}_{Rj}+(m_{\stackrel{~}{\nu }}^2)_{ij}\stackrel{~}{\nu }_{Ri}^{}\stackrel{~}{\nu }_{Rj}`$ (5)
$`+(A_\nu ^{ij}h_2\stackrel{~}{\nu }_{Ri}^{}\stackrel{~}{l}_{Lj}+A_e^{ij}h_1\stackrel{~}{e}_{Ri}^{}\stackrel{~}{l}_{Lj}+{\displaystyle \frac{1}{2}}B_\nu ^{ij}\stackrel{~}{\nu }_{Ri}^{}\stackrel{~}{\nu }_{Rj}^{}+h.c.),`$
where $`\stackrel{~}{l}_L`$, $`\stackrel{~}{e}_R`$, and $`\stackrel{~}{\nu }_R`$ represent the left-handed slepton, and the right-handed charged slepton, and the right-handed neutrino. Also, $`h_1`$ and $`h_2`$ are the doublet Higgs bosons. In the minimal supergravity scenario the SUSY breaking masses for sleptons, squarks, and the Higgs bosons are universal at the gravitational scale ($`M_{\mathrm{grav}}10^{18}`$GeV), and the SUSY breaking parameters associated with the supersymmetric Yukawa couplings or masses ($`A`$ or $`B`$ parameters) are proportional to the Yukawa coupling constants or masses. Then, the SUSY breaking parameters in Eq. (5) are given as
$`(m_{\stackrel{~}{L}}^2)_{ij}=(m_{\stackrel{~}{e}}^2)_{ij}=(m_{\stackrel{~}{\nu }}^2)_{ij}=\delta _{ij}m_0^2,`$
$`A_\nu ^{ij}=f_{\nu _{ij}}a_0,A_e^{ij}=f_{e_{ij}}a_0,`$
$`B_\nu ^{ij}=M_{\nu _i\nu _j}b_0,`$ (6)
at the tree level.
In order to know the values of the SUSY breaking parameters at the low energy, we have to include the radiative corrections to them. While we evaluate them by solving the RGE’s, we discuss only the qualitative behavior of the solution using the logarithmic approximation here. The SUSY breaking masses of squarks, sleptons, and the Higgs bosons at the low energy become heavier by gauge interactions at one-loop level, and the corrections are flavor-independent. On the other hand, Yukawa interactions reduce the diagonal SUSY breaking mass squareds and the radiative corrections are flavor-dependent. Then, if the Yukawa coupling is lepton-flavor violating, the radiative correction to the SUSY breaking parameters is also lepton-flavor violating. The LFV off-diagonal components for $`(m_{\stackrel{~}{L}}^2)`$, $`(m_{\stackrel{~}{e}}^2)`$, and $`A_e`$ in the SUSY SM with the right-handed neutrinos are given at the low energy as
$`(m_{\stackrel{~}{L}}^2)_{ij}`$ $``$ $`{\displaystyle \frac{1}{8\pi ^2}}(3m_0^2+a_0^2)V_{Dki}^{}V_{Dlj}f_{\nu _k}f_{\nu _l}U_{km}^{}U_{lm}\mathrm{log}{\displaystyle \frac{M_{\mathrm{grav}}}{M_{\nu _m}}},`$
$`(m_{\stackrel{~}{e}}^2)_{ij}`$ $``$ $`0,`$
$`A_e^{ij}`$ $``$ $`{\displaystyle \frac{3}{8\pi ^2}}a_0f_{e_i}V_{Dki}^{}V_{Dlj}f_{\nu _k}f_{\nu _l}U_{km}^{}U_{lm}\mathrm{log}{\displaystyle \frac{M_{\mathrm{grav}}}{M_{\nu _m}}},`$ (7)
where $`ij`$. In these equations, the off-diagonal components of $`(m_{\stackrel{~}{L}}^2)`$ and $`A_e`$ are generated radiatively while those of $`(m_{\stackrel{~}{e}}^2)`$ are not. This is because the right-handed leptons have only one kind of the Yukawa interaction $`f_e`$ and we can always take a basis where $`f_e`$ is diagonal. This is a characteristic of the SUSY SM with the right-handed neutrinos. <sup>c</sup><sup>c</sup>c In the minimal SU(5) SUSY GUT, the right-handed sleptons receive the LFV masses through the Yukawa interaction of colored Higgs, but not the left-handed ones $`^\mathrm{?}`$ $`^\mathrm{?}`$. In the SO(10) SUSY GUT and the non-minimal SU(5) SUSY GUT both sleptons may have the LFV masses $`^\mathrm{?}`$.
The magnitudes of the off-diagonal components of $`(m_{\stackrel{~}{L}}^2)`$ and $`A_e`$ are sensitive to $`f_{\nu _i}`$ and $`V_D`$, while not to $`U`$. This is because the off-diagonal components of $`U`$ are small when the hierarchy among the right-handed neutrino masses is large, and then we will take $`U=\mathrm{𝟏}`$ in the following discussion for simplicity. In the following sections we will evaluate the values of $`f_{\nu _i}`$ and $`V_D`$ from the neutrino oscillation data.
## 3 The atmospheric neutrino result and $`\tau ^\pm \mu ^\pm \gamma `$
In this section we discuss the branching ratio of $`\tau ^\pm \mu ^\pm \gamma `$ using the atmospheric neutrino result. From the zenith-angle dependence of $`\nu _e`$ and $`\nu _\mu `$ fluxes measured by the superKamiokande it is natural that the atmospheric neutrino anomaly comes from the neutrino oscillation between $`\nu _\mu `$ and $`\nu _\tau `$, and the neutrino mass-squared difference and mixing angle are expected as
$`\mathrm{\Delta }m_{\nu _\mu \nu _\tau }^210^{(23)}\mathrm{eV}^2,`$
$`\mathrm{sin}^22\theta _{\nu _\mu \nu _\tau }^{}{}_{}{}^{>}0.8.`$ (8)
Assuming that the neutrino masses is hierarchical as $`m_{\nu _\tau }m_{\nu _\mu }m_{\nu _e}`$, the tau neutrino mass is given as
$`m_{\nu _\tau }(3\times 10^21\times 10^1)\mathrm{eV},`$ (9)
and if the tau neutrino Yukawa coupling constant $`f_{\nu _\tau }`$ is as large as that of the top quark, the right-handed tau neutrino $`M_{\nu _\tau }`$ is about $`10^{1415}`$GeV.
In order to evaluate the event rate for $`\tau ^\pm \mu ^\pm \gamma `$, we have to know the value of $`V_{D\tau \mu }`$, which is not necessary the same as the $`\mathrm{sin}\theta _{\nu _\mu \nu _\tau }`$. However, it is expected that it is also of the order of one as explained bellow.
Let us consider only the tau and the mu neutrino masses for simplicity. In this case we parameterize two unitary matrices $`V_D`$ and $`V_M`$ as
$`V_D=\left(\begin{array}{cc}\mathrm{cos}\theta _D& \mathrm{sin}\theta _D\\ \mathrm{sin}\theta _D& \mathrm{cos}\theta _D\end{array}\right),`$ $`V_M=\left(\begin{array}{cc}\mathrm{cos}\theta _M& \mathrm{sin}\theta _M\\ \mathrm{sin}\theta _M& \mathrm{cos}\theta _M\end{array}\right).`$ (14)
The observed large angle $`\theta _{\nu _\mu \nu _\tau }`$ is a sum of $`\theta _D`$ and $`\theta _M`$. However, in order to derive $`\theta _M\pi /4`$ we need a fine-tune among the independent Yukawa coupling constants and the mass parameters. The neutrino mass matrix $`(\overline{m}_\nu )`$ for the second and the third generations (Eq. (3)) is given as
$`(\overline{m}_\nu )`$ $``$ $`\left(\begin{array}{cc}\frac{m_{\nu _\mu D}^2}{M_{\nu _\mu \nu _\mu }}& \frac{m_{\nu _\mu D}m_{\nu _\tau D}}{M_{\nu _\mu \nu _\tau }}\frac{M_{\nu _\mu \nu _\tau }^2}{M_{\nu _\mu \nu _\mu }M_{\nu _\tau \nu _\tau }}\\ \frac{m_{\nu _\mu D}m_{\nu _\tau D}}{M_{\nu _\mu \nu _\tau }}\frac{M_{\nu _\mu \nu _\tau }^2}{M_{\nu _\mu \nu _\mu }M_{\nu _\tau \nu _\tau }}& \frac{m_{\nu _\tau D}^2}{M_{\nu _\tau \nu _\tau }}\end{array}\right).`$ (17)
If the following relations are valid,
$$\frac{m_{\nu _\tau D}^2}{M_{\nu _\tau \nu _\tau }}\frac{m_{\nu _\mu D}^2}{M_{\nu _\mu \nu _\mu }}\frac{m_{\nu _\mu D}m_{\nu _\tau D}}{M_{\nu _\mu \nu _\tau }},$$
(18)
$`m_{\nu _\tau }m_{\nu _\mu }`$ and $`\theta _M\pi /4`$ can be derived. However, the relation among the independent coupling constants and masses is not natural without some mechanism or symmetry. Also, if $`m_{\nu _\tau D}m_{\nu _\mu D}`$ similar to the quark sector, the mixing angle $`\theta _M`$ tends to be suppressed as
$$\mathrm{tan}2\theta _M2\left(\frac{m_{\nu _\mu D}}{m_{\nu _\tau D}}\right)\left(\frac{M_{\nu _\mu \nu _\tau }}{M_{\nu _\mu \nu _\mu }}\right).$$
(19)
Therefore, in the following discussion we assume that the large mixing angle between $`\nu _\tau `$ and $`\nu _\mu `$ comes from $`\theta _D`$ and that $`V_M`$ is a unit matrix.
Large $`V_{D\tau \mu }`$ leads to the non-vanishing $`(m_{\stackrel{~}{L}}^2)_{\tau \mu }`$ and $`A_e^{\tau \mu }`$, which result in a finite $`\tau ^\pm \mu ^\pm \gamma `$ event rate via diagrams involving them. They are given as
$`(m_{\stackrel{~}{L}}^2)_{\tau \mu }`$ $``$ $`{\displaystyle \frac{1}{16\pi ^2}}(3m_0^2+a_0^2)\mathrm{sin}2\theta _Df_{\nu _\tau }^2\mathrm{log}{\displaystyle \frac{M_{\mathrm{grav}}}{M_{\nu _\tau }}},`$
$`A_e^{\tau \mu }`$ $``$ $`{\displaystyle \frac{3}{16\pi ^2}}a_0\mathrm{sin}2\theta _Df_\tau f_{\nu _\tau }^2\mathrm{log}{\displaystyle \frac{M_{\mathrm{grav}}}{M_{\nu _\tau }}}.`$ (20)
As will be shown, if $`f_{\nu _\tau }`$ is of the order of one, the branching ratio of $`\tau \mu \gamma `$ may reach the present experimental bound.
Let us evaluate the branching ratios of $`\tau \mu \gamma `$. The amplitude of the $`e_i^+e_j^+\gamma `$ ($`i>j`$) takes a form
$`T=eϵ^\alpha (q)\overline{v}_i(p)i\sigma _{\alpha \beta }q^\beta (A_L^{(ij)}P_L+A_R^{(ij)}P_R)v_j(pq),`$ (21)
where $`p`$ and $`q`$ are momenta of $`e_i`$ and photon, and the event rate is given by
$`\mathrm{\Gamma }(e_ie_j\gamma )={\displaystyle \frac{e^2}{16\pi }}m_{e_i}^3(|A_L^{(ij)}|^2+|A_R^{(ij)}|^2).`$ (22)
Here, we neglect the mass of $`e_j`$. The amplitude is not invariant for the SU(2)<sub>L</sub> and U(1)<sub>Y</sub> symmetries and the chiral symmetries of leptons. Then the coefficients $`A_L^{(ij)}`$ and $`A_R^{(ij)}`$ are proportional to the charged lepton masses. Since the mismatch between the left-handed slepton and the charged lepton mass eigenstates is induced in the SUSY SM with the right-handed neutrinos, $`A_L^{(ij)}`$ is much larger than $`A_R^{(ij)}`$. Also, when $`\mathrm{tan}\beta (v_2/v_1)`$ is large, the contribution to $`A_L^{(ij)}`$ proportional to $`f_{e_i}v_2(=\sqrt{2}m_{e_i}\mathrm{tan}\beta )`$ becomes dominant. Then, the dominant contribution to $`\tau ^\pm \mu ^\pm \gamma `$ is from the diagram of Fig. (1) for $`\mathrm{tan}\beta _{}^>1`$.
In Fig. (2) we show the branching ratio of $`\tau ^\pm \mu ^\pm \gamma `$ as a function of the Dirac neutrino mass for tau neutrino $`m_{\nu _\tau D}`$ (the right-handed tau neutrino mass $`M_{\nu _\tau }`$). Here, $`m_{\nu _\tau }=0.07`$eV, $`\mathrm{sin}2\theta _D=1`$. Also, we take $`m_{\stackrel{~}{e}_L}=170`$GeV and the wino mass 130GeV. The other gaugino masses are determined by the GUT relation for the gaugino masses for simplicity. Also, we impose the radiative breaking condition of the SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub> gauge symmetries with $`\mathrm{tan}\beta =3,10,30`$ and the Higgsino mass parameter positive. The branching ratio is proportional to $`m_{\nu _\tau D}^4`$ ($`M_{\nu _\tau }^2`$). The current experimental bound is $`Br1.1\times 10^6`$ $`^\mathrm{?}`$, and some region is excluded by it. If $`10^8`$ can be reached in the future experiments, such as B factories, we can probe $`m_{\nu _\tau D}>20(80)`$GeV for $`\mathrm{tan}\beta =30(3)`$. Then, if the Dirac tau neutrino mass is as large as the top quark mass, we may observe $`\tau ^\pm \mu ^\pm \gamma `$ there.
## 4 The solar neutrino result and $`\mu ^+e^+\gamma `$
In this section we discuss the relation between the solar neutrino result and $`\mu ^+e^+\gamma `$, assuming that the solar neutrino deficit comes from the $`\nu _e\nu _\mu `$ oscillation. The relation is more complicated compared with that between the atmospheric neutrino result and $`\tau ^\pm \mu ^\pm \gamma `$.
There are three candidates for the solution of the solar neutrino deficit if it comes from neutrino oscillation. The MSW solution $`^\mathrm{?}`$ due to the matter effect in the sun gives the natural explanation, and the observation favors
$`\mathrm{\Delta }m_{\nu _e\nu _Y}^210^{(45)}\mathrm{eV}^2\text{ or}10^7\mathrm{eV}^2,`$
$`\mathrm{sin}^22\theta _{\nu _e\nu _Y}^{}{}_{}{}^{>}0.5,`$ (23)
or
$`\mathrm{\Delta }m_{\nu _e\nu _Y}^210^5\mathrm{eV}^2,`$
$`\mathrm{sin}^22\theta _{\nu _e\nu _Y}10^{(23)}.`$ (24)
If the solar neutrino anomaly comes from so-called ’just so’ solution $`^\mathrm{?}`$, the neutrino oscillation in vacuum, the mass-squared difference and mixing angle are expected as $`^\mathrm{?}`$
$`\mathrm{\Delta }m_{\nu _e\nu _Y}^210^{(1011)}\mathrm{eV}^2,`$
$`\mathrm{sin}^22\theta _{\nu _e\nu _Y}^{}{}_{}{}^{>}0.5.`$ (25)
Assuming that the neutrino masses hierarchical as $`m_{\nu _\tau }m_{\nu _\mu }m_{\nu _e}`$, it is natural to consider $`\nu _Y=\nu _\mu `$. If one of the large angle solutions for the solar neutrino anomaly is true, the large mixing $`\theta _{\nu _\mu \nu _e}`$ may imply the LFV large mixing for sleptons between the first- and the second-generations. Similar to the atmospheric neutrino case, it is natural to consider that the large mixing angle between $`\nu _\mu `$ and $`\nu _e`$ in the MSW solution or the ’just so’ solution for the solar neutrino anomaly comes from $`V_D`$.
The amplitude for $`\mu ^+e^+\gamma `$ is proportional to $`(m_{\stackrel{~}{L}}^2)_{\mu e}`$, and it has two contributions in the SUSY SM with right-handed neutrinos as
$`(m_{\stackrel{~}{L}}^2)_{\mu e}`$ $``$ $`{\displaystyle \frac{1}{8\pi ^2}}(3m_0^2+a_0^2)\times `$ (26)
$`\left(V_{D\tau \mu }^{}V_{D\tau e}f_{\nu _\tau }^2\mathrm{log}{\displaystyle \frac{M_{\mathrm{grav}}}{M_{\nu _\tau }}}+V_{D\mu \mu }^{}V_{D\mu e}f_{\nu _\mu }^2\mathrm{log}{\displaystyle \frac{M_{\mathrm{grav}}}{M_{\nu _\mu }}}\right).`$
Here, we assume $`f_{\nu _\tau }f_{\nu _\mu }f_{\nu _e}`$, and the term proportional to $`f_{\nu _e}^2`$ is neglected. Unfortunately, we do not have information about $`V_{D\tau e}`$ and we can not evaluate the first term in Eq. (26). On the other hand, we can evaluate the second term if $`V_{D\mu e}`$ can be determined from the solar neutrino result. Then, in the following, we evaluate the event rate for $`\mu ^+e^+\gamma `$ assuming $`V_{D\tau e}=0`$. Notice that though this gives the conservative value for the event rate, there are also possibilities where the event rate is larger or smaller due to the finite $`V_{D\tau e}`$.
Let us evaluate $`\mu ^+e^+\gamma `$. The forms of the amplitude and the event rate are the same as those of $`\tau ^\pm \mu ^\pm \gamma `$ (Eqs. (21,22)). As mentioned above, if the solar neutrino anomaly comes from the MSW effect or the vacuum oscillation with the large angle, $`V_{D\mu e}`$ is expected to be large. This may lead to large $`(m_{\stackrel{~}{L}}^2)_{\mu e}`$. In Fig. (3-a), under the condition that
$$V_D=\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}& \frac{1}{2}& \frac{1}{2}\\ \frac{1}{\sqrt{2}}& \frac{1}{2}& \frac{1}{2}\\ 0& \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}\end{array}\right),$$
(27)
we show the branching ratio of $`\mu ^+e^+\gamma `$ as a function of $`m_{\nu _\mu D}`$ ($`M_{\nu _\mu }`$). We take $`m_{\nu _\mu }=4.0\times 10^3`$eV, which is consistent with the MSW solution. The other input parameters are taken to be the same as in Fig. (2). The branching ratio is promotional to $`m_{\nu _\mu D}^4`$ ($`M_{\nu _\mu }^2`$). For $`\mathrm{tan}\beta =30(3)`$, the branching ratio reaches the experimental bound ($`\mathrm{Br}(\mu e\gamma )<1.2\times 10^{11}`$ $`^\mathrm{?}`$) when $`m_{\nu _\mu D}4(10)`$GeV. A future experiment at PSI is expected to reach 10<sup>-14</sup> $`^\mathrm{?}`$. This corresponds to $`m_{\nu _\mu D}0.5(2)`$GeV. If we take $`m_{\nu _\mu }=1.0\times 10^5`$eV expected by the ’just so’ solution (Fig. (3-b)), the branching ratio becomes slightly smaller for a fixed $`m_{\nu _\mu D}`$ since the log factor in Eq. (26) is smaller.
If the solar neutrino anomaly comes from the MSW solution with the small mixing, we cannot distinguish whether the mixing comes from $`V_D`$ or $`V_M`$ even if using argument of naturalness. If it comes from $`V_D`$, the branching ratio is smaller by about 1/100 compared with that in the MSW solution with the large mixing, as shown in Fig. (3-c). In Fig. (3-c) we assume that
$$V_D=\left(\begin{array}{ccc}1& 0.04& 0.03\\ 0.04& 0.79& 0.59\\ 0& 0.60& 0.80\end{array}\right)$$
(28)
and $`m_{\nu _\mu }=2.2\times 10^3`$eV. Other input parameters are the same as Fig. (2).
Finally we consider the $`\mu ^+e^+e^{}e^+`$ process and the $`\mu `$-$`e`$ conversion in nuclei. For these processes the penguin type diagrams tend to dominate over the others, so the behavior of the decay rate is similar to that of $`\mu ^+e^+\gamma `$. For the $`\mu ^+e^+e^{}e^+`$ process the following approximate relation holds,
$`\mathrm{Br}(\mu 3e)`$ $``$ $`{\displaystyle \frac{\alpha }{8\pi }}{\displaystyle \frac{8}{3}}\left(\mathrm{log}{\displaystyle \frac{m_\mu ^2}{m_e^2}}{\displaystyle \frac{11}{4}}\right)\mathrm{Br}(\mu e\gamma )`$ (29)
$``$ $`7\times 10^3\mathrm{Br}(\mu e\gamma ).`$ (30)
For the $`\mu `$-$`e`$ conversion rate $`\mathrm{\Gamma }(\mu e)`$ a similar relation holds at the $`\mathrm{tan}\beta _{}^>1`$ region,
$$\mathrm{\Gamma }(\mu e)16\alpha ^4Z_{\mathrm{eff}}^4Z|F(q^2)|^2\mathrm{Br}(\mu e\gamma ).$$
(31)
Here $`Z`$ is the proton number in the nucleus, and $`Z_{\mathrm{eff}}`$ is the effective charge, $`F(q^2)`$ the nuclear form factor at the momentum transfer $`q`$. The $`\mu `$-$`e`$ conversion rate normalized by the muon capture rate in Ti nucleus, $`R(\mu ^{}e^{};{}_{22}{}^{48}\mathrm{Ti})`$, is approximately
$$R(\mu ^{}e^{};{}_{22}{}^{48}\mathrm{Ti})6\times 10^3\mathrm{Br}(\mu e\gamma ).$$
(32)
The MECO collaboration proves that they have a technology to reach $`R(\mu ^{}e^{};{}_{22}{}^{48}\mathrm{Ti})<10^{16}`$ $`^\mathrm{?}`$. Furthermore, now there are active discussions of the high intensity muon source, and we may reach to a level of $`10^{18}`$ if the muon storage is constructed and $`10^{(1920)}`$ muons per a year are produced$`^\mathrm{?}`$. This is comparable to $`\mathrm{Br}(\mu e\gamma )10^{16}`$, and we can probe the region $`m_{\nu _\mu D}0.2(0.5)`$GeV ($`M_{\nu _\mu }10^{10}`$ $`(10^{11})`$GeV) in the MSW solution with the large angle.
## 5 Conclusion
In this article we discuss the charged lepton-flavor violating processes, $`\mu ^+e^+\gamma `$ and $`\tau ^\pm \mu ^\pm \gamma `$, using the current neutrino experimental data, in the SUSY SM with the right-handed neutrinos. While this model has many unknown parameters, these processes may be accessible in near future experiment. The LFV search will give new insights to the origin of the neutrino masses.
## References
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# Foundations of self-consistent particle-rotor models and of self-consistent cranking models
## I Introduction
The aims of this paper are to study anew the foundations of the particle-rotor model and of the cranking model . The particle-rotor model (PRM) was introduced as an angular momentum-conserving phenomenological description of odd deformed nuclei. Because of its relative ease of application and, on the whole, quite remarkable success, it has been applied even up to the present (for instance ), with various alterations of detail, to a myriad of applications, over a lifetime of more than four and a half decades. Among the extensions, we mention in particular that to the description of triaxial nuclei , the original model having been formulated for axially symmetric nuclei.
In one of the textbooks , p. 109, we find, after a glowing appraisal of the success of the model, the following statement: “However, until now a clear-cut microscopic derivation has been missing.” In fact, a microscopic derivation had been given earlier , based on the Kerman-Klein (KK) method . The microscopic foundation of the axially symmetric PRM was studied more recently in , starting from a semi-microscopic version of the KK approach, and compared in accuracy, for several examples of well-deformed nuclei, both with its more accurate progenitor and with the inherently less accurate cranking approximation.
The cranking model was originally introduced into nuclear physics , within the framework of a prescribed single-particle model, to deal with the enigma presented by the first values encountered for the moments of inertia of deformed nuclei. An extended version , the one considered in most applications until recent years, was based on the self-consistent mean-field theory of a deformed rotating object. This early work was designed primarily to provide formulas for the moment of inertia.
The full range of applicability of the self-consistent cranking model, as well as its limitations, was realized in the so-called cranked shell model (CSM) , that has been widely applied to the analysis of band-crossing and other high-spin phenomena. (For a current list of references, especially reviews, see .) The formulations under discussion, which apply to axially symmetric nuclei, assume that collective rotation occurs about a principal axis perpendicular to the symmetry axis. Such a formulation is referred to currently as principal-axis cranking (PAC) as opposed to a recent generalization, called tilted-axis cranking (TAC) . In the latter, even in the axial case, the system may rotate about an axis in a principal plane of the assumed (quadrupole) intrinsic shape, and for the triaxial case about an arbitrary (dynamically determined) direction with respect to the principal axes.
A second aim of the present paper is to establish the relationship of the cranking models, including the recent generalized versions, to a microscopic theory. The previous literature on this subject is modest in extent. The standard references are , the major results of which are reproduced and discussed in . Briefly, starting from a formulation of the microscopic theory by means of generator coordinates, the energy is evaluated approximately as a power series in the angular momentum by a method due to Kamlah , valid for large deformations. When the variational method is applied to the lowest non-trivial approximation of this procedure, it can be shown that the cranking theory is a solution of the resulting equations. This is summarized by stating that cranking is a solution, involving a semi-classical approximation, of the method of variation after projection as opposed to the exact procedure of variation before projection.
To our knowledge, the only other studies of this subject are those based on the KK method, a brief treatment of the case of rotation in a plane that predates the above-cited work and two studies that postdated them, one again on the problem of rotation in a plane and the second a restricted study of the triaxial case . (Some discussion of the cranking limit, also based on a variant of the KK method, can be found in .) Up to now we have never presented a full account of the three-dimensional treatment either for axial or for triaxial nuclei. Since our methods contain features distinct from those found in the standard literature , and in view of the renewed interest in generalized cranking models , the publication of a detailed account is perhaps justified, even at this late date.
The foundations of the study are presented in Sec. II. We utilize a shell model Hamiltonian, widely employed for medium and heavy nuclei, with two-particle interactions in which the latter are separated into two parts clearly distinguished as multipole and pairing forces, respectively. The advantage of such a model is that the (c-number) equations of motion that can be derived from it by the KK method are completely rigorous. It is a simplifying feature for the further study to recognize that these equations can be derived from a variational principle that we called the trace variational principle, suggested in our earliest paper and developed more fully in and in . This variational principle has several noteworthy features: (i) It is formulated for the many-body problem in the language of second quantization. (ii) The quantities varied are not wave functions, but rather a suitably chosen set of matrix elements, in our case coefficients of fractional parentage (to be discussed at the appropriate point of Sec. II). (iii) Rather than involving the Rayleigh-Ritz principle for one state at a time, the functional to be varied is the trace of energy expectation values over a prescribed space of states.
\[It turns out that not all aspects of our formulation are novel. Thus an incomplete version of the trace variational principle is to be found in one of the initial series of papers on matrix mechanics , in which the variational parameters are matrix elements of the coordinates and momenta. This application to particle quantum mechanics was discovered and developed independently by us in several accounts of which the most recent is . A version of the trace variational principle can, furthermore, be found is in a classic text in mathematical physics . Here the formulation is close to standard Rayleigh-Ritz, in that the quantities varied are wave function. This formulation has found its way into the theory of density functionals and even been generalized to include the case where the trace is replaced by a different weighted diagonal sum . Most recently the trace variational principle for fields has appeared in a quaternion generalization of quantum mechanics .\]
The theory is elaborated in Sec. II only as far as is required for the remaining body of the text. Further development is presented in Appendix A, with an eye to formulating algorithms that can eventually be applied to the study of the self-consistent problem posed by the formulation in Sec. II. We turn to applications in Sec. III, where we derive the self-consistent PRM from the variational principle associated with the KK equations. (With one possible exception , we are unaware of any recent work, other than our own, that has examined the foundations of the PRM.)
The formalism presented in Sec. III does not lend itself naturally to a derivation of the self-consistent cranking theory, which should be a limit of the self-consistent PRM. In Sec. IV we describe an alternative derivation of the PRM, following ideas first advanced briefly in , that does lead directly to the cranking limit. The considerations of Secs. III and IV apply to axially symmetric nuclei. Both treatments are extended to the case of triaxial nuclei in Sec. V. Further discussion of results and conclusions are given in Sec. VI.
## II Equations of motion and variational principle
We choose a shell-model Hamiltonian in the form
$`H`$ $`=`$ $`h_aa_\alpha ^{}a_\alpha +{\displaystyle \frac{1}{2}}F_{\alpha \gamma \delta \beta }a_\alpha ^{}a_\gamma a_\beta ^{}a_\delta +{\displaystyle \frac{1}{2}}G_{\alpha \gamma \beta \delta }a_\alpha ^{}a_\gamma ^{}a_\delta a_\beta .`$ (1)
In this standard model, the $`a_\alpha ,a_\alpha ^{}`$ are the destruction, creation operators for fermions in the shell-model mode $`\alpha =(nljm\tau )`$ ($`\tau `$ distinguishing neutrons from protons); $`F_{\alpha \gamma \delta \beta }`$ describes multipole forces and $`G_{\alpha \gamma \delta \beta }`$ pairing forces. In this version, all multipolarities allowed by angular momentum conservation are included, though in practice we limit ourselves to the lowest few multipoles of each type. We shall also consistently use the summation convention, except when we wish to highlight some set of indices. With the help of the definitions
$`F_{\alpha \gamma \delta \beta }`$ $`=`$ $`s_\gamma (j_am_aj_cm_c|LM_L)s_\beta (j_dm_dj_bm_b|LM_L)F_{acdb}(L),`$ (2)
$`G_{\alpha \gamma \delta \beta }`$ $`=`$ $`(j_am_aj_cm_c|LM_L)(j_dm_dj_bm_b|LM_L)G_{acdb}(L),`$ (3)
$`s_\gamma `$ $`=`$ $`(1)^{j_cm_c}=\sqrt{2j_c+1}(j_cm_cj_cm_c|00),`$ (4)
where $`(jmj^{}m^{}|LM)`$ is a Clebsch-Gordon (CG) coefficient, the operator equations of motion can be obtained in the form
$`[a_\alpha ,H]`$ $`=`$ $`h_a^{}a_\alpha +F_{\alpha \alpha ^{}\beta ^{}\beta }a_\alpha ^{}a_\beta ^{}a_\beta ^{}+G_{\alpha \alpha ^{}\beta \beta ^{}}a_\alpha ^{}^{}a_\beta ^{}a_\beta ,`$ (5)
$`h_a^{}`$ $`=`$ $`h_a{\displaystyle \frac{1}{2}}F_{abab}{\displaystyle \frac{2L+1}{2j_a+1}},`$ (6)
$`[a_{\overline{\alpha }}^{},H]`$ $`=`$ $`h_a^{\prime \prime }a_{\overline{\alpha }}^{}F_{\beta \beta ^{}\alpha ^{}\overline{\alpha }}a_\beta ^{}^{}a_\beta a_\alpha ^{}^{}G_{\beta \overline{\beta }^{}\alpha ^{}\overline{\alpha }}a_\alpha ^{}a_\beta ^{}a_{\overline{\beta }^{}}^{},`$ (7)
$`h_a^{\prime \prime }`$ $`=`$ $`h_a^{}+2{\displaystyle \frac{2L+1}{2j_a+1}}G_{abab}(L).`$ (8)
Here, for example, $`\overline{\alpha }=(j_a,m_a)`$.
To develop a dynamical scheme, we turn to the problem of obtaining equations for the matrix elements of Eqs. (5) and (7). We designate a state of interest of an odd nucleus as $`|JM\nu `$, where $`J`$ is the total angular momentum, $`M`$ is its $`z`$ component, and $`\nu `$ are the remaining quantum numbers necessary for unique specification of the state. Neighboring even nuclei are specified, correspondingly, as $`|\overline{IMn}`$, referring to a heavier neighbor, and $`|\underset{¯}{IMn}`$, referring to a lighter neighbor. Below we shall then derive equations for the matrix elements, referred to as CFP (coefficients of fractional parentage)
$`JM\nu |a_\alpha |\overline{IM_In}`$ $`=`$ $`V_{JM\nu }(\alpha IM_In),`$ (9)
$`JM\nu |a_{\overline{\alpha }}^{}|\underset{¯}{IM_In}`$ $`=`$ $`U_{JM\nu }(\alpha IM_In).`$ (10)
We shall require the full notation when we turn to applications in the next section. For the formal developments of this section, we utilize a compressed notation, with
$$JM\nu i,IM_Inn.$$
(11)
With new symbols defined and discussed below, we thus obtain the equations
$`_iV_i(\alpha n)`$ $`=`$ $`(ϵ_a^{}E_{\overline{n}}^{})V_i(\alpha n)+F_{\alpha \alpha ^{}\beta ^{}\beta }[V_i^{}^{}(\beta n^{})V_i^{}(\beta ^{}n)]V_i(\alpha ^{}n^{})`$ (13)
$`+G_{\alpha \overline{\alpha }^{}\beta \overline{\beta }^{}}[U_i^{}^{}(\beta ^{}n^{})V_i^{}(\beta n)]U_i(\alpha ^{}n^{}),`$
$`_iU_i(\alpha n)`$ $`=`$ $`(ϵ_a^{\prime \prime }E_{\underset{¯}{n}}^{})U_i(\alpha n)F_{\overline{\beta }\overline{\beta }^{}\overline{\alpha }^{}\overline{\alpha }}[U_i^{}^{}(\beta n^{})U_i^{}(\alpha ^{}n)]U_i(\beta ^{}n^{})`$ (15)
$`+G_{\overline{\alpha }\alpha ^{}\overline{\beta }^{}\beta }[V_i^{}^{}(\beta n^{})U_i^{}(\beta ^{}n)]V_i(\alpha ^{}n^{}).`$
In the definitions that follow, we understand that $`E_i`$ is the energy of the state $`|i`$ and that $`E_{\overline{n}}`$ and $`E_{\underset{¯}{n}}`$ are, correspondingly the energies of the neighboring even states, with the subscript $`0`$ standing either for the ground state, or for the lowest energy state considered, which for conciseness we shall continue to refer to as the ground state. We thus encounter the quantities
$`_i`$ $`=`$ $`E_i+{\displaystyle \frac{1}{2}}(E_{\overline{0}}+E_{\underset{¯}{0}}),`$ (16)
$`ϵ_a^{}`$ $`=`$ $`h_a^{}\lambda ,`$ (17)
$`\lambda `$ $`=`$ $`{\displaystyle \frac{1}{2}}(E_{\overline{0}}E_{\underset{¯}{0}}),`$ (18)
$`E_n^{}`$ $`=`$ $`E_nE_0.`$ (19)
The physical significance of the quantities defined in Eqs. (16)–(19) is evident. $`_i`$ are the negatives of the energies of the odd nucleus relative to the ground-state energies of its even neighbors, $`ϵ_a`$, variously primed, are single-particle energies measured relative to the chemical potential $`\lambda `$, and $`E_n^{}`$ are excitation energies of the appropriate even nuclei. Finally in achieving the form of Eq. (15), we have assumed that $`F`$ and $`G`$ are real. Given the Hamiltonian (1), Eqs. (13) and (15) are an exact set of consequences that define a non-linear eigenvalue problem with eigenvalue $`_i`$. The elements on the right hand sides of these equations define an effective Hamiltonian that will be discussed in considerable further detail in the course of this work.
We display next a functional, $``$, whose vanishing first variations yield the equations of motion, namely,
$``$ $`=`$ $`ϵ_a^{}|V_i(\alpha n)|^2ϵ_a^{\prime \prime }|U_i(\alpha n)|^2`$ (26)
$`{\displaystyle \frac{1}{2}}F_{\alpha \alpha ^{}\beta ^{}\beta }[V_i^{}^{}(\beta n^{})V_i^{}(\beta ^{}n)][V_i^{}(\alpha n)V_i(\alpha ^{}n^{})]`$
$`+G_{\alpha \overline{\alpha }^{}\beta ^{}\overline{\beta }}[U_i^{}^{}(\beta n^{})V_i^{}(\beta ^{}n)][V_i^{}(\alpha n)U_i(\alpha ^{}n^{})]`$
$`{\displaystyle \frac{1}{2}}F_{\overline{\beta }\overline{\alpha }^{}\overline{\beta }^{}\overline{\alpha }}[U_i^{}^{}(\beta n^{})U_i^{}(\beta ^{}n)][U_i^{}(\alpha n)U_i(\alpha ^{}n^{})`$
$`_i[|V_i(\alpha n)|^2+U_i(\alpha n)|^2]`$
$`E_{\overline{n}}^{}|V_i(\alpha n)|^2E_{\underset{¯}{n}}^{}|U_i(\alpha n)|^2`$
$`𝒢_i[|V_i(\alpha n)|^2+|U_i(\alpha n)|^2].`$
One verifies that the equations of motion (13) and (15) emerge, respectively, from the requirements
$$\frac{\delta }{\delta V_i^{}(\alpha n)}=\frac{\delta }{\delta U_i^{}(\alpha n)}=0.$$
(27)
It is natural to inquire at this point if the functional $``$ has any simple physical significance, in particular, if it is related to a Rayleigh-Ritz principle. To answer this question, we evaluate the sum
$$\mathrm{Tr}(\overline{\mathrm{H}}+\underset{¯}{\mathrm{H}})=\underset{\mathrm{n}}{}[\overline{\mathrm{n}}|\mathrm{H}|\overline{\mathrm{n}}+\underset{¯}{\mathrm{n}}|\mathrm{H}|\underset{¯}{\mathrm{n}}].$$
(28)
The evaluation of this sum with the aim of eventually recognizing the relevant pieces of $``$ requires, in addition to the standard tool of completeness, some algebraic rearrangement of the trace involving the lighter system, just as was necessary in the equations of motion. We then find that the interaction terms match exactly those in Eq. (26), but that the single particle terms do not. Instead we find
$`h_a^{}`$ $``$ $`h_a\overline{h}_a,`$ (29)
$`h_a^{\prime \prime \prime }`$ $``$ $`h_a+2{\displaystyle \frac{2L+1}{2j_a+1}}G_{acac}(L)+\sqrt{{\displaystyle \frac{2j_b+1}{2j_a+1}}}F_{aabb}(0)\underset{¯}{h}_a.`$ (30)
We are thus tempted to replace the functional $``$, as basis for the theory, by a functional that contains the new single-particle energies. We do not make this change because it destroys the simple physical significance of the Lagrange multiplier terms in Eq. (26) to which we next turn our attention. In practice, these extra single-particle terms are often ignored anyway.
We consider then the Lagrange-multiplier terms that appear in Eq. (26). The relevant question concerns the constraints that have been imposed on the variations. Since
$$\underset{i\alpha }{}|V_i(\alpha n)|^2=\underset{\alpha }{}\overline{n}|a_\alpha ^{}a_\alpha |\overline{n}=\overline{n}|\widehat{N}|\overline{n},$$
(31)
where $`\widehat{N}`$ is the number operator, we see that the excitation energies $`E_{\overline{n}}^{}`$ enter as Lagrange multipliers for the conservation of nucleons in the heavier even nucleus. Similarly the term involving the sum over the $`|U_i(\alpha n)|^2`$ expresses (to an additive constant) the conservation of nucleons in the lighter system. Finally, we show that the eigenvalue $`_i`$ is (no surprise here) a Lagrange multiplier for an appropriate normalization condition. To see this we take the matrix element in the state $`|i`$ of the summed anticommutator,
$$\underset{\alpha }{}\{a_\alpha ,a_\alpha ^{}\}=\mathrm{\Omega },\mathrm{\Omega }=\underset{j_a}{}(2j_a+1)=\underset{a}{}\mathrm{\Omega }_a,$$
(32)
and thus find
$$\frac{1}{\mathrm{\Omega }}\underset{\alpha n}{}[|V_i(\alpha n)|^2+|U_i(\alpha n)|^2]=1.$$
(33)
Orthogonality constraints on the solutions need not be imposed, since they follow directly from the equations of motion.
There is more to the story, however. We must note that Eq. (33) is only a sum of required normalization conditions. From the summed anticommutator for each level,
$$\underset{m_a}{}\{a_\alpha ,a_\alpha ^{}\}=\mathrm{\Omega }_a,,$$
(34)
we have
$$\frac{1}{\mathrm{\Omega }_a}\underset{m_an}{}[|V_i(\alpha n)|^2+U_i(\alpha n)|^2]=1.$$
(35)
If Eqs. (13) and (15) described a linear eigenvalue problem, it would be impossible to impose the additional normalization conditions represented by Eq. (35). For the general non-linear problem, there is no a priori inconsistency; the satisfaction of these conditions will be a part of any fully satisfactory algorithm. The form of the normalization condition (35) suggests, furthermore, that it may be both useful and natural to rescale the CFP,
$`V_i(\alpha n)`$ $`=`$ $`\sqrt{2j_a+1}v_i(\alpha n),`$ (36)
$`U_i(\alpha n)`$ $`=`$ $`\sqrt{2j_a+1}u_i(\alpha n).`$ (37)
There is considerably more to the formal theory than what has been presented thus far. However, we have all the tools needed for the further development in the text and thus relegate the additional theoretical considerations to Appendix A.
## III Derivation of particle-rotor model and its cranking limit: axially symmetric case
As a first illustration of the formalism presented in the previous section, we assume that the even (core) nuclei are in a single axially-symmetric band $`|IM_IK`$, where $`K`$ is the component of the angular momentum along the figure axis. There are at least two cases where it makes some physical sense to isolate a single $`K`$ value, where it is the ground-state band with $`K=0`$, or where the band has a large $`K`$-value and we are dealing with an isomeric state.
We first use rotational invariance to study the structure of the amplitudes $`V`$ and $`U`$ defined in Eqs. (9) and (10), respectively. For this purpose we introduce a complete set of states $`|R`$ localized in the Euler angles, $`R=(\alpha \beta \gamma )`$, where $`\alpha ,\beta `$ are the usual polar and azimuthal angles, respectively, and write
$`|IM_IK`$ $`=`$ $`{\displaystyle 𝑑R|RR|IM_IK}`$ (38)
$`=`$ $`\left({\displaystyle \frac{2I+1}{8\pi ^2}}\right)^{\frac{1}{2}}{\displaystyle 𝑑R|RD_{M_IK}^{(I)}(R)}.`$ (39)
The identification of a scalar product of many-body states with the Wigner $`D`$ function is not a trivial statement, but is rather an essential element in the definition of the model to be studied. In fact, the designation $`|R`$ for the many-body state is insufficiently detailed and is made more explicit by the statement
$$|R=U(R)|\widehat{0}K,$$
(40)
where $`|\widehat{0}K`$ is an axially symmetric intrinsic state spinning with angular momentum $`K`$ about its symmetry axis, and $`U(R)`$ is the unitary rotation operator in the many-body space defined by the Euler angles that specify the rotation $`R`$. For such a state, we thus note the relation, with $`(\alpha \beta \gamma )=(\widehat{n}\gamma )`$
$$U(\widehat{n}\gamma )|\widehat{0}K\mathrm{exp}(iK\gamma )=U(\widehat{n}0)|0K.$$
(41)
The introduction of strictly localized states is, of course, an idealization that ignores the reality of band termination, but it is a standard approximation for well-deformed nuclei.
\[The previous discussion and that which follows does not take into account $``$ invariance, the invariance of the quadrupole shape under a rotation of $`\pi `$ about a principal axis. To include this symmetry in the discussion, we replace the state $`|IM_IK`$ by an eigenfunction of $``$,
$$|IM_IK]=\frac{1}{2}\{|IM_IK+(1)^{I+K}|IM_IK\}.$$
(42)
We then imitate the arguments starting on p. 8 of . The task is to sort and collect the extra terms that appear both in the equations of motion and in formulas for one and two-particle observables. \]
When Eq. (39) is substituted into the definition (9) of $`V`$, and use is made of the definitions to be given below, we are thereby led to the study of an amplitude such as
$`JM\nu |a_\alpha |R`$ $`=`$ $`JM\nu |U(R)U^1(R)a_\alpha U(R)|\widehat{0}K`$ (43)
$`=`$ $`{\displaystyle \underset{M^{}}{}}JM\nu |U|JM^{}\nu JM^{}\nu |U^1a_\alpha U|\widehat{0}K`$ (44)
$`=`$ $`{\displaystyle \underset{M^{}\kappa _a}{}}D_{MM^{}}^{(J)}(R)D_{m_a\kappa _a}^{(j_a)}(R)\chi _{JM^{}\nu }(j_a\kappa _a,K)(1)^{j_a+\kappa _a},`$ (45)
where the previous manipulations have utilized the following relations and definitions (of which the first two are standard):
$`JM|U(R)|JM^{}`$ $`=`$ $`D_{MM^{}}^{(J)}(R),`$ (46)
$`U^1(R)a_{jm}U(R)`$ $`=`$ $`{\displaystyle \underset{\kappa }{}}a_{j\kappa }D_{m\kappa }^{(j)}(R),`$ (47)
$`JM\nu |a_{jm}|\widehat{0}`$ $``$ $`(1)^{j+m}\chi _{JM\nu }(jm,K)`$ (48)
The phase in (48) has been introduced for algebraic convenience.
With the help of the integral of a product of three $`D`$ functions (, Eq. (4.6.2)) and the application of standard symmetry properties of CG coefficients (Eqs. (3.5.15) and (3.5.16) of the same reference), we find for the CFP defined in Eq. (9),
$`V_{JM\nu }(\alpha IM_IK)`$ $`=`$ $`{\displaystyle \underset{\kappa _a}{}}\sqrt{{\displaystyle \frac{8\pi ^2}{2j_a+1}}}(1)^{JM}`$ (51)
$`\times (IM_IJM|j_am_a)(JK\kappa _aj_a\kappa _a|IK)`$
$`\times (1)^{j_a+\kappa _a}\chi _{JK\kappa _a\nu }(j_a\kappa _a,K).`$
A similar analysis carried out for the amplitude $`U`$ defined in (10) yields the result
$`U_{JM\nu }(\alpha IM_IK)`$ $`=`$ $`{\displaystyle \underset{\kappa _a}{}}\sqrt{{\displaystyle \frac{8\pi ^2}{2j_a+1}}}`$ (54)
$`\times (1)^{JM+j_a\kappa _a+j_a+m_a}(IM_IJM|j_am_a)`$
$`\times (JK\kappa _aj_a\kappa _a|IK)\varphi _{JK\kappa _a\nu }(j_a\kappa _a,K),`$
$`\varphi _{JM\nu }(j_a\kappa _a,K)`$ $`=`$ $`JM\nu |a_{j_a\kappa _a}^{}|\widehat{0}K.`$ (55)
It is most succinct to base further discussion on the variational principle (26). We evaluate this expression when the core collective states are restricted to the members of a single band of an axial rotor, and the states of the odd nucleus are any states that can arise from the coupling. Returning to a full nomenclature, this calls for the identifications
$`\overline{n}\overline{IM_IK},\underset{¯}{n}\underset{¯}{IM_IK},`$ (56)
$`iJM\nu .`$ (57)
We are assuming here that there are corresponding bands in the two even nuclei that couple to the given odd nucleus. In the following we shall also suppress the bar and underline in the CFP, understanding them from context, but continue to emphasize this distinction in the excitation energies $`E^{}`$.
We consider first the evaluation of all terms in the variational principle depending on $`|V_i(\alpha n)|^2`$, which includes some of the single-particle terms and some of the Lagrange-multiplier terms. Writing $`\overline{E}^{}(IK)`$ for $`E_{\overline{n}}^{}`$, we consider in particular the combination
$$(_{J\nu }+\overline{E}^{}(IK))|V_{JM\nu }(\alpha IM_IK)|^2.$$
(58)
For the evaluation of the CFP $`V`$, we utilize Eq. (51), renormalized, however, by a factor of $`(1/\sqrt{8\pi ^2})`$ in order that the reciprocal of such factors do not appear in the answer, i. e., we rescale $`\chi `$ by this factor. For the evaluation of expressions involving the CFP $`U`$, we shall utilize Eq. (54) similarly renormalized. This rescaling will be understood throughout the remainder of this paper.
The rest of this section consists of a relatively detailed account of the evaluation of the variational sum in the “intrinsic” system. Subsequent variation will led to the self-consistent version of the strong coupling PRM. Toward this end, as part of the definition of an axial rotor, we assume that, equally for the barred and underlined quantities,
$`E^{}(IK)`$ $``$ $`E^{}(\stackrel{}{I}^2K^2)=E^{}(\widehat{I}_1^2+\widehat{I}_2^2+\widehat{I}_3^2K^2),`$ (59)
where we have introduced intrinsic components of the angular momentum. The arrow indicates the replacement of an eigenvalue by an operator. This is done by making use of the appropriate one of the CG coefficients, understood as a scalar product, that appear in Eq. (9), as follows
$$E^{}(IK)(JK\kappa _aj_a\kappa _a|IK)=(JK\kappa _aj_a\kappa _a|E^{}(\stackrel{}{I}^2K^2)|IK).$$
(60)
Furthermore, in each term of the sum $`I`$ is coupled with some $`j_a`$ to a value of $`J`$. From the structure of the CFP, it follows that we may replace $`\stackrel{}{I}`$ by $`\stackrel{}{J}+\stackrel{}{j}_a`$ in Eq. (60) and write (with $`j_aj)`$,
$`\overline{E^{}}(\stackrel{}{I}^2K^2)`$ $`=`$ $`\overline{E^{}}[(\stackrel{}{J}+\stackrel{}{j})^2K^2]`$ (61)
$`=`$ $`\overline{E^{}}(\stackrel{}{J}^2K^2)+{\displaystyle \frac{\overline{E^{}}}{\widehat{J}_i}}\widehat{j}_i+{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2\overline{E}^{}}{\widehat{J}_i^2}}+\mathrm{}.`$ (62)
It is not necessary for these considerations that $`E(IK)`$ have the simple form of a rotor spectrum, only that it be a function as indicated. The first term of Eq. (62) may be replaced by an eigenvalue $`\overline{E}^{}[J(J+1)K^2)]`$, and the second term, the Coriolis coupling will be evaluated below. The third and possible higher order terms will not be studied here, but will be included in applications. Other than the Coriolis coupling and higher-order terms, the contribution from Eq. (58) and of the remaining single-particle terms takes the form
$`{\displaystyle \underset{J\nu j_a\kappa _a}{}}(ϵ_a^{}\overline{\epsilon }_{J\nu })|\chi _{JK\kappa _a\nu }(j_a\kappa _a,K)|^2,`$ (63)
$`\overline{\epsilon }_{J\nu }=_{J\nu }\overline{E}^{}[J(J+1)K^2].`$ (64)
For the further evaluation we can write in full generality
$$\frac{E^{}}{J_i}f(\stackrel{}{J}^2K^2)J_i,$$
(65)
where in the simplest case $`f`$ is just the reciprocal of the moment of inertia. Using the matrix elements of the raising and lowering operators, and collecting terms, we find, in connection with the second term of Eq. (62), the Coriolis coupling term, the following sum to evaluate
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{M_Im_aMIJ\kappa _a\kappa _a^{}}{}}\chi _{JK\kappa _a\nu }(j_a\kappa _a,K)\chi _{JK\kappa _a^{}\nu }^{}(j_a\kappa _a^{},K)(2j_a+1)^1`$ (66)
$`\times (1)^{2j_a+\kappa _a+\kappa _a^{}}|(IM_IJM|j_am_a)|^2(JK\kappa _a^{}j_a\kappa _a^{}|IK)`$ (67)
$`\times \overline{f}(JK)[\sqrt{(JK+\kappa _a)(J+K\kappa _a+1)}`$ (68)
$`\sqrt{(j_a+\kappa _a)(j_a\kappa _a+1)}(JK\kappa _a+1j_a\kappa _a1|IK)`$ (69)
$`+\sqrt{(J+K\kappa _a)(JK+\kappa _a+1)}\sqrt{(j_a\kappa _a)(j_a+\kappa _a+1)}`$ (70)
$`\times (JK\kappa _a1j_a\kappa _a+1|IK)+2(K\kappa _a)\kappa _a(JK\kappa _aj_a\kappa _a|IK)].`$ (71)
With the help of the standard orthonormalization conditions for CG coefficients this reduces to the final result for the term under study,
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{J\nu j_a\kappa _a}{}}\chi _{JK\kappa _a\nu }^{}(j_a\kappa _a,K)\overline{f}(JK)`$ (72)
$`\times [\sqrt{(J+K\kappa _a)(JK+\kappa _a+1)}\sqrt{(j_a\kappa _a)(j_a+\kappa _a+1)}`$ (73)
$`\times \chi _{JK\kappa _a1\nu }(j_a\kappa _a+1,K)`$ (74)
$`+\sqrt{(JK+\kappa _a)(J+K\kappa _a+1)}\sqrt{(j_a+\kappa _a)(j_a\kappa _a+1)}`$ (75)
$`\times \chi _{JK\kappa _a+1\nu }(j_a\kappa _a1,K)`$ (76)
$`+2(K\kappa _a)\kappa _a\chi _{JK\kappa _aj_a\kappa _a}(j_a\kappa _a,K)].`$ (77)
We outline briefly the corresponding calculation of the single-particle terms associated with the $`U`$ coefficients. Only the following change is necessary: All barred energies associated with the heavier of the two neighboring even nuclei are replaced by underlined energies associated with the lighter of the two neighbors. For the terms corresponding to Eq. 63), we find
$$\underset{J\nu j_a\kappa _a}{}(ϵ_a^{\prime \prime }+\underset{¯}{\epsilon }_{J\nu })|\varphi _{JK\kappa _a\nu }(j_a\kappa _a,K)|^2.$$
(78)
For the Coriolis coupling term, we find
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{J\nu j_a\kappa _a}{}}\varphi _{JK\kappa _a\nu }^{}(j_a\kappa _a,K)\underset{¯}{f}(JK)`$ (79)
$`\times [\sqrt{(J+K\kappa _a)(JK+\kappa _a+1)}\sqrt{(j_a\kappa _a)(j_a+\kappa _a+1)}`$ (80)
$`\times \varphi _{JK\kappa _a1\nu }(j_a\kappa _a+1,K)`$ (81)
$`+\sqrt{(JK+\kappa _a)(J+K\kappa _a+1)}\sqrt{(j_a+\kappa _a)(j_a\kappa _a+1)}`$ (82)
$`\times \varphi _{JK\kappa _a+1\nu }(j_a\kappa _a1,K)+(K\kappa _a)\kappa _a\varphi _{JK\kappa _a\nu }(j_a\kappa _a,K)].`$ (83)
We turn to the contribution of the interaction terms. We evaluate a typical quartic contribution to Eq. (27), for example
$`{\displaystyle \underset{\chi }{}}{\displaystyle \frac{1}{2}}F_{\alpha \alpha ^{}\beta ^{}\beta }V_{JM\nu }(\beta ^{}IM_IK)V_{JM\nu }^{}(\beta I^{}M_I^{}K)`$ (84)
$`\times V_{J^{}M^{}\nu ^{}}(\alpha ^{}I^{}M_I^{}K)V_{J^{}M^{}\nu ^{}}^{}(\alpha IM_IK)`$ (85)
$`=F_{aa^{}b^{}b}(L)(j_am_aj_a^{}m_a^{}|Lm_am_a^{})(j_b^{}m_b^{}j_bm_b|Lm_b^{}m_b)`$ (86)
$`\times (1)^{j_a^{}m_a^{}+j_bm_b}(1)^{j_a+\kappa _a+j_a^{}+\kappa _a^{}+j_b+\kappa _b+j_b^{}+\kappa _b^{}}`$ (87)
$`\times [(2j_a+1)(2j_a^{}+1)(2j_b+1)(2j_b^{}+1)]^{\frac{1}{2}}`$ (88)
$`\times \chi _{JK\kappa _b^{}\nu }(j_b^{}\kappa _b^{})\chi _{JK\kappa _b\nu }^{}(j_b\kappa _b)\chi _{J^{}K\kappa _a^{}\nu ^{}}(j_a^{}\kappa _a^{})\chi _{J^{}K\kappa _a\nu ^{}}^{}(j_a\kappa _a)`$ (89)
$`\times (IM_IJm_b^{}M_I|j_b^{}m_b^{})(I^{}M_I^{}Jm_bM_I^{}|j_bm_b)`$ (90)
$`\times (I^{}M_I^{}J^{}m_a^{}M_I^{}|j_a^{}m_a^{})(IM_IJ^{}m_aM_I|j_am_a)`$ (91)
$`\times (JK\kappa _b^{}j_b^{}\kappa _b^{}|IK)(JK\kappa _bj_b\kappa _b|I^{}K)`$ (92)
$`\times (J^{}K\kappa _a^{}j_a^{}\kappa _a^{}|I^{}K)(J^{}K\kappa _aj_a\kappa _a|IK).`$ (93)
To obtain this form from the corresponding term of Eq. (26) we have utilized Eqs. (51) and (2). From angular momentum conservation, we have the relations
$$m_b^{}M_I=m_bM_I^{},m_a^{}M_I^{}=m_aM_I.$$
(94)
The next step is to perform the sums over the magnetic quantum numbers $`m_a`$ and $`m_b`$. For this purpose, we use Eq. (6.6.27) of with assists from Eqs. (3.5.14)-(3.5.17) of the same reference. At this point we obtain the sum (henceforth we simplify the arguments of the single-particle amplitudes, $`j_a\kappa _a,Ka`$, etc.)
$`{\displaystyle \underset{\chi }{}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}F_{aa^{}b^{}b}(1)^{j_a+\kappa _a+\kappa _a^{}+j_b^{}+\kappa _b^{}+\kappa _b+J+J^{}}(2L+1)`$ (103)
$`\times \chi _{JK\kappa b^{}\nu }(b^{})\chi _{jK\kappa _b\nu }^{}(b)\chi _{J^{}K\kappa _a^{}\nu ^{}}(a^{})\chi _{J^{}K\kappa _a\nu ^{}}^{}(a)`$
$`\times (JK\kappa _b^{}j_b^{}\kappa _b^{}|IK)(JK\kappa _bj_b\kappa _b|I^{}K)`$
$`\times (J^{}K\kappa _a^{}j_a^{}\kappa _a^{}|I^{}K)(J^{}K\kappa _aj_a\kappa _a|IK)`$
$`\times \left\{\begin{array}{ccc}j_a^{}& j_a& L\\ I& I^{}& J^{}\end{array}\right\}\left\{\begin{array}{ccc}j_b& j_b^{}& L\\ I& I^{}& J\end{array}\right\}.`$
The final summations that we expect to be able to do in general are over $`I`$ and $`I^{}`$. Toward this end we need to re express the $`6j`$ symbols that occur in Eq. (103) in terms of $`6j`$ symbols that depend separately on $`I`$ and $`I^{}`$. This can be done by the application of Eq. (6.2.12) of . We quote the value of $`_\chi `$ that results from this transformation:
$`{\displaystyle \underset{\chi }{}}={\displaystyle \frac{1}{2}}{\displaystyle \underset{II^{}\mathrm{}}{}}F_{aa^{}b^{}b}(L)(1)^{\kappa _a+\kappa _b+\kappa _a^{}+\kappa _b^{}+j_a^{}+j_b+L+I+I^{}+I^{\prime \prime }}`$ (104)
$`\times (2L+1)(2I^{\prime \prime }+1)\left\{\begin{array}{ccc}j_a^{}& j_a& L\\ j_b^{}& j_b& I^{\prime \prime }\end{array}\right\}\left\{\begin{array}{ccc}j_a& I& J^{}\\ J& I^{\prime \prime }& j_b^{}\end{array}\right\}\left\{\begin{array}{ccc}j_a^{}& J^{}& I^{}\\ J& j_b& I^{\prime \prime }\end{array}\right\}`$ (111)
$`\times (JK\kappa _b^{}j_b^{}\kappa _b^{}|IK)(J^{}K\kappa _aj_a\kappa _a|IK)`$ (112)
$`\times (JK\kappa _bj_b\kappa _b|I^{}K)(J^{}K\kappa _a^{}j_a^{}\kappa _a^{}|I^{}K)`$ (113)
$`\times \chi _{JK\kappa _b^{}\nu }(j_b^{}\kappa _b^{})\chi _{JK\kappa _b\nu }^{}(j_b\kappa _b)`$ (114)
$`\times \chi _{J^{}K\kappa _a^{}\nu ^{}}(j_a^{}\kappa _a^{})\chi _{J^{}K\kappa _a\nu ^{}}^{}(j_a\kappa _a).`$ (115)
We are finally in a position to do the sums over $`I`$ and $`I^{}`$ by the application of Eq. (6.2.6) of . This leads to our final exact result. It turns out that the evaluation of the remaining interaction terms parallels that just described, differing only in the coupling constants and the single-particle functions that occur. For economy of expression, it is this total result that we finally quote:
$`{\displaystyle }={\displaystyle \underset{I^{\prime \prime }JJ^{}\nu \nu ^{}j_a\mathrm{}\kappa _a\mathrm{}}{}}[{\displaystyle \frac{1}{2}}F_{aa^{}b^{}b}(L)\chi _{JK\kappa _b^{}\nu }(b^{})\chi _{JK\kappa _b\nu }^{}(b)\chi _{J^{}K\kappa _a^{}\nu ^{}}(a^{})\chi _{J^{}K\kappa _a\nu ^{}}^{}(a)`$ (116)
$`{\displaystyle \frac{1}{2}}F_{ba^{}b^{}a}(L)\varphi _{JK\kappa _b^{}\nu }(b^{})\varphi _{JK\kappa _b\nu }^{}(b)\varphi _{J^{}K\kappa _a^{}\nu ^{}}(a^{})\varphi _{J^{}K\kappa _a\nu ^{}}^{}(a)`$ (117)
$`+G_{aa^{}b^{}b}\chi _{JK\kappa _b^{}\nu }(b^{})\varphi _{JK\kappa _b\nu }^{}(b)\varphi _{J^{}K\kappa _a^{}\nu ^{}}(a^{})\chi _{J^{}K\kappa _a\nu ^{}}^{}(a)]`$ (118)
$`\times (1)^{\kappa _b+\kappa _b^{}1+j_a^{}+j_b+L+I^{\prime \prime }}`$ (119)
$`\times {\displaystyle \frac{(2L+1)}{(2J^{}+1)}}\left\{\begin{array}{ccc}j_a^{}& j_a& L\\ j_b^{}& j_b& I^{\prime \prime }\end{array}\right\}`$ (122)
$`\times (j_a\kappa _aj_b^{}\kappa _b^{}|I^{\prime \prime }\kappa _b^{}\kappa _a)(j_a^{}\kappa _a^{}j_b\kappa _b|I^{\prime \prime }\kappa _b\kappa _a^{})`$ (123)
$`\times (I^{\prime \prime }\kappa _b^{}\kappa _aJK\kappa _b^{}|J^{}K\kappa _a)(I^{\prime \prime }\kappa _b\kappa _a^{}JK\kappa _b|J^{}K\kappa _a^{}).`$ (124)
We now collect the results of the calculations presented in this section. The strong coupling limit of the functional $``$ for the axial case is given by the sum of Eqs. (63), (77), (78), (83), and (124). By varying in turn with respect to $`\chi _{JK\kappa _a\nu }^{}(a)`$ and $`\varphi _{JK\kappa _a\nu }^{}`$, we obtain the equations of motion
$`\overline{\epsilon }_{J\nu }\chi _{JK\kappa _a\nu }(a)`$ $`=`$ $`ϵ_a^{}\chi _{JK\kappa _a\nu }(a){\displaystyle \frac{1}{2}}\overline{f}(JK)\sqrt{(J+K\kappa _a)(JK+\kappa _a+1)}`$ (135)
$`\times \sqrt{(j_a\kappa _a)(j_a+\kappa _a+1)}\chi _{JK\kappa _a1\nu }(a)`$
$`+\sqrt{(JK+\kappa _a)(J+K\kappa _a+1)}\sqrt{(j_a+\kappa _a)(j_a\kappa _a+1)}\chi _{JK\kappa _a+1\nu }(a)`$
$`+2(K\kappa _a)\kappa _a\chi _{JK\kappa _a\nu }(a)`$
$`+[F_{aa^{}b^{}b}(L)\chi _{J^{}K\kappa _b^{}\nu ^{}}(b^{})\chi _{J^{}K\kappa _b\nu ^{}}^{}(b)\chi _{JK\kappa _a^{}\nu }(a^{})`$
$`+G_{aa^{}b^{}b}(L)\chi _{J^{}K\kappa _b^{}\nu ^{}}(b^{})\varphi _{J^{}K\kappa _b\nu ^{}}^{}(b)\varphi _{JK\kappa _a^{}\nu }(a^{})]`$
$`\times {\displaystyle \frac{2L+1}{2J+1}}\left\{\begin{array}{ccc}j_a^{}& j_a& L\\ j_b^{}& j_b& I\end{array}\right\}`$
$`\times (j_a\kappa _aj_b^{}\kappa _b^{}|I\kappa _b^{}\kappa _a)(j_a^{}\kappa _a^{}j_b\kappa _b|I\kappa _b\kappa _a^{})`$
$`\times (I\kappa _b^{}\kappa _aJ^{}K\kappa _b^{}|JK\kappa _a)(I\kappa _b\kappa _a^{}J^{}K\kappa _b|JK\kappa _a^{}),`$
$`\underset{¯}{\epsilon }_{J\nu }\varphi _{JK\kappa _a\nu }(a)`$ $`=`$ $`ϵ_a^{\prime \prime }\varphi _{JK\kappa _a\nu }(a){\displaystyle \frac{1}{2}}\underset{¯}{f}(JK)\sqrt{(J+K\kappa _a)(JK+\kappa _a+1)}`$ (146)
$`\times \sqrt{(j_a\kappa _a)(j_a+\kappa _a+1)}\varphi _{JK\kappa _a1\nu }(a)`$
$`+\sqrt{(JK+\kappa _a)(J+K\kappa _a+1)}\sqrt{(j_a+\kappa _a)(j_a\kappa _a+1)}\varphi _{JK\kappa _a+1\nu }(a)`$
$`+2(K\kappa _a)\kappa _a\varphi _{JK\kappa _a\nu }(a)`$
$`[F_{ba^{}b^{}a}(L)\varphi _{J^{}K\kappa _b^{}\nu ^{}}(b^{})\varphi _{J^{}K\kappa _b\nu ^{}}^{}(b)\varphi _{JK\kappa _a^{}\nu }(a^{})`$
$`+G_{aa^{}b^{}b}(L)\varphi _{J^{}K\kappa _b^{}\nu ^{}}(b^{})\chi _{J^{}K\kappa _b\nu ^{}}^{}(b)\chi _{JK\kappa _a^{}\nu }(a^{})]`$
$`\times {\displaystyle \frac{2L+1}{2J+1}}\left\{\begin{array}{ccc}j_a^{}& j_a& L\\ j_b^{}& j_b& I\end{array}\right\}`$
$`\times (j_a\kappa _aj_b^{}\kappa _b^{}|I\kappa _b^{}\kappa _a)(j_a^{}\kappa _a^{}j_b\kappa _b|I\kappa _b\kappa _a^{})`$
$`\times (I\kappa _b^{}\kappa _aJ^{}K\kappa _b^{}|JK\kappa _a)(I\kappa _b\kappa _a^{}J^{}K\kappa _b|JK\kappa _a^{}).`$
We add the normalization conditions for the particle-rotor model that follow from Eq. (35). We find
$$\underset{\kappa _a}{}[|\chi _{JK\kappa _a\nu }(a)|^2+|\varphi _{JK\kappa _a\nu }(a)|^2]=2j_a+1.$$
(147)
## IV Alternative derivation and its cranking limit:axial case
In addition to the PRM, self-consistent or otherwise, we are interested in the cranking theory, valid in the limit in which a single-particle angular momentum $`j_a`$ may be neglected compared to the collective angular momentum. In principle, we should be able to derive this limit from the form of the theory developed in Sec. II. However, the interaction terms, as derived, do not provide a natural pathway to the limit sought. Therefore we start anew in this section, but concentrate on deriving an approximate version of the PRM in which an expansion in $`(j/J)`$ has been made, the main difference compared to the previous calculation residing in the treatment of the interaction terms. We derive an approximate version of the PRM and then introduce the additional approximation necessary to reach the cranking limit.
For present purposes it is convenient to work in coordinate-spin-isospin space, designated by $`x`$. We work with amplitudes that we refer to as coordinate coefficients of fractional parentage (CCFP),
$`V_{JM\nu }(xIM_IK)`$ $`=`$ $`JM\nu |\widehat{\psi }(x)|\overline{IM_IK},`$ (148)
$`U_{JM\nu }(xIM_IK)`$ $`=`$ $`JM\nu |\widehat{\psi }^{}(x)|\underset{¯}{IM_IK},`$ (149)
where $`\widehat{\psi }(x)`$ is the nucleon destruction operator at the space-spin-isospin point $`x`$. In terms of these amplitudes, we rewrite the variational functional $``$ of Eq. (26) as
$``$ $`=`$ $`[ϵ(xx^{})_{J\nu }\delta (xx^{})\overline{E}^{}(IK)\delta (xx^{})]V_{JM\nu }(x^{}IM_IK)V_{JM\nu }^{}(xIM_IK)`$ (157)
$`[ϵ(xx^{})_{J\nu }\delta (xx^{})\underset{¯}{E}^{}(IK)\delta (xx^{})]U_{JM\nu }(xIM_IK)U_{JM\nu }^{}(x^{}IM_IK)`$
$`+{\displaystyle \frac{1}{2}}F(xx^{}x^{\prime \prime }x^{\prime \prime \prime })V_{J^{}M^{}\nu ^{}}(x^{\prime \prime }IM_IK)V_{J^{}M^{}\nu ^{}}^{}(x^{\prime \prime \prime }I^{}M_I^{}K)`$
$`\times V_{JM\nu }(x^{}I^{}M_I^{}K)V_{JM\nu }^{}(xIM_IK)`$
$`+G(xx^{}x^{\prime \prime }x^{\prime \prime \prime })V_{J^{}M^{}\nu ^{}}(x^{\prime \prime }IM_IK)U_{J^{}M^{}\nu ^{}}^{}(x^{\prime \prime \prime }I^{}M_I^{}K)`$
$`\times U_{JM\nu }(x^{}I^{}M_I^{}K)V_{JM\nu }^{}(xIM_IK)`$
$`{\displaystyle \frac{1}{2}}F(x^{\prime \prime \prime }x^{}x^{\prime \prime }x)U_{J^{}M^{}\nu ^{}}(x^{\prime \prime }IM_IK)VU_{J^{}M^{}\nu ^{}}^{}(x^{\prime \prime \prime }I^{}M_I^{}K)`$
$`\times U_{JM\nu }(x^{}I^{}M_I^{}K)U_{JM\nu }^{}(xIM_IK).`$
We have set $`ϵ^{}=ϵ^{\prime \prime }=ϵ`$ and shall adhere to this simplification for the remainder of our presentation. To carry out the transformation to Eq. (157), we have made use of a special mode transformation to a basis in which $`ϵ(xx^{})`$ is diagonal,
$`a_\alpha `$ $`=`$ $`\phi _\alpha ^{}(x)\widehat{\psi }(x),`$ (158)
$`ϵ_a\phi _\alpha (x)`$ $`=`$ $`ϵ(xx^{})\phi _\alpha (x^{}),`$ (159)
$`F(xx^{}x^{\prime \prime }x^{\prime \prime \prime })`$ $`=`$ $`F_{\alpha \gamma \delta \beta }\phi _\alpha ^{}(x)\phi _\gamma (x^{})\phi _\delta ^{}(x^{\prime \prime })\phi _\beta (x^{\prime \prime \prime }),`$ (160)
$`G(xx^{}x^{\prime \prime }x^{\prime \prime \prime })`$ $`=`$ $`G_{\alpha \gamma \beta \delta }\phi _\delta ^{}(x^{\prime \prime \prime })\phi _\beta ^{}(x^{})\phi _\gamma (x^{})\phi _\alpha (x).`$ (161)
The major device of the present derivation is to transform from angular momentum eigenfunctions to eigenfunctions localized in angle space, a technique that has already been exploited in Sec. IV. We base the developments on expressions for the CCFP that are derived by the same initial transformations that led to Eqs. (51) and (54), namely
$`\left(\begin{array}{c}V_{JM\nu }(xM_IK)\\ U_{JM\nu }(xM_IK)\end{array}\right)`$ $`=`$ $`{\displaystyle 𝑑RD_{MM^{}}^{(J)}(R)\left(\begin{array}{c}\chi _{JM^{}\nu }(Rx,K)\\ \varphi _{JM^{}\nu }(Rx,K)\end{array}\right)\sqrt{\frac{2I+1}{8\pi ^2}}D_{M_IK}^{(I)}(R)},`$ (166)
$`\chi _{JM\nu }(Rx,K)`$ $`=`$ $`JM\nu |\widehat{\psi }(Rx)|\widehat{0}K,`$ (167)
$`\varphi _{JM\nu }(Rx,K)`$ $`=`$ $`JM\nu |\widehat{\psi }^{}(Rx)|\widehat{0}K,`$ (168)
$`\widehat{\psi }(Rx)`$ $`=`$ $`U^1(R)\widehat{\psi }(x)U(R).`$ (169)
When we substitute Eq. (166) into the first two terms of Eq. (157), we encounter the restricted completeness relation
$$\underset{IM}{}D_{MK}^{(I)}(R)D_{MK}^{(I)}(R^{})\frac{2I+1}{8\pi ^2}=\delta (\widehat{n}\widehat{n}^{})\mathrm{exp}[iK(\gamma \gamma ^{})]\frac{1}{2\pi }.$$
(170)
To perform the integral over $`\gamma ^{}`$, we note the relation (41), and in the further calculation we utilize basic properties of the $`D`$ functions, following from their definition, Eq. (46), the first one also used extensively for the interaction term to be computed below,
$`{\displaystyle \underset{M}{}}D_{M^{}M}^{(J)}(R^1)D_{MM^{\prime \prime }}^{(J)}(R^{})`$ $`=`$ $`D_{M^{}M^{\prime \prime }}^{(J)}(R^1R^{}),`$ (171)
$`D_{MM^{}}^{(J)}(0)`$ $`=`$ $`\delta _{M,M^{}}.`$ (172)
Finally, invoking the rotational invariance of the Hamiltonian, which for the term under consideration means that $`ϵ(RxRx^{})=ϵ(xx^{})`$, we obtain the result (recall that factors of $`8\pi ^2`$ are suppressed)
$$[ϵ(xx^{})_{J\nu }\delta (xx^{})]\chi _{JM\nu }(x^{})\chi _{JM\nu }^{}(x).$$
(173)
We note the corresponding contributions
$$[ϵ(xx^{})_{J\nu }\delta (xx^{})]\varphi _{JM\nu }(x)\varphi _{JM\nu }^{}(x^{}).$$
(174)
Applying Eqs. (171) and (172) to the interaction term, we find at the stage that the sums over $`IM`$ and $`I^{}M^{}`$ have been carried out,
$`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑R𝑑R^{}F(xx^{}x^{\prime \prime }x^{\prime \prime \prime })D_{M^{}M^{\prime \prime }}^{(J)}(R^1R^{})D_{M^{iv}M^{\prime \prime \prime }}^{(J^{})}(R^1R^{})}`$ (175)
$`\times \chi _{JM^{}\nu }(Ry^{})\chi _{JM^{\prime \prime }\nu }^{}(R^{}y)\chi _{J^{}M^{\prime \prime \prime }\nu ^{}}(R^{}x^{})\chi _{J^{}\nu ^{iv}\nu ^{}}^{}(Rx).`$ (176)
Introducing the definition
$$R^1R^{}=,$$
(177)
and replacing the integral over $`R^{}`$ by an integral $``$, we could do the integrals exactly by decomposing the amplitudes $`\chi `$ into irreducible tensors. We have resisted the temptation to do this, since a full calculation was carried out in Sec. III. It is more illuminating, as well as simpler, to proceed approximately by expanding $``$ about the unit matrix where ever it appears as the argument of a $`\chi `$ function. This brings in at each order angular momentum operators acting on single-particle wave-functions and therefore dimensionally is the source of the expansion in $`(j/J)`$. For the interaction term the cranking limit will arise from the leading term of this expansion.
With the help of the rotational invariance of the interaction, in the present instance the relation for example $`F(RxRx^{}Rx^{\prime \prime }Rx^{\prime \prime \prime })=F(xx^{}x^{\prime \prime }x^{\prime \prime \prime })`$, and the orthonormality relations of the $`D`$ functions, we reach the result
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{JMM^{}\nu \nu ^{}}{}}{\displaystyle \frac{1}{2J+1}}F(xx^{}x^{\prime \prime }x^{\prime \prime \prime })\chi _{JM\nu ^{}}(x^{\prime \prime })\chi _{JM^{}\nu ^{}}^{}(x^{\prime \prime \prime })\chi _{JM^{}\nu }(x^{})\chi _{JM\nu }^{}(x),`$ (178)
which is almost the cranking limit. The remaining interaction terms may be written down by inspection, namely
$`{\displaystyle \underset{JMM^{}\nu \nu ^{}}{}}{\displaystyle \frac{1}{2J+1}}[G(xx^{}x^{\prime \prime }x^{\prime \prime \prime })\chi _{JM\nu ^{}}(x^{\prime \prime })\varphi _{JM^{}\nu ^{}}^{}(x^{\prime \prime \prime })\varphi _{JM^{}\nu }(x^{})\chi _{JM\nu }^{}(x)`$ (179)
$`{\displaystyle \frac{1}{2}}F(x^{\prime \prime \prime }x^{}x^{\prime \prime }x)\varphi _{JM\nu ^{}}(x^{\prime \prime })\varphi _{JM^{}\nu ^{}}^{}(x^{\prime \prime \prime })\varphi _{JM^{}\nu }(x^{})\varphi _{JM\nu }^{}(x)].`$ (180)
It remains for us to calculate the Coriolis coupling terms. Remarking that $`\overline{E}^{}(IK)`$ is an eigenvalue of $`D_{MK}^{(I)}`$,
$$\overline{E}^{}(IK)D_{MK}^{(I)}=\overline{E}^{}(\stackrel{}{I}_{\mathrm{op}}^2K^2)D_{MK}^{(I)},$$
(181)
using completeness and integrating by parts, we reach an intermediate stage of the calculation, in the form,
$`{\displaystyle 𝑑R[\overline{E}^{}(\stackrel{}{I}_{\mathrm{op}}^2K^2)D_{MM^{}}^{(J)}(R)\chi _{JM^{}\nu }(Rx)]D_{MM^{}}^{(J)}\chi ^{}JM^{\prime \prime }\nu (Rx)}.`$ (182)
By distribution of the derivatives, $`\stackrel{}{I}_{\mathrm{op}}\stackrel{}{J}_{\mathrm{op}}+\stackrel{}{j}_{\mathrm{op}}`$, by noting the relation
$$\stackrel{}{j}_{\mathrm{op}}(R)\chi (Rx)=\stackrel{}{j}(x)\chi (Rx),$$
(183)
and by expansion in powers of $`(j/J)`$, Eq. (182) becomes to first order (in the intrinsic system)
$`{\displaystyle \underset{JM\nu }{}}\overline{E}^{}[J(J+1)K^2]|\chi _{JM\nu }(x)|^2`$ (184)
$`+[{\displaystyle \frac{\overline{E}^{}}{J_i}}D_{MM^{}}^{(J)}(R)][j_i(x)\chi _{JM^{}\nu }(Rx)]D_{MM^{\prime \prime }}^{(J)}(R)\chi _{JM^{\prime \prime }\nu }^{}(Rx)+\mathrm{}`$ (185)
The second term, the Coriolis coupling may be evaluated further by substituting the general relation (65) and the matrix elements for intrinsic components,
$`J_{}D_{MM^{}}^{(J)}`$ $`=`$ $`\sqrt{(J+M^{})(JM^{}+1)}D_{MM^{}1}^{(J)},`$ (186)
$`J_+D_{MM^{}}^{(J)}`$ $`=`$ $`\sqrt{(JM^{})(J+M^{}+1)}D_{MM^{}+1}^{(J)},`$ (187)
yielding
$`{\displaystyle \underset{JM\nu }{}}\overline{f}(JK)\{{\displaystyle \frac{1}{2}}\sqrt{(JM)(J+M+1)}j_+\chi _{JM+1\nu }(x)`$ (188)
$`+{\displaystyle \frac{1}{2}}\sqrt{(J+M)(JM+1)}j_{}\chi _{JM1\nu }(x)+Mj_3\chi _{JM\nu }(x)\}\chi _{JM\nu }^{}(x)\}.`$ (189)
The corresponding calculation for the other Coriolis term yields the sum
$`{\displaystyle \underset{JM\nu }{}}\underset{¯}{E}^{}[J(J+1)K^2]|\varphi _{JM\nu }(x)|^2`$ (190)
$`+{\displaystyle \underset{JM\nu }{}}\underset{¯}{f}(JK)\{{\displaystyle \frac{1}{2}}\sqrt{(JM)(J+M+1)}j_+\varphi _{JM+1\nu }(x)`$ (191)
$`+{\displaystyle \frac{1}{2}}\sqrt{(J+M)(JM+1)}j_{}\varphi _{JM1\nu }(x)+Mj_3\varphi _{JM\nu }(x)\}\varphi ^{}_{JM\nu }(x).`$ (192)
The signs in the Coriolis terms (189) and (192) appear to be reversed compared to those encountered in Eqs. (77) and (83), but when due account is taken of Eq. (183), there is no inconsistency.
We are finally ready to discuss the cranking limit. The essential observation is that once the expansion to leading order in $`(j/J)`$ has been made both in the Coriolis coupling and in the interaction terms, the resulting approximate functional $``$ presents itself as a single sum over $`J`$. However, angular momentum is still conserved at this juncture. We lose angular momentum conservation by assuming that consistent with the condition $`(j/J)<<1`$ we may identify $`M`$ and $`K`$, i. e., we may neglect the angular momentum transferred to or from the particle, and write, furthermore, ($`\omega `$ defined below)
$`\chi _{JK\nu }(x)`$ $``$ $`\sqrt{2J(\omega )+1}\chi _{\omega \nu }(x),`$ (193)
$`\chi _{JK\pm 1\nu }(x)`$ $``$ $`\sqrt{2J(\omega )+1}C_{}\chi _{\omega \nu }(x),`$ (194)
i. e., the amplitudes differing in $`K`$ from the “central value” by a unit are assumed proportional to the central amplitude (which is defined as the cranking amplitude) up to scale factors $`C_{}`$ discussed below. Similar definitions hold for the $`\varphi `$ amplitudes. The factor $`\sqrt{2J+1}`$ is inserted for convenience, as will be evident from Eq. (202) given below.
These assumptions suggest the following definitions of the components of the angular frequency (overline and underline understood)
$`\omega _{}(K)`$ $`=`$ $`f(JK)C_{}\sqrt{(JK+1)(J\pm K)},`$ (195)
$`\omega _3`$ $`=`$ $`f(JK)K.`$ (196)
The introduction of the factors $`C_{}`$ may appear gratuitous at first sight, but it is needed, as will become especially evident when we treat the triaxial case, to guarantee that in the cranking limit the theorem that the angular velocity is proportional to the angular momentum is valid .
Remembering the definition (64) and reinstating Cartesian intrinsic coordinates for the Coriolis coupling terms, we obtain the cranking variational expression
$`[/(2J(\omega )+1)]`$ $`=`$ $`ϵ(xx^{})\chi _{\omega \nu }(x^{})\chi _{\omega \nu }^{}(x)ϵ(xx^{})\varphi _{\omega \nu }(x)\varphi _{\omega \nu }^{}(x^{})`$ (202)
$`+(\overline{\omega }_ij_i\chi _{\omega \nu }(x))\chi _{\omega \nu }^{}(x)+(\underset{¯}{\omega }_ij_i\varphi _{\omega \nu }(x))\varphi _{\omega \nu }^{}(x)`$
$`+{\displaystyle \frac{1}{2}}F(xx^{}x^{\prime \prime }x^{\prime \prime \prime })\chi _{\omega \nu ^{}}(x^{\prime \prime })\chi _{\omega \nu ^{}}^{}(x^{\prime \prime \prime })\chi _{\omega \nu }(x^{})\chi _{\omega \nu }^{}(x)`$
$`+G(xx^{}x^{\prime \prime }x^{\prime \prime \prime })\chi _{\omega \nu ^{}}(x^{\prime \prime })\varphi _{\omega \nu ^{}}^{}(x^{\prime \prime \prime })\chi _{\omega \nu }(x^{})\varphi _{\omega \nu }^{}(x)`$
$`{\displaystyle \frac{1}{2}}F(x^{\prime \prime \prime }x^{}x^{\prime \prime }x)\varphi _{\omega \nu ^{}}(x^{\prime \prime })\varphi _{\omega \nu ^{}}^{}(x^{\prime \prime \prime })\varphi _{\omega \nu }(x^{})\varphi _{\omega \nu }^{}(x)`$
$`\overline{\epsilon }_{\omega \nu }\chi _{\omega \nu }(x)\chi _{\omega \nu }^{}(x)\underset{¯}{\epsilon }_{\omega \nu }\varphi _{\omega \nu }(x)\varphi _{\omega \nu }^{}(x).`$
The equations of motion that follow are number-conserving, and according to the definitions (195) and (196) allow solutions with principal axis cranking
## V Triaxial rotor: core-particle coupling model and cranking limit
In this section, we assume that states of interest of neighboring even nuclei can be described phenomenologically by a Hamiltonian
$$_c=\frac{1}{2}a_iI_i^2+\frac{1}{4}a_{ij}\{I_i^2,I_j^2\}+\mathrm{}.$$
(203)
In the calculations to be described below, we shall retain only the first term of $`_c`$. The underlying model arises as follows: We assume that we can identify states of the appropriate even nucleus as $`|IM_In\sigma `$, which we read as the $`n`$th state of angular momentum $`I`$ belonging to a triaxial intrinsic structure $`\sigma `$. We also define a rotated intrinsic state
$$|R\sigma =U(R)|\widehat{0}\sigma .$$
(204)
It is part of the definition of the model that the scalar product
$$R\sigma |IM_In\sigma F_{M_In}^{(I)}(R)$$
(205)
satisfies the eigenvalue equation
$$_cF_{M_In}^{(I)}=E^{}(In)F_{M_In}^{(I)}.$$
(206)
Further useful equations satisfied by or defining the model include
$`|IM_In\sigma `$ $`=`$ $`|IM_IK\sigma c_{Kn}^{(I\sigma )}`$ (207)
$`\delta _{nn^{}}`$ $`=`$ $`{\displaystyle \underset{K}{}}c_{Kn}^{(I\sigma )}c_{Kn^{}}^{(I\sigma )}`$ (208)
$`\delta _{KK^{}}`$ $`=`$ $`{\displaystyle \underset{n}{}}c_{Kn}^{(I\sigma )}c_{K^{}n}^{(I\sigma )},`$ (209)
$`R\sigma |IMK\sigma `$ $`=`$ $`\sqrt{{\displaystyle \frac{2I+1}{8\pi ^2}}}D_{MK}^{(I)}(R).`$ (210)
We turn to the evaluation of the terms in the variational functional $``$. We shall follow the methods of both Sec. III and Sec. IV, depending on the aim of a particular fragment of the calculation. Starting from the representation
$`V_{JM\nu }(\alpha IM_In\sigma )`$ $`=`$ $`JM\nu |a_\alpha |R\sigma F_{M_In}^{(I)}(R)`$ (211)
$`=`$ $`JM\nu |a_\alpha |R\sigma D_{M_IK}^{(I)}(R)c_{Kn}^{(I\sigma )},`$ (212)
we can derive a formula for the current version of the CFP $`V`$ that is analogous to Eq. (51), namely
$`V_{JM\nu }(\alpha ;IM_In\sigma )`$ $`=`$ $`{\displaystyle \underset{K\kappa _a}{}}\sqrt{{\displaystyle \frac{8\pi ^2}{2j_a+1}}}(1)^{JM}`$ (215)
$`\times (IM_IJM|j_am_a)(JKj_a\kappa _a|IK+\kappa _a)`$
$`\times (1)^{j_a+\kappa _a}\chi _{JK\nu }(j_a\kappa _a\sigma )c_{K+\kappa _an}^{(I\sigma )},`$
$`(1)^{j+m}\chi _{JK\nu }(jm\sigma )`$ $`=`$ $`JK\nu |a_{jm}|\widehat{0}\sigma .`$ (216)
The corresponding formula for the CFP $`U`$ is
$`U_{JM\nu }(\alpha IM_In\sigma )`$ $`=`$ $`{\displaystyle \underset{K\kappa _a}{}}\sqrt{{\displaystyle \frac{8\pi ^2}{2j_a+1}}}`$ (219)
$`\times (1)^{JM+j_a\kappa _a+j_a+m_a}(IM_IJM|j_am_a)`$
$`\times (JKj_a\kappa _a|IK+\kappa _a)\varphi _{JK\nu }(j_a\kappa _a\sigma )c_{K+\kappa _an}^{(I\sigma )},`$
$`\varphi _{JM\nu }(j_a\kappa _a\sigma )`$ $`=`$ $`JM\nu |a_{j_a\kappa _a}^{}|\widehat{0}\sigma .`$ (220)
With these formulas, we find the contributions of the simplest single-particle terms to take the form, in the shell-model or mode representation,
$`{\displaystyle \underset{JK\nu j_a\kappa _a}{}}[(ϵ_a\overline{}_{J\nu })|\chi _{JK\nu }(a)|^2`$ (221)
$`(ϵ_a+\underset{¯}{}_{J\nu }|\varphi _{JK\nu }(a)|^2].`$ (222)
We study next the term involving the Lagrange multiplier $`\overline{E}^{}(In)`$. Here it is convenient to carry out the calculation by a method analogous to that utilized beginning with Eq. (181). With the help of the defining equation (206) and a subsequent integration by parts, we have first of all
$`\overline{E}^{}(In)V_{JM\nu }(IM_In\sigma )`$ $`=`$ $`{\displaystyle 𝑑R[_c(\widehat{I}_i)JM\nu |a_\alpha |R\sigma ]F_{M_In}^{(I\sigma )}(R)}.`$ (223)
With the help of the completeness relation
$$\underset{IM_In}{}F_{M_In}^{(I\sigma )}(R)F_{IM_In}^{(I\sigma )}(R^{})=\delta (RR^{}),$$
(224)
we thus find for the total term
$$\overline{E}^{}(In\sigma )|V_{JM\nu }(\alpha IM_In\sigma )|^2=dR[_c(\widehat{I}_i)JM\nu |a_\alpha |R\sigma ]JM\nu |a_\alpha |R\sigma ^{}.$$
(225)
The square bracket may be reexpressed as
$$_c(\widehat{I}_i)JM\nu |a_\alpha |R\sigma =[_c(\widehat{I}_i)D_{MK}^{(J)}(R)D_{m_a\kappa _a}^{(j_a)}(R)]JK\nu |a_{j_a\kappa _a}|\widehat{0}\sigma .$$
(226)
As far as the application of $`_c`$ in (226) is concerned, we then write
$`_c(\widehat{I}_i)`$ $``$ $`_c(\widehat{J}_i+\widehat{j}_i)`$ (227)
$`=`$ $`_c(\widehat{J}_i)+{\displaystyle \frac{_c}{\widehat{J}_i}}\widehat{j}_i+\mathrm{},`$ (228)
and work only to the order indicated explicitly.
At the same time it is convenient to rewrite
$`_c(\widehat{I}_i)`$ $`=`$ $`{\displaystyle \frac{1}{4}}b_1(I_+I_{}+I_{}I_+)+{\displaystyle \frac{1}{4}}b_2(I_+^2+I_{}^2)+{\displaystyle \frac{1}{2}}b_3I_3^2`$ (229)
$`a_1=b_1+b_2,`$ $`a_2=b_1b_2,a_3=b_3.`$ (230)
It is now straightforward to calculate the contributions arising from the two terms of Eq. (228). For the first term we find
$`8\pi ^2\{[{\displaystyle \frac{1}{2}}b_1[J(J+1)K^2]+{\displaystyle \frac{1}{2}}b_3K^2]|\chi _{JK\nu }(a)|^2`$ (231)
$`{\displaystyle \frac{1}{4}}b_2\sqrt{(JK+2)(JK+1)(J+K1)(J+K)}\chi _{JK2\nu }(a)\chi _{JK\nu }^{}(a)`$ (232)
$`+{\displaystyle \frac{1}{4}}b_2\sqrt{(J+K+2)(J+K+1)(JK1)(JK)}`$ (233)
$`\times \chi _{JK+2\nu }(a)\chi _{JK\nu }^{}(a)\},`$ (234)
and for the second term,
$`8\pi ^2[{\displaystyle \frac{1}{2}}b_1\sqrt{(JK+1)(J+K)(j_a+\kappa _a+1)(j_a\kappa _a)}\chi _{JK1\nu }(j_a\kappa _a+1\sigma )`$ (235)
$`+{\displaystyle \frac{1}{2}}b_1\sqrt{(J+K+1)(JK)(j_a\kappa _a+1)(j_a+\kappa _a)}\chi _{JK+1\nu }(j_a\kappa _a1\sigma )`$ (236)
$`+{\displaystyle \frac{1}{2}}b_2\sqrt{(JK+1)(J+K)(j_a\kappa _a+1)(j_a+\kappa _a)}\chi _{JK1\nu }(j_a\kappa _a1\sigma )`$ (237)
$`+{\displaystyle \frac{1}{2}}b_2\sqrt{(J+K+1)(JK)(j_a+\kappa _a+1)(j_a\kappa _a)}\chi _{JK+1\nu }(j_a\kappa _a+1\sigma )`$ (238)
$`+b_3K\kappa _a\chi _{JK\nu }(j_a\kappa _a\sigma )]\chi ^{}_{JK\nu }(j_a\kappa _a\sigma ).`$ (239)
Both of these terms can be identified as familiar structures. By means of this identification we shall have achieved both a simpler form for the particle-rotor formalism and for its limiting case, the cranking formalism. First consider Eq. (234). Note that the content of Eqs. (206) and (207) can be rewritten as
$`_cD_{M_IK}^{(I)}`$ $`=`$ $`D_{M_IK^{}}^{(I)}(_c)_{K^{}K},`$ (240)
$`()_{KK^{}}c_{K^{}n}^{(I\sigma )}`$ $`=`$ $`E^{}(In)c_{Kn}^{(I\sigma )}.`$ (241)
This eigenvalue equation was associated with even nuclei and thus with integer values of the angular momentum. By analytic continuation, we can define a corresponding eigenvalue equation for odd nuclei as follows:
$$(_c(\widehat{J}_i))_{KK^{}}c_{K^{}\tau }^{(J)}=E^{}(J\tau )c_{K\tau }^{(J)},$$
(242)
where $`J,K`$ are now half-integral. We then see that if we introduce a new set of particle amplitudes $`\chi _{J\tau \nu }`$ by means of the equation
$$\chi _{JK\nu }(j\kappa )=c_{K\tau }^{(J)}\chi _{J\tau \nu }(j\kappa ),$$
(243)
we can transform Eq. (234) into the form
$$E^{}(J\tau )|\chi _{J\tau \nu }(a)|^2.$$
(244)
Finally, as we did for the axial case, we can combine energy terms by means of a definition
$$\epsilon _{J\nu }=_{J\nu }+E^{}(J\tau ).$$
(245)
We turn our attention next to Eq. (239). We note first that this expression is an expanded version of
$$[a_i\widehat{J}_i\widehat{j}_i\chi _{JK\nu }(a)]\chi _{JK\nu }^{}(a),$$
(246)
where $`\widehat{J}_i`$ acts on the value of $`K`$ and $`\widehat{j}_i`$ acts on the value of $`\kappa _a`$. Transforming to the new amplitudes $`\chi _{J\tau \nu }`$, expression (246) becomes
$$[a_i\widehat{J}_i\widehat{j}_i\chi _{J\tau \nu }(a)]\chi _{J\tau \nu }^{}(a),$$
(247)
where now
$`\widehat{J}_i\chi _{J\tau \nu }`$ $`=`$ $`\chi _{J\tau ^{}\nu }(J\tau ^{}|J_i|J\tau ),`$ (248)
$`(J\tau ^{}|\widehat{J_i}|J\tau )`$ $`=`$ $`c_{K^{}\tau ^{}}^{(J)}(JK^{}|J_i|JK)c_{K\tau }^{(J)}.`$ (249)
For the purpose of taking the cranking limit and comparing the forms derived in Sec. IV, we rewrite the results found so far and the corresponding terms involving $`\varphi `$ amplitudes in coordinate space. For this we require only Eqs. (243), the corresponding equations
$`\chi _{JK\nu }(x)`$ $`=`$ $`JK\nu |\widehat{\psi }(x)|\widehat{0}\sigma `$ (250)
$`=`$ $`c_{K\tau }^{(J)}\chi _{J\tau \nu }(x),`$ (251)
and the similar equations for the terms involving $`\varphi `$. We thus find the contributions
$`\overline{\epsilon }_{J\nu }|\chi _{J\tau \nu }(x)|^2+ϵ(xx^{})\chi _{J\tau \nu }(x^{})\chi _{J\tau \nu }^{}(x)+\overline{a}_i[J_ij_i(x)\chi _{J\tau \nu }(x)]\chi _{J\tau \nu }^{}(x)`$ (252)
$`\underset{¯}{\epsilon }_{J\nu }|\varphi _{J\tau \nu }(x)|^2+ϵ(xx^{})\varphi _{J\tau \nu }(x^{})\varphi _{J\tau \nu }^{}(x)+\underset{¯}{a}_i[J_ij_i(x)\varphi _{J\tau \nu }(x)]\varphi _{J\tau \nu }^{}(x).`$ (253)
The cranking limit of these terms may now be taken by means of the replacements that generalize Eqs. (194) and (193),
$$\chi _{J\tau ^{}\nu }(x)\sqrt{2J+1}C_\tau (\tau ^{})\chi _{\omega \nu }(x),$$
(254)
and $`\tau ^{}`$ refers to $`\tau `$ or any of the values coupled to $`\tau `$ by the matrices of $`\widehat{J}_i`$, with $`C_\tau (\tau )=1`$. This is the essential blurring of angular momentum conservation that takes us from the conserving particle-rotor approximation to the cranking approximation. It allows us as well to define the components of the angular velocity in generalization of Eq. (193).
$`\overline{\omega }_i(\tau )`$ $`=`$ $`\overline{a}_i{\displaystyle \underset{\tau ^{}}{}}C_\tau (\tau ^{})(J\tau ^{}|J_i|J\tau ),`$ (255)
$`\tau ^{}`$ $`=`$ $`\tau ^{}(\tau ).`$ (256)
As usual, there are corresponding equations for the amplitudes $`\varphi `$.
We may thus replace Eq. (253) by its cranking limit
$`(2J+1)[\overline{\epsilon }_{\omega \nu }|\chi _{\omega \nu }(x)|^2+ϵ(xx^{})\chi _{\omega \nu }(x^{})\chi _{\omega \nu }^{}(x)+\overline{\omega }_ij_i(x)\chi _{\omega \nu }(x)]\chi _{\omega \nu }^{}(x)`$ (257)
$`\underset{¯}{\epsilon }_{\omega \nu }|\varphi _{\omega \nu }(x)|^2+ϵ(xx^{})\varphi _{\omega \nu }(x^{})\varphi _{\omega \nu }^{}(x)+\underset{¯}{\omega }_ij_i(x)\varphi _{\omega \nu }(x)]\varphi ^{}_{\omega \nu }(x),`$ (258)
which is indistinguishable in form from the corresponding terms of Eq. (202).
It remains for us to compute the contributions of the interaction terms. We consider first an exact calculation analogous to that carried out in Sec. III, starting from the representations (215) and (219) for the CFP in the triaxial case. It is straightforward to follow the calculations that begin with Eq. (93) and culminate with the result (124), as soon as one utilizes the orthonormality relations involving the coefficients $`c_{Kn}^{(I)}`$ at the first step. The final result is
$`{\displaystyle }={\displaystyle \underset{I^{\prime \prime }JJ^{}KK^{}\nu \nu ^{}j_a\mathrm{}\kappa _a\mathrm{}}{}}[{\displaystyle \frac{1}{2}}F_{aa^{}b^{}b}(L)\chi _{JK\kappa _b^{}\nu }(b^{})\chi _{JK^{}\kappa _b\nu }^{}(b)\chi _{J^{}K^{}\kappa _a^{}\nu ^{}}(a^{})\chi _{J^{}K\kappa _a\nu ^{}}^{}(a)`$ (259)
$`{\displaystyle \frac{1}{2}}F_{ba^{}b^{}a}(L)\varphi _{JK\kappa _b^{}\nu }(b^{})\varphi _{JK^{}\kappa _b\nu }^{}(b)\varphi _{J^{}K^{}\kappa _a^{}\nu ^{}}(a^{})\varphi _{J^{}K\kappa _a\nu ^{}}^{}(a)`$ (260)
$`+G_{aa^{}b^{}b}\chi _{JK\kappa _b^{}\nu }(b^{})\varphi _{JK^{}\kappa _b\nu }^{}(b)\varphi _{J^{}K^{}\kappa _a^{}\nu ^{}}(a^{})\chi _{J^{}K\kappa _a\nu ^{}}^{}(a)]`$ (261)
$`\times (1)^{\kappa _b+\kappa _b^{}1+j_a^{}+j_b+L+I^{\prime \prime }}`$ (262)
$`\times {\displaystyle \frac{(2L+1)}{(2J^{}+1)}}\left\{\begin{array}{ccc}j_a^{}& j_a& L\\ j_b^{}& j_b& I^{\prime \prime }\end{array}\right\}`$ (265)
$`\times (j_a\kappa _aj_b^{}\kappa _b^{}|I^{\prime \prime }\kappa _b^{}\kappa _a)(j_a^{}\kappa _a^{}j_b\kappa _b|I^{\prime \prime }\kappa _b\kappa _a^{})`$ (266)
$`\times (I^{\prime \prime }\kappa _b^{}\kappa _aJK\kappa _b^{}|J^{}K\kappa _a)(I^{\prime \prime }\kappa _b\kappa _a^{}JK^{}\kappa _b|J^{}K^{}\kappa _a^{}).`$ (267)
Superficially, the change compared to Eq. (124) is that instead of a fixed value of $`K`$, we have a double sum over $`K`$ and $`K^{}`$. The same expression holds for a finite number of interacting $`K`$ bands provided the sums are restricted correspondingly.
Finally, we consider the calculation of the interaction term by the method of Sec. IV, needed to obtain the cranking limit. Here, in place of Eqs. (166)-(168), we utilize the forms
$`\left(\begin{array}{c}V_{JM\nu }(xM_In)\\ U_{JM\nu }(xM_In)\end{array}\right)`$ $`=`$ $`{\displaystyle 𝑑RD_{MM^{}}^{(J)}(R)\left(\begin{array}{c}\chi _{JM^{}\nu }(Rx,\sigma )\\ \varphi _{JM^{}\nu }(Rx,\sigma )\end{array}\right)\sqrt{\frac{2I+1}{8\pi ^2}}F_{M_In}^{(I)}(R)},`$ (272)
$`\chi _{JM\nu }(Rx,\sigma )`$ $`=`$ $`JM\nu |\widehat{\psi }(Rx)|\widehat{0}\sigma ,`$ (273)
$`\varphi _{JM\nu }(Rx,\sigma )`$ $`=`$ $`JM\nu |\widehat{\psi }^{}(Rx)|\widehat{0}\sigma .`$ (274)
Once the full completeness relation (224) is utilized instead of the restricted completeness relation (170), the calculation mimics the one carried out in Sec. IV. In terms of the amplitudes $`\chi _{J\tau \nu }`$ and $`\varphi _{J\tau \nu }`$, the result is
$`{\displaystyle \underset{J\tau \tau ^{}\nu \nu ^{}}{}}{\displaystyle \frac{1}{2J+1}}\{{\displaystyle \frac{1}{2}}F(xx^{}x^{\prime \prime }x^{\prime \prime \prime })\chi _{J\tau \nu ^{}}(x^{\prime \prime })\chi _{J\tau ^{}\nu ^{}}^{}(x^{\prime \prime \prime })\chi _{J\tau ^{}\nu }(x^{})\chi _{J\tau \nu }^{}(x)`$ (275)
$`+G(xx^{}x^{\prime \prime }x^{\prime \prime \prime })\chi _{J\tau \nu ^{}}(x^{\prime \prime })\varphi _{J\tau ^{}\nu ^{}}^{}(x^{\prime \prime \prime })\varphi _{J\tau ^{}\nu }(x^{})\chi _{J\tau \nu }^{}(x)`$ (276)
$`{\displaystyle \frac{1}{2}}F(x^{\prime \prime \prime }x^{}x^{\prime \prime }x)\varphi _{J\tau \nu ^{}}(x)\varphi _{J\tau ^{}\nu ^{}}^{}(x^{})\varphi _{J\tau ^{}\nu }(x^{\prime \prime \prime })\varphi _{J\tau \nu }^{}(x^{\prime \prime })\}.`$ (277)
The cranking limit of this expression is indistinguishable from the corresponding terms of Eq. (202) just as was the case for the single-particle terms (258). Thus the form of the cranking variational principle for the triaxial is indistinguishable from that for the axial case and need not be written again. It is understood, however, that we are dealing with full three-dimensional cranking, and that the single-particle wave functions have suitably modified symmetry.
## VI Summary and discussion
We have studied the microscopic foundations of the particle-rotor model and of the cranking model for both axial and triaxial nuclei. The microscopic model was chosen in a form in which the interaction is given at the outset as a sum of multipole and pairing forces. We carried out the study from the point of view of the Kerman-Klein method based on the equations of motion for single fermion operators, and this choice of interaction has the advantage that the c-number equations of motion for the CFP are formally exact. These equations of motion and an associated variational principle, worked out in Sec. II, form the basis for the remaining considerations.
\[In the earliest papers on the KK approach , a more general shell-model interaction was used in the derivation of c-number equations. An essential part of the derivation involved the introduction of the physical arguments needed to separate this interaction into multipole and pairing contributions. As a consequence of the limitations of this procedure, the equations of motion found from it are not exact. Nevertheless, the final equations are formally equivalent to those utilized in this paper. The explanation for this concordance is that in the approach in this paper, the “error” involved in the separation has already been built into the starting Hamiltonian, as a further compromise, widely accepted, in the definition of the microscopic theory.\]
As the first application, we derived in Sec. III a self-consistent particle-rotor model for axially symmetric nuclei. The derivation was carried out using basic ideas developed in , where, starting from a semi-microscopic version of the theory, we derived the standard non-self-consistent version of the particle-rotor model. The present discussion complements the previous one in the sense that it starts from the beginning and carries the reasoning up to the edge of the semi-microscopic form.
We began this work with the prejudice that, as apposed to previous treatments, a natural path to the cranking model involved passing through the particle-rotor model. Though we were ultimately able to confirm this prejudice, the version of the particle-rotor model derived in Sec. III, though a useful one for applications , does not appear to be useful in the further transition to cranking. For this purpose we must be able to expand all contributions in powers of $`(j/J)`$, the ratio of a characteristic single-particle angular momentum to the collective angular momentum. We have not discovered such an expansion for the interaction forms derived in this first treatment. Therefore, in Sec. IV we start anew, utilizing an approach already described briefly for two-dimensional rotations in an early publication . Rather than pushing through to a formally exact result, we stop the calculation at the leading order of the small parameter, and thus obtain an approximate version of the particle-rotor model that still conserves angular momentum, but is only a step away from the cranking limit. This further step violates angular momentum conservation by the way in which an angular velocity is introduced to replace the collective angular momentum. In Sec. V the considerations of both previous sections are generalized to the triaxial case.
Several special features of our treatment should be highlighted. For the axial case, as soon as the neighboring even nuclei are represented by bands with non-vanishing $`K`$ values, we have tilted cranking in its simplest form. A fortiori, in the triaxial case we derive the possibility of full three-dimensional cranking. Within our mode of analysis, these statements may be taken to have the status of theorems. Another feature of our derivations of cranking models is that number conservation is maintained.
Nevertheless, in the light of recent developments associated with tilted cranking , possible limitations on our work have to be addressed. Superficially, our results apply to one quasiparticle spectra of odd nuclei, whereas the current focus of interest is on at least two quasiparticle spectra of even nuclei, and even more on multi-quasiparticle states. In principle, however, these examples are covered by our considerations. Thus the two quasiparticle case is readily derived from the formalism developed in Appendix A. The multiple quasiparticle case is covered if one replaces the reference ground states of the even nuclei by suitably chosen band heads of two quasiparticle bands. Details of such calculations are best addressed within the framework of specific applications.
## A Further development of the theory: Iterative solution schemes
The theory developed in Sec. II was sufficient for the purposes of the remainder of the body of this work, a study of the strong-coupling limit. However, we cannot resist the temptation to show how this theory can be developed further, and the full architecture used to suggest algorithms for the solution of the non-linear problem thus defined. We point out before proceeding, moreover, that applications of algorithms similar to some of those to be described were carried out in our early work .
To begin the extended development, it is helpful to introduce a more concise representation of the equations of motion (13) and (15) for the CFP by defining the vector
$$\mathrm{\Psi }_i(\alpha n)=\left(\begin{array}{c}V_i(\alpha n)\\ U_i(\alpha n)\end{array}\right).$$
(A1)
The equations of motion can then be written
$$_i\mathrm{\Psi }_i(\alpha n)=(\alpha n,\beta n^{})\mathrm{\Psi }_i(\beta n^{}),$$
(A2)
where $`(\alpha n,\beta n^{})`$ is the Hermitian matrix
$`(\alpha n,\beta n^{})=`$ (A3)
$`\left(\begin{array}{cc}(ϵ_a^{}E_{\overline{n}}^{})\delta _{nn^{}}\delta _{\alpha \beta }+\overline{\mathrm{\Gamma }}(\alpha n,\beta n^{})& \mathrm{\Delta }(\alpha n,\beta n^{})\\ \mathrm{\Delta }^{}(\beta n^{},\alpha n)& (ϵ_a^{\prime \prime }E_{\underset{¯}{n}}^{})\delta _{nn^{}}\delta _{\alpha \beta }+\underset{¯}{\mathrm{\Gamma }}(\alpha n,\beta n^{})\end{array}\right),`$ (A6)
and the potentials are defined as
$`\overline{\mathrm{\Gamma }}(\alpha n,\gamma n^{})`$ $`=`$ $`F_{\alpha \gamma \delta \beta }[V_i^{}(\beta n^{})V_i(\delta n)]`$ (A7)
$`=`$ $`F_{\alpha \gamma \delta \beta }_{11}(\delta n,\beta n^{}),`$ (A8)
$`\underset{¯}{\mathrm{\Gamma }}(\alpha n,\gamma n^{})`$ $`=`$ $`F_{\overline{\delta }\overline{\beta }\overline{\alpha }\overline{\gamma }}[U_i^{}(\beta n^{})U_i(\delta n)]`$ (A9)
$`=`$ $`F_{\overline{\delta }\overline{\beta }\overline{\alpha }\overline{\gamma }}\mathrm{tr}{\displaystyle \frac{1}{2}}(1\tau _3)_{22}(\delta n,\beta n^{}),`$ (A10)
$`\mathrm{\Delta }(\alpha n,\gamma n^{})`$ $`=`$ $`G_{\alpha \overline{\gamma }\beta \overline{\delta }}[U_i^{}(\delta n^{})V_i(\beta n)]`$ (A11)
$`=`$ $`G_{\alpha \overline{\gamma }\beta \overline{\delta }}_{12}(\beta n,\delta n^{}).`$ (A12)
Here we have utilized a generalized density matrix, $`(\alpha n,\beta n^{})`$, defined as
$$(\alpha n,\beta n^{})=\mathrm{\Psi }_i(\alpha n)\mathrm{\Psi }_i^{}(\beta n^{})=\left(\begin{array}{cc}_{11}& _{12}\\ _{21}& _{22}\end{array}\right),$$
(A13)
satisfying the idempotent condition
$$^2=\mathrm{\Omega }.$$
(A14)
The structure displayed in the preceeding paragraph suggests that if we knew the potentials $`\overline{\mathrm{\Gamma }},\underset{¯}{\mathrm{\Gamma }},\mathrm{\Delta }`$, the chemical potential $`\lambda `$, and the excitation energies $`E_{\overline{n}}^{},E_{\underset{¯}{n}}^{}`$, we could view Eq. (A2) as a linear eigenvalue problem, with given Hamiltonian $``$, with eigenvalues $`_i`$ and with solutions normalized according to Eq. (33). In turn, this suggests at least the first elements of an iterative scheme for the solution of Eq. (A2), where given the $`\nu th`$ approximation, $`^{(\nu )}`$ to $``$, we define the $`(\nu +1)st`$ approximation as the solution of the linear eigenvalue equation
$$_i^{(\nu +1)}\mathrm{\Psi }_i^{(\nu +1)}=^{(\nu )}\mathrm{\Psi }_i^{(\nu +1)}.$$
(A15)
Thus, with the help of Eq. (A13), we can construct a (tentative, see below) $`(\nu +1)st`$ approximation, $`^{(\nu +1)}`$, to the density matrix. Consequently, we can calculate a next approximation to the “potentials”.
But how do we generate higher approximations to the energies, including $`\lambda `$, of the neighboring even nuclei also needed for the next approximation to $``$? Actually the formalism provides several alternatives on how to proceed at this point. We describe first a method which adds the minimum number of conditions sufficient to determine the energies. This is to calculate, e. g.,
$`E_{\overline{n}}`$ $`=`$ $`\overline{n}|H|\overline{n}E_{\overline{n}}(^{(\nu +1)})=E_{\overline{n}}^{(\nu +1)},`$ (A16)
$`E_{\underset{¯}{n}}`$ $`=`$ $`\underset{¯}{n}|H|\underset{¯}{n}E_{\underset{¯}{n}}(^{(\nu +1)})=E_{\underset{¯}{n}}^{(\nu +1)},`$ (A17)
from which we can form the averages and differences needed for the calculation of $`^{(\nu +1)}`$.
Before we can actually proceed to the next iteration, we must take another step first recognized in our early work on the pairing problem . The general theory requires that if we are to interpret the diagonal elements of $`H`$ as eigenvalues, then the off-diagonal elements of $`H`$ must vanish. The latter may be calculated at any stage of approximations by the same techniques used for the diagonal elements. In general, we shall find at any intermediate stage of the calculation that these conditions are not fully satisfied. To rectify this deficiency, we must therefore carry out orthogonal transformations in the spaces $`|\overline{n}`$ and $`|\underset{¯}{n}`$ in order to eliminate the off-diagonal elements. These are equivalent to (different) linear transformations of the CFP $`V`$ and $`U`$. Such transformations are norm preserving.
We have now (almost) defined a cycle of the present algorithm. To achieve a full measure of self-consistency requires yet an additional link in the chain of reasoning. This is because to this point, the vectors $`\mathrm{\Psi }_i`$ cannot be guaranteed to satisfy the normalization conditions (35). In general we shall find
$$\underset{m_an}{}\mathrm{\Psi }_i^{}(\alpha n)\mathrm{\Psi }_i(\alpha n)=\mathrm{\Lambda }_a\mathrm{\Omega }_a,$$
(A18)
whereas full self-consistency requires $`\mathrm{\Lambda }_a=1`$. We can rectify the deficiency by rescaling the solution
$$\mathrm{\Psi }_i^{(\nu )}(\alpha n)\sqrt{\mathrm{\Lambda }_a}\mathrm{\Psi }_i^{(\nu )}(\alpha n),$$
(A19)
redefining the potentials, e. g.,
$$\overline{\mathrm{\Gamma }}(\alpha n,\gamma n^{})=F_{\alpha \gamma \delta \beta }\sqrt{\mathrm{\Lambda }_b\mathrm{\Lambda }_d}[V_i^{(\nu )}(\beta n^{})V_i^{(\nu )}(\delta n)],$$
(A20)
and similarly for other potentials. This requires us to fit in an additional set of iterations until the values $`\mathrm{\Lambda }_a=1`$ are achieved.
We next describe an alternative algorithm in which the equations of motion and normalization conditions are used in the same way as in the algorithm just described. However the transition from $`^{(\nu )}`$ to $`^{(\nu +1)}`$ is done differently and requires further development of the formalism. The first step is to combine the equations of motion (A2) with their complex conjugate equations so as to eliminate the eigenvalues $`_i`$. We thereby obtain the following equations for the generalized density matrix $``$,
$$0=(\alpha n,\gamma n^{\prime \prime })(\gamma n^{\prime \prime },\beta n^{})(\alpha n,\gamma n^{\prime \prime })(\gamma n^{\prime \prime },\beta n^{}),$$
(A21)
i. e., we find the vanishing of the commutator, $`[,]=0`$.
Before undertaking the exposition of the algorithm, we exhibit an alternative derivation of Eq. (A21) utilizing a variant of the variational principle (26), (26). Consider the functional
$$𝒟=𝒢\mathrm{\Theta }(\alpha n,\beta n^{})[^2(\beta n^{},\alpha n)\mathrm{\Omega }(\beta n^{},\alpha n)],$$
(A22)
where the “new” constraint with Lagrange multiplier matrix $`\mathrm{\Theta }`$ is for normalization in the density matrix form. Since
$$\frac{\delta 𝒢}{\delta (\beta n^{},\alpha n)}=(\alpha n,\beta n^{}),$$
(A23)
it follows that the variational condition applied to Eq. (A22) yields the equation
$$\mathrm{\Theta }\mathrm{\Theta }+\mathrm{\Theta }\mathrm{\Omega }=0.$$
(A24)
From this condition, Eq. (A21) is readily derived by forming the appropriate commutator.
However, Eq. (A24) contains additional information that we shall exploit below. Indeed, we shall consider the possibility of constructing an algorithm on the basis of Eq. (A24), but first we describe one that utilizes Eq. (A21) for the density matrix. Recalling the first algorithm described, let us imagine ourselves at the point where we have an approximation to the density matrix $``$ that has been determined from the iterative procedure associated with the equations of motion (A2), and that there remains the problem of computing the next approximation to the energies of the even neighbors. An alternative to the procedure that starts with Eq. (A16) is to note that the equation of motion in the density matrix form (A21) provides a sufficient set of equations to determine the excitation energies when supplemented by the number conservation conditions
$`0`$ $`=`$ $`{\displaystyle \underset{\alpha i}{}}|V_i(\alpha n)|^2\overline{N}`$ (A25)
$`=`$ $`{\displaystyle \underset{\alpha }{}}_{11}(\alpha n,\alpha n)\overline{N}`$ (A26)
$`0`$ $`=`$ $`{\displaystyle \underset{\alpha i}{}}|U_i(\alpha n)|^2\mathrm{\Omega }+\underset{¯}{N}`$ (A27)
$`=`$ $`{\displaystyle \underset{\alpha }{}}_{22}(\alpha n,\alpha n)\mathrm{\Omega }+\underset{¯}{N}`$ (A28)
Just as we did previously for the Hamiltonian, we should check that the number operator matrices $`_\alpha _{11}(\alpha n,\alpha n^{})`$ and $`_\alpha _{22}(\alpha n,\alpha n^{})`$ are diagonal and, if necessary, diagonalize them so as to improve the convergence of the procedure.
So far we have suggested two possible algorithms which used the equations of motion (A2), attached in each case to a different method for calculating the energies of the even nuclei. We next propose two additional algorithms, also distinguished by one of the two methods of computing the energy in which, however, we replace the equations of motion for the $`\mathrm{\Psi }`$ vectors by a condition related to the density-matrix formulation. Starting from Eq. (A24) we derive the pair
$``$ $`=`$ $`\mathrm{\Theta }=,`$ (A29)
which imply that
$$=\mathrm{\Theta }\mathrm{\Omega }.$$
(A30)
Substituting this result into Eq. (A24), we obtain
$$\frac{1}{\mathrm{\Omega }}\{,\}=0.$$
(A31)
We describe an algorithm based on Eq. (A31) and the normalization condition $`^2=\mathrm{\Omega }`$. Suppose that in a $`\nu th`$ approximation, we have
$`(^{(\nu ))})^2\mathrm{\Omega }^{(\nu )}`$ $`=`$ $`0,`$ (A32)
$`^{(\nu )}{\displaystyle \frac{1}{\mathrm{\Omega }}}\{^{(\nu )},^{(\nu )}\}`$ $`=`$ $`𝒴^{(\nu )}.`$ (A33)
As an illustration of the method of steepest descents, we now choose, in order to improve the value of the density matrix
$$\delta ^{(\nu )}=\eta 𝒴^{(\nu )},$$
(A34)
where $`\eta `$ is an arbitrary small parameter that may be set to unity for $`𝒴^{(\nu )}`$ sufficiently small. We can check that this choice preserves the norm to first order, since it satisfies the required condition
$$\{^{(\nu )},\delta ^{(\nu )}\}\mathrm{\Omega }\delta ^{(\nu )}=0,$$
(A35)
that follows from Eq. (A32). To see this, note that Eq. (A35) implies that
$$^{(\nu )}\delta ^{(\nu )}^{(\nu )}=0,$$
(A36)
which besides the trivial solution, is satisfied by Eq. (A34), as follows from Eq. (A33).
This as far as we can go without attacking specific models.
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# Geometric Algebra in Quantum Information Processing
## 1. Introduction
Quantum mechanics attaches physical significance to representations of the rotation group which differ substantially from those studied in classical geometry. Much of the mystery surrounding it is due to this fact. The enormous interest recently generated by proposals to build a *quantum computer* \[Llo95, EJ96, Ste98, WC98, Bro99, BD00\] has focussed attention on the simplest possible quantum system: a two-state system or *qubit*. Our understanding of qubits is based on two distinct geometric models of their states and transformations:
* A complex projective line under the action of $`\mathrm{𝖲𝖴}(2)`$.
* A Euclidean unit $`2`$-sphere under the action of $`\mathrm{𝖲𝖮}(3)`$.
The first is used almost exclusively in fundamental quantum physics, while the second (“classical”) model is used extensively in certain applications, e.g. nuclear magnetic resonance (NMR) spectroscopy \[HSTC00\]. In particular, in quantum computing a qubit represents a binary $`0`$ or $`1`$ as its state corresponds to one of a pair of conjugate (orthogonal) projective points $`|\mathrm{\hspace{0.17em}0}`$ or $`|\mathrm{\hspace{0.17em}1}`$. These in turn correspond to a pair of diametrical points on the unit sphere, which determine the alignment of the qubit with or against the corresponding axis of quantization.
Formally, these two models are related by stereographic projection of the Riemann (unit) sphere onto the Argand plane, the points of which are the ratios of the homogeneous coordinates of points on the projective line (see e.g. \[Alt86, FH81\]). While this elegant construction describes the mapping between the two representations in a geometric fashion, it does not unite them in a single mathematical structure. This paper provides an informal account of how this is done by geometric (aka Clifford) algebra; in addition, it describes an extension of this formalism to multi-qubit systems, and shows that it provides a concise and lucid means of describing the operations of quantum information processing \[SCH98, HSTC00\]. Significantly, this extension is most naturally derived via the geometric algebra of Minkowski space-time \[DLG93\], which has also been shown to be an efficient formalism within which to study a very wide range of problems in classical \[Hes99, Jan89\], relativistic \[Hes66, Bay96\] and fundamental quantum \[DLG<sup>+</sup>96\] physics. More complete and rigorous accounts may be found in these references, and in \[HCST00, Hav01\].
## 2. Euclidean Geometry and Spinors
Let $`_3`$ be a three-dimensional Euclidean vector space whose inner product is denoted by $`(𝒂,𝒃)𝒂𝒃`$. The Clifford or *geometric algebra* of $`_3`$ is the associative algebra generated by $`_3`$ over $``$ such that $`𝒂^2=𝒂^2𝒂𝒂`$ for all $`𝒂_3`$. This algebra will be referred to in the following as the *Pauli algebra*, and denoted by $`𝒢_3`$. The interesting thing about this algebra is its geometric interpretation, which will now be described.
To begin, note that every nonzero vector $`𝒂_3`$ has an inverse $`𝒂/𝒂^2`$. In addition, a simple application of the law of cosines shows that the inner product of $`𝒂`$ with any other vector $`𝒃_3`$ is given by the symmetric part of the geometric product:
(2.1)
$$\begin{array}{cc}\hfill \frac{1}{2}(𝒂𝒃+𝒃𝒂)=& \frac{1}{2}((𝒂+𝒃)^2𝒂^2𝒃^2)\hfill \\ \hfill =& \frac{1}{2}(𝒂+𝒃^2𝒂^2𝒃^2)=𝒂𝒃\hfill \end{array}$$
The antisymmetric part, by way of contrast, is called the *outer product*, and denoted by $`(𝒂,𝒃)𝒂𝒃(𝒂𝒃𝒃𝒂)/2`$. Since the outer product of two vectors $`𝒂𝒃`$ is invariant under inversion in the origin, it cannot itself be a vector. The space $`𝒂𝒃|𝒂,𝒃_3`$ therefore carries an inequivalent representation of the orthogonal group $`𝖮(3)`$, and its elements are accordingly called *bivectors*. These are most naturally interpreted as oriented plane segments, instead of oriented line segments like vectors in $`_3`$. If we similarly define the outer product of a vector with a bivector and require it to be associative, i.e.
(2.2)
$$𝒂(𝒃𝒄)\frac{1}{2}(𝒂𝒃𝒄𝒄𝒃𝒂)(𝒂𝒃)𝒄$$
($`𝒂,𝒃,𝒄_3`$), then it can be shown via straightforward though somewhat lengthy calculations that this product of three vectors is totally antisymmetric, meaning that the outer product generates the well-known exterior algebra $`_3`$ (cf. \[Hes99, Rie58\]). The outer product of three vectors is called a *trivector*, and (since it changes sign under inversion) is most appropriately interpreted as an oriented space segment or volume element.
The general properties of inner and outer products in the geometric algebras of arbitrary metric vector spaces can be worked out along these lines in a coordinate-free fashion \[HS84\]. The remainder of this section will focus on how the Pauli algebra is used to describe the quantum mechanics of qubits. In this application it is more common to work with a fixed orthonormal basis $`𝝈_𝗑,𝝈_𝗒,𝝈_𝗓_3`$. Quantum mechanics, however, views these basis vectors in a very different way from that taken above, in that they are regarded as *operators* on a two-dimensional Hilbert space $`^2`$ (see e.g. \[Sak94\]). These operators, in turn, are usually identified with the Pauli matrices
(2.3)
$$𝝈_𝗑\underset{¯}{𝝈}_𝗑\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right],𝝈_𝗒\underset{¯}{𝝈}_𝗒\left[\begin{array}{cc}0& ı\\ ı& 0\end{array}\right],𝝈_𝗓\underset{¯}{𝝈}_𝗓\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right],$$
where $`ı`$ is an imaginary unit ($`ı^2=1`$), the underline signifies that the associated symbol is a matrix, and throughout this paper the symbol “$``$” should be read as “is represented by” or “is equivalent to”. The connection between the two viewpoints lies the fact that these matrices satisfy the defining relations of the *abstract* Pauli algebra $`𝒢_3`$, namely
(2.4)
$$(𝝈_\mu )^2=\mathrm{𝟏}1,𝝈_\mu 𝝈_\nu =𝝈_\nu 𝝈_\mu (\mu ,\nu \{𝗑,𝗒,𝗓\},\mu \nu ),$$
and hence constitute a faithful matrix representation of it. This shows, in particular, that $`𝒢_3`$ is $`8`$-dimensional as a *real* linear space.<sup>1</sup><sup>1</sup>1 In the quantum mechanics literature, the notation $`\underset{¯}{𝒂}\underset{¯}{\overset{}{𝝈}}`$ is often used for $`_\mu a_\mu \underset{¯}{𝝈}_\mu `$. Because $`𝒂`$ is geometrically just a vector in $`_3`$ (*not* a matrix for it in a basis-dependent representation of $`𝒢_3`$), this is an abuse of the dot-notation for the Euclidean inner product, which is otherwise perhaps the most consistently used notation in all of science. This abuse of notation will not be perpetrated in this paper.
In most physical situations, these operators (times $`\mathrm{}`$) represent measurements of the intrinsic angular momentum of the qubits, and hence are regarded as generators of rotations in the Lie algebra $`\mathrm{𝗌𝗈}(3)`$ over $``$ satisfying the commutator relation
(2.5)
$$\frac{1}{2}[𝝈_𝗑,𝝈_𝗒]=ı𝝈_𝗓,$$
and its cyclic permutations. In terms of geometric algebra, the left-hand side is just the outer product of the vectors. The right-hand side is somewhat harder to interpret, because the Pauli algebra is defined over the real numbers. The trick is to observe that, in terms of the matrix representation, $`\underset{¯}{𝝈}_𝗑\underset{¯}{𝝈}_𝗒\underset{¯}{𝝈}_𝗓=ı\underset{¯}{\mathrm{𝟏}}`$. Thus by interpreting the abstract imaginary $`ı`$ as the trivector $`𝜾𝝈_𝗑𝝈_𝗒𝝈_𝗓=𝝈_𝗑𝝈_𝗒𝝈_𝗓`$, the angular momentum relations become a triviality:
(2.6)
$$𝝈_𝗑𝝈_𝗒=𝝈_𝗑𝝈_𝗒=𝝈_𝗑𝝈_𝗒(𝝈_𝗓)^2=𝜾𝝈_𝗓$$
More generally, the vector cross product is related to the outer product by
(2.7)
$$𝒂\times 𝒃=\frac{𝜾}{2}(𝒂𝒃𝒃𝒂)=𝜾(𝒂𝒃),$$
from which it may be seen that multiplication by the unit trivector $`𝜾`$ maps vectors to orthogonal bivectors and vice versa. Since they span a one-dimensional space but change sign under inversion in the origin, trivectors can also be regarded as *pseudo-scalars*. Perhaps the most important thing which geometric algebra contributes to physics are geometric interpretations for the imaginary units which it otherwise uses blindly.
If we denote the induced bivector basis by
(2.8)
$$𝑰𝝈_𝗒𝝈_𝗓,𝑱𝝈_𝗓𝝈_𝗑,𝑲𝝈_𝗑𝝈_𝗒,$$
it is readily seen that these basis bivectors likewise square to $`1`$. On multiplying the angular momentum generating relations through by $`1=𝜾^2`$, we obtain
(2.9)
$$𝑱𝑰=𝑲,𝑰𝑲=𝑱,𝑲𝑱=𝑰,\text{and}𝑲𝑱𝑰=1.$$
This shows that these basis bivectors generate a subalgebra of $`𝒢_3`$ isomorphic to Hamilton’s quaternions \[Alt86, Alt89\], which is also known as the *even subalgebra* $`𝒢_3^+`$ (since it is generated by the products of even numbers of vectors). It is well-known that the quaternions’ multiplicative group is $`^{}\mathrm{𝖲𝖴}(2)`$, which implies that the even subalgebra should be closely related to rotations. This relationship will now be worked out explicitly.
Consider the result of conjugating a vector $`𝒙`$ by a vector $`𝒂`$, i.e.
(2.10)
$$𝒂𝒙𝒂^1=𝒂(𝒙_{}+𝒙_{})𝒂^1=𝒂𝒂^1𝒙_{}𝒂𝒂^1𝒙_{},$$
where we have split $`𝒙=(𝒙𝒂^1+𝒙𝒂^1)𝒂𝒙_{}+𝒙_{}`$ into its parts parallel and perpendicular to $`𝒂`$. This shows that $`𝒂𝒙𝒂^1`$ is the *reflection* of $`𝒙`$ in the plane orthogonal to $`𝒂`$. From the well-known fact that the composition of two reflections is a rotation by *twice* the lessor angle between their planes and about these planes’ line of intersection, it follows that conjugating a vector by an element of the even subalgebra just rotates it accordingly:
(2.11)
$$(𝒃𝒂)𝒙(𝒃𝒂)^1=𝒃𝒂𝒙𝒂^1𝒃^1=\frac{𝒃𝒂𝒙𝒂𝒃}{𝒂^2𝒃^2}$$
Let $`𝒖𝒂/𝒂`$, $`𝒗𝒃/𝒃`$ and $`𝑹𝒗𝒖`$ be the corresponding *unit* quaternion. Then $`𝑹=\mathrm{cos}(\theta /2)𝜾𝒓\mathrm{sin}(\theta /2)`$ where $`\mathrm{cos}(\theta /2)=𝒖𝒗`$ and $`𝜾𝒓𝒖𝒗/𝒖𝒗`$. Moreover, the inverse $`(𝒗𝒖)^1`$ is now simply the *reverse* $`𝒖𝒗(𝒗𝒖)^{}`$, which in turn corresponds to the *conjugate* quaternion $`𝑹^{}\mathrm{cos}(\theta /2)+𝜾𝒓\mathrm{sin}(\theta /2)`$. This reversal operation on $`𝒢_3^+`$ extends to a well-defined anti-automorphism of $`𝒢_3`$, which corresponds to Hermitian conjugation in its representation by Pauli matrices. On splitting $`𝒙`$ into its parts parallel $`𝒙_{}`$ and perpendicular $`𝒙_{}`$ to $`𝒓`$ as above, the rotation may now be written as $`𝑹𝒙𝑹^{}=𝑹(𝒙_{}+𝒙_{})𝑹^{}=𝒙_{}+𝒙_{}(𝑹^{})\text{}^2`$
(2.12)
$$\begin{array}{cc}\hfill =& 𝒙_{}+𝒙_{}(\mathrm{cos}^2(\theta /2)\mathrm{sin}^2(\theta /2)+2𝜾𝒓\mathrm{cos}(\theta /2)\mathrm{sin}(\theta /2))\hfill \\ \hfill =& 𝒙_{}+𝒙_{}(\mathrm{cos}(\theta )+𝜾𝒓\mathrm{sin}(\theta )),\hfill \end{array}$$
and so may be viewed as multiplication of $`𝒙_{}`$ by the “complex number” $`\mathrm{cos}(\theta )+𝜾𝒓\mathrm{sin}(\theta )`$ in the Argand plane defined by the bivector $`𝜾𝒓`$.
By collecting even and odd powers in its Taylor series, it may be seen that any unit quaternion can be written as the exponential of a bivector orthogonal to the axis of rotation $`𝒓`$:
(2.13)
$$e^{𝜾𝒓\theta /2}=1𝜾𝒓\frac{\theta }{2}\frac{1}{2}\left(\frac{\theta }{2}\right)^2+\mathrm{}=\mathrm{cos}(\theta /2)𝜾𝒓\mathrm{sin}(\theta /2)$$
This is formally analogous to a complex exponential, and is also in accord with our previous observation that the space of bivectors is isomorphic to the Lie algebra $`\mathrm{𝗌𝗎}(2)\mathrm{𝗌𝗈}(3)`$ under the commutator product. The pair $`[\mathrm{cos}(\theta /2);\mathrm{sin}(\theta /2)𝒓]`$ are known as *Euler-Rodrigues* parameters for the rotation group $`\mathrm{𝖲𝖮}(3)`$; since $`[\mathrm{cos}(\theta /2);\mathrm{sin}(\theta /2)𝒓]`$ determines the same rotation, this parametrization is two-to-one. A one-to-one parametrization is obtained from the *outer exponential*, i.e.
(2.14)
$$^{𝜾𝒓\tau }=1𝜾𝒓\tau \frac{1}{2}𝒓𝒓\tau ^2+\mathrm{}=1𝜾𝒓\tau \text{(since }𝒓𝒓=0\text{)}.$$
The squared norm of this outer exponential is $`1+\tau ^2`$, so that the normalized outer exponential equals the usual exponential if we set $`\tau =\mathrm{tan}(\theta /2)`$. Because $`𝒕\mathrm{tan}(\theta /2)𝒓`$ is the four-dimensional stereographic projection of $`[\mathrm{cos}(\theta );\mathrm{sin}(\theta )𝒓]`$ from $`\theta =\pi `$, it has been called the *stereographic* parameter for $`\mathrm{𝖲𝖮}(3)`$. Note, however, that this parametrization does not include rotations by $`\pi `$.
Another two-to-one parametrization of rotations is given by the *Cayley-Klein* parameters $`[\psi _1;\psi _2]^2`$, where
(2.15)
$$\begin{array}{cc}\hfill \psi _1& \mathrm{cos}(\theta /2)ı\mathrm{sin}(\theta /2)(𝒓𝝈_𝗓),\hfill \\ \hfill \psi _2& \mathrm{sin}(\theta /2)(𝒓𝝈_𝗑)+ı\mathrm{sin}(\theta /2)(𝒓𝝈_𝗒).\hfill \end{array}$$
The corresponding $`\mathrm{𝖲𝖴}(2)`$ matrix is simply
(2.16)
$$\underset{¯}{𝚿}=\left[\begin{array}{cc}\psi _1& \psi _2^{}\\ \psi _2& \psi _1^{}\end{array}\right].$$
It follows that the complex column vector $`|\psi [\psi _1;\psi _2]`$ itself transforms under left-multiplication with matrices in $`\mathrm{𝖲𝖴}(2)`$, which is commonly described in quantum mechanics by calling it a *spinor*. In particular, the spinors $`|\mathrm{\hspace{0.17em}0}[1;0]`$ and $`|\mathrm{\hspace{0.17em}1}[0;1]`$ are those commonly used in quantum computing to store binary information. Since the Cayley-Klein parameters uniquely determine the $`\mathrm{𝖲𝖴}(2)`$ matrix, however, we can just as well regard spinors as entities *in* $`\mathrm{𝖲𝖴}(2)`$, e.g. $`|\mathrm{\hspace{0.17em}0}\underset{¯}{\mathrm{𝟏}}`$ and $`|\mathrm{\hspace{0.17em}1}ı\underset{¯}{𝝈}_𝗒`$. The usual action of $`\mathrm{𝖲𝖴}(2)`$ on spinors then becomes the left-regular action of $`\mathrm{𝖲𝖴}(2)`$ on itself.
The representation of $`\mathrm{𝖲𝖴}(2)`$ used above depends upon the choice of coordinate system: Changing to a different the coordinate system gives a different (though equivalent) representation. Recalling that $`\mathrm{𝖲𝖴}(2)`$ is isomorphic to the multiplicative group of unit elements (quaternions) in the even subalgebra $`𝒢_3^+`$, a coordinate-free or *geometric* interpretation of spinors is obtained by regarding them as elements of $`𝒢_3^+`$ itself. This interpretation of spinors as entities in ordinary Euclidean geometry was first pointed out by Hestenes over thirty years ago \[Hes66\], but physicists persist in putting operators and operands into separate spaces, and in working with a matrix representation instead of directly with the geometric entities themselves. The perceived nonintuitive nature of quantum mechanics is due in large part to the resulting confusion over the geometric meaning of the objects with which it deals, which is spelled out explicitly in geometric algebra.
As another example, consider how the *density operator* of an “ensemble” of qubits can be interpreted in geometric algebra. This operator $`𝝆`$ is usually defined via a matrix representation as $`\underset{¯}{𝝆}\overline{|\psi \psi |}`$, where the overline denotes the average over the ensemble. As first observed by von Neumann, this matrix contains all the information needed to compute the ensemble average expectation values of the qubit observables, since
(2.17)
$$\overline{\psi |\underset{¯}{𝝈}_\mu |\psi }=\overline{\mathrm{tr}(\underset{¯}{𝝈}_\mu |\psi \psi |)}=\mathrm{tr}(\underset{¯}{𝝈}_\mu \overline{|\psi \psi |})=\mathrm{tr}(\underset{¯}{𝝈}_\mu \underset{¯}{𝝆})$$
($`\mu \{𝗑,𝗒,𝗓\}`$). To translate this into geometric algebra, we set the second column of $`\underset{¯}{𝚿}`$ to zero by right-multiplying it by the idempotent matrix $`\underset{¯}{𝑬}_+(\underset{¯}{\mathrm{𝟏}}+\underset{¯}{𝝈}_𝗓)/2`$, i.e.
(2.18)
$$\underset{¯}{𝚿}\underset{¯}{𝑬}_+\left[\begin{array}{cc}\psi _1& \psi _2^{}\\ \psi _2& \psi _1^{}\end{array}\right]\left[\begin{array}{cc}1& 0\\ 0& 0\end{array}\right]=\left[\begin{array}{cc}\psi _1& 0\\ \psi _2& 0\end{array}\right].$$
This corresponds to projecting $`𝚿𝒢_3^+`$ onto a *left-ideal* in $`𝒢_3`$, and allows the dyadic product $`|\psi \psi |`$ in Eq. (2.17) to be written as:
(2.19)
$$|\psi \psi |=\left[\begin{array}{cc}\psi _1& 0\\ \psi _2& 0\end{array}\right]\left[\begin{array}{cc}\psi _1^{}& \psi _2^{}\\ 0& 0\end{array}\right](\underset{¯}{𝚿}\underset{¯}{𝑬}_+)(\underset{¯}{𝚿}\underset{¯}{𝑬}_+)^{}$$
Thus the interpretation of the density operator in geometric algebra is
(2.20)
$$𝝆=\overline{(𝚿𝑬_+)(𝚿𝑬_+)^{}}=\overline{𝚿𝑬_+𝚿^{}}=\frac{1}{2}\left(1+\overline{𝚿𝝈_𝗓𝚿^{}}\right)$$
(cf. \[SLD99\]). The vector part $`𝒑𝝆_1=\overline{𝚿𝝈_𝗓𝚿^{}}`$ is called the *polarization vector* (in optics, its components are known as the *Stokes parameters* \[BBDH93\], while in NMR it is known as the *Bloch vector* after the pioneer of NMR who rediscovered it \[Blo46\]). Its length is $`𝒑1`$ with equality if and only if all members of the ensemble are in the same state $`𝚿`$. In this case the ensemble is said to be in a *pure* state, and the density operator is itself an idempotent $`(1+𝒑)/2`$, where $`𝒑𝚿𝝈_𝗓𝚿^{}`$. For an ensemble in a general *mixed* state, the length of the ensemble-average polarization vector measures the degree of alignment among the (unit length) polarization vectors of the individual members of the ensemble, and is called the *polarization* of the ensemble.
In many physical situations there is a natural reference direction; for example, in NMR computing the qubits are spin $`1/2`$ atomic nuclei whose intrinsic magnetic dipoles have been polarized by the application of a strong magnetic field \[HCST00\]. From a geometric perspective, however, the density operator is just the sum of a scalar and a vector, which for a pure state is related to the corresponding “spinor” by rotation of a fixed reference vector (conventionally taken to be $`𝝈_𝗓`$ as above) by $`𝚿`$. Since the trace in the standard matrix representation is simply twice the scalar part $`\underset{¯}{}_0`$ of the corresponding expression in geometric algebra, the ensemble-average expectation value
(2.21)
$$\frac{1}{2}\mathrm{tr}(\underset{¯}{𝝈}_\mu \underset{¯}{𝝆})𝝈_\mu 𝝆_0=𝝈_\mu 𝚿𝝈_𝗓𝚿^{}_0=𝝈_\mu (𝚿𝝈_𝗓𝚿^{})=𝝈_\mu 𝒑$$
is just the component of the polarization vector along the $`\mu `$-th axis. Unlike the *strong* measurements usually considered in quantum texts, where measurement of $`𝝈_\mu `$ yields one of the random outcomes $`\pm 1`$ with probabilities $`(1\pm 𝝈_\mu 𝒑)/2`$ and leaves the system in the corresponding state $`𝒑=\pm 𝝈_\mu `$, *weak* measurements of ensemble-average expectation values can be made with only negligible perturbations to the ensemble as a whole \[Per93\]. This is in fact how quantum mechanical systems are usually manifest at the macroscopic level!
To see how all this relates to conventional wisdom, observe that the polarization vector of a pure state may be written in terms of the Cayley-Klein parameters as
(2.22)
$$𝒑=2\mathrm{}(\psi _1^{}\psi _2)𝝈_𝗑+2\mathrm{}(\psi _1^{}\psi _2)𝝈_𝗒+(|\psi _1|^2|\psi _2|^2)𝝈_𝗓.$$
Its stereographic projection from $`𝝈_𝗓`$ onto the $`𝝈_𝗑𝝈_𝗒`$ plane is therefore
(2.23)
$$\frac{2\mathrm{}(\psi _1^{}\psi _2)𝝈_𝗑+2\mathrm{}(\psi _1^{}\psi _2)𝝈_𝗒}{1+|\psi _1|^2|\psi _2|^2}.$$
Multiplying by $`𝝈_𝗑`$ and simplifying the denominator using $`|\psi _1|^2+|\psi _2|^2=1`$ yields
(2.24)
$$\frac{\mathrm{}(\psi _1^{}\psi _2)+\mathrm{}(\psi _1^{}\psi _2)𝑲}{|\psi _1|^2}$$
where $`𝑲=𝝈_𝗑𝝈_𝗒`$ is a square-root of $`1`$. This is the same as the ratio $`\psi _2/\psi _1`$ save for the use of $`ı`$ instead of $`𝑲`$ as the imaginary unit, which explains formally why $`\mathrm{𝖲𝖮}(3)`$ acts on the polarization vector in the same way that $`\mathrm{𝖲𝖴}(2)`$ acts on the ratio of the Cayley-Klein parameters \[Alt86, FH81\].
## 3. Space-Time Geometry and Multiparticle Spinors
The above interpretations apply only to single qubits (or to ensembles consisting of noninteracting and identical qubits). Extending them to systems of interacting and distinguishable qubits may be done in a physically significant fashion by considering the geometric algebra of *space-time* (or Minkowski space) $`_{1,3}`$. This algebra, known as the *Dirac* algebra and denoted by $`𝒢_{1,3}`$, may be defined by the generating relations among an orthonormal basis analogous to Eq. (2.4):
(3.1)
$$\begin{array}{cc}& 𝜸_𝗍^2=1,𝜸_\mu ^2=1(\mu \{𝗑,𝗒,𝗓\}),\hfill \\ & 𝜸_\mu 𝜸_\nu =𝜸_\nu 𝜸_\mu (\mu ,\nu \{𝗍,𝗑,𝗒,𝗓\},\mu \nu )\hfill \end{array}$$
The corresponding geometric algebra separates into five inequivalent representations under the action of the full Lorentz group $`𝖮(1,3)`$, i.e.
(3.2)
$$\begin{array}{cc}\hfill 1& (\text{scalars, }1\text{-dimensional})\hfill \\ \hfill 𝜸_\mu & (\text{vectors, }4\text{-dimensional})\hfill \\ \hfill 𝜸_\mu 𝜸_\nu & (\text{bivectors, }6\text{-dimensional})\hfill \\ \hfill 𝜸_\mu 𝜸_\nu 𝜸_\eta & (\text{trivectors, }4\text{-dimensional})\hfill \\ \hfill 𝜸_𝗍𝜸_𝗑𝜸_𝗒𝜸_𝗓& (\text{pseudo-scalars, }1\text{-dimensional}),\hfill \end{array}$$
where $`\mu ,\nu ,\eta \{𝗍,𝗑,𝗒,𝗓\}`$ with $`\mu \nu \eta \mu `$, for a total dimension of $`16`$.
The important point for our purposes is that the even subalgebra of the Dirac algebra $`𝒢_{1,3}^+`$ is isomorphic to the Pauli algebra $`𝒢_3`$ \[Hes66\]. This isomorphism may be constructed by choosing bases $`𝜸_\mu 𝒢_{1,3}`$ and $`𝝈_\mu 𝒢_3`$, and defining an invertible linear mapping by
(3.3)
$$𝝈_\mu 𝒢_3𝜸_\mu 𝜸_𝗍𝒢_{1,3}^+(\mu \{𝗑,𝗒,𝗓\}).$$
These so-called *relative spatial vectors* $`𝜸_\mu 𝜸_𝗍`$ satisfy the relations in Eq. (2.4), since
(3.4)
$$\begin{array}{cc}\hfill (𝝈_\mu )^2& (𝜸_\mu 𝜸_𝗍)^2=𝜸_\mu (𝜸_𝗍)^2𝜸_\mu =(𝜸_\mu )^2=1\hfill \\ \hfill 𝝈_\mu 𝝈_\nu & (𝜸_\mu 𝜸_𝗍)(𝜸_\nu 𝜸_𝗍)=𝜸_\nu (𝜸_\mu 𝜸_𝗍)𝜸_𝗍=(𝜸_\nu 𝜸_𝗍)(𝜸_\mu 𝜸_𝗍)\hfill \\ \hfill & 𝝈_\nu 𝝈_\mu (\mu ,\nu \{𝗑,𝗒,𝗓\},\mu \nu ),\hfill \end{array}$$
and hence generate an algebra isomorphic to $`𝒢_3`$. As bivectors in $`𝒢_{1,3}`$, however, they also generate $`𝒢_{1,3}^+`$, since
(3.5)
$$𝜸_\mu 𝜸_\nu =𝜸_\mu (𝜸_𝗍)^2𝜸_\nu =(𝜸_\mu 𝜸_𝗍)(𝜸_\nu 𝜸_𝗍)𝝈_\mu 𝝈_\nu $$
($`\mu ,\nu \{𝗑,𝗒,𝗓\}`$, $`\mu \nu `$), and similarly
(3.6)
$$𝜾𝜸_𝗍𝜸_𝗑𝜸_𝗒𝜸_𝗓=(𝜸_𝗑𝜸_𝗍)(𝜸_𝗒𝜸_𝗍)(𝜸_𝗓𝜸_𝗍)𝝈_𝗑𝝈_𝗒𝝈_𝗓.$$
Thus $`𝜸_\mu 𝜸_𝗍𝝈_\mu `$ ($`\mu \{𝗑,𝗒,𝗓\}`$) induces an algebra isomorphism as claimed, and when the bases are understood we may identify $`𝝈_\mu 𝜸_\mu 𝜸_𝗍`$.
The choice of time-like vector $`𝜸_𝗍_{1,3}`$ in fact determines an inertial frame up to spatial rotation, in which the *time* $`t`$ and *place* $`𝒔`$ of an event $`𝐞`$ in that frame are given by
(3.7)
$$t+𝒔=𝐞𝜸_𝗍+𝐞𝜸_𝗍=𝐞𝜸_𝗍$$
(note that upright case is used for the space-time vector $`𝐞_{1,3}`$). Thus the invariant interval between events separated by the space-time vector $`𝐞`$ is $`𝐞^2=𝐞𝜸_𝗍^2𝐞=(t+𝒔)(t𝒔)=t^2𝒔^2`$ as usual, while the relative velocity between events whose space-time velocities are $`𝜸_𝗍`$ and $`𝐯𝐞/\tau `$ is
(3.8)
$$𝒗=\frac{𝒔}{t}=\frac{𝒔}{\tau }\frac{\tau }{t}\left(\frac{𝐞}{\tau }𝜸_𝗍\right)\left(\frac{𝐞}{\tau }𝜸_𝗍\right)^1=\frac{𝐯𝜸_𝗍}{𝐯𝜸_𝗍},$$
so that $`𝒗𝜸_𝗍`$ lies on an affine hyperplane in space-time.
A great deal of physics can be done in a manifestly Lorentz covariant fashion using the Dirac algebra. For example, the electromagnetic field at a given point in space-time corresponds to an arbitrary bivector $`𝐅_2_{1,3}`$, called the *Faraday bivector*, and the covariant form of the Lorentz force equation is
(3.9)
$$m\dot{𝐯}=q𝐅𝐯,$$
where $`m`$ is the rest mass, $`q`$ the charge and $`𝐯`$ the space-time velocity. (This is another example of the general rule that, in geometric algebra, the generators of motion are bivectors \[DHSvA93\].) The usual frame-dependent form is recovered by splitting the quantities in this equation by $`𝜸_𝗍`$ as above \[Jan89\]; in particular, the Faraday bivector splits into an electric and a magnetic field as $`𝐅𝑬+𝜾𝑩`$, where
(3.10)
$$𝑬=(𝐅𝜸_𝗍)𝜸_𝗍\text{and}𝜾𝑩=(𝐅𝜸_𝗍)𝜸_𝗍.$$
The space-time reverse will be denoted by a tilde, e.g. in the present case $`\stackrel{~}{𝐅}=𝐅`$. This is related to the spatial (or Pauli) reverse by $`𝐅^{}=𝑬𝜾𝑩=𝜸_𝗍\stackrel{~}{𝐅}_{}𝜸_𝗍`$. Both operations agree on the Pauli-even subalgebra, but the spatial reverse *not* Lorentz coveriant since it depends on a particular $`𝜸_𝗍`$.
Returning to our previous discussion of the density operator, we observe that the space-time form of the density operator of a single qubit polarized along $`𝗓`$ can be written as
(3.11)
$$𝝆=\frac{1}{2}(1+\alpha 𝝈_𝗓)=\frac{1}{2}(𝜸_𝗍+\alpha 𝜸_𝗓)𝜸_𝗍\mathit{\varrho }𝜸_𝗍,$$
where $`1\alpha 1`$ is the polarization and $`𝜸_𝗍`$ determines the local inertial frame. It follows that the Lorentz covariant form of the density operator is a time-like vector $`\mathit{\varrho }_{1,3}`$. Under a Lorentz boost $`𝑳=\mathrm{exp}(\lambda 𝝈_𝗓/2)\mathrm{𝖲𝖮}(1,3)`$ along $`𝝈_𝗓`$, therefore, the relativistic density operator $`\mathit{\varrho }`$ transforms to $`\mathit{\varrho }^{}\frac{1}{2}𝑳(𝜸_𝗍+\alpha 𝜸_𝗓)\stackrel{~}{𝑳}=`$
(3.12)
$$\frac{1}{2}\left(\mathrm{cosh}(\lambda )𝜸_𝗍\mathrm{sinh}(\lambda )𝜸_𝗓+\alpha (\mathrm{cosh}(\lambda )𝜸_𝗓\mathrm{sinh}(\lambda )𝜸_𝗍)\right).$$
This implies that in the unaccelerated frame (with renormalization by $`\mathit{\varrho }^{}𝜸_𝗍`$),
(3.13)
$$\begin{array}{cc}\hfill 𝝆^{}=& \frac{\mathit{\varrho }^{}𝜸_𝗍}{\mathit{\varrho }^{}𝜸_𝗍}=\frac{\mathit{\varrho }^{}𝜸_𝗍+\mathit{\varrho }^{}𝜸_𝗍}{\mathit{\varrho }^{}𝜸_𝗍}\hfill \\ \hfill =& \frac{1}{2}\left(1+\frac{\alpha \mathrm{cosh}(\lambda )\mathrm{sinh}(\lambda )}{\mathrm{cosh}(\lambda )\alpha \mathrm{sinh}(\lambda )}𝝈_𝗓\right).\hfill \end{array}$$
It follows that the polarization itself transforms as
(3.14)
$$\alpha ^{}=\frac{\alpha \mathrm{cosh}(\lambda )\mathrm{sinh}(\lambda )}{\mathrm{cosh}(\lambda )\alpha \mathrm{sinh}(\lambda )}.$$
If we assume the qubit is at equilibrium with a heat bath, statistical mechanics tells us that $`\alpha =\mathrm{tanh}(\beta ϵ/2)`$ where $`\beta =1/(k_𝖡T)`$ is the inverse temperature and $`ϵ`$ is the energy difference between the $`|\mathrm{\hspace{0.17em}0}`$ and $`|\mathrm{\hspace{0.17em}1}`$ states \[Tol38\]. Then the addition formulae for $`\mathrm{cosh}`$ and $`\mathrm{sinh}`$ give
(3.15)
$$\alpha ^{}=\mathrm{tanh}(\beta ϵ/2\lambda ),$$
so the apparent equilibrium polarization depends on velocity. These results are not to be found in the classic treatise on relativistic thermodynamics \[Tol34\].
We will now construct a Lorentz covariant *multiparticle theory* of qubit systems in the simplest possible way, by taking a direct sum of copies of space-time (regarded as a vector space, rather than an algebra), one for each of the $`N`$ qubits, i.e.
(3.16)
$$_{q=1}^N𝜸_𝗍^q,𝜸_𝗑^q,𝜸_𝗒^q,𝜸_𝗓^q,$$
and considering the associated geometric algebra $`𝒢_{N,3N}`$. Then the even subalgebras of different particle spaces $`pq`$ *commute*, since (in any given bases)
(3.17)
$$𝝈_\mu ^p𝝈_\nu ^q=𝜸_\mu ^p(𝜸_\nu ^q𝜸_𝗍^q)𝜸_𝗍^p=𝝈_\nu ^q𝝈_\mu ^p$$
for all $`\mu ,\nu \{𝗑,𝗒,𝗓\}`$, so that the algebra generated by the even subalgebras is isomorphic to a tensor product of these algebras, written as
(3.18)
$$(𝒢_{1,3}^+)^N𝒢_3^N(𝒢_3)^N.$$
This construction of the tensor product was first used by Clifford as a means of studying the tensor products of quaternion algebras \[Cli78\]; van der Waerden has in fact called it a Clifford algebra of the second kind \[vdW85\]. As a means of justifying the tensor product of nonrelativistic quantum mechanics in terms of the underlying geometry of space-time, however, it is a much more recent development \[DLG93\].
A key feature of quantum mechanics, which is needed for quantum computers to be able to solve problems more efficiently than their classical counterparts, is an exponential growth in the dimension of the Hilbert space of a multi-qubit system with the number of particles involved. The complex dimension of the Hilbert space $`()^N`$ of an $`N`$-qubit system is in fact $`2^N`$, and the space of operators (linear transformations) on $`()^N`$ therefore has *real* dimension $`2^{2N+1}`$. The above construction yields a space of “operators” $`𝒢_3^N`$ whose real dimension also grows exponentially, but as $`2^{3N}`$. The extra degrees of freedom are due to the presence of a different unit pseudo-scalar $`𝜾^q`$ in every particle space. They can easily be removed by multiplying through by an idempotent element called the *correlator*:
(3.19)
$$𝑪\frac{1}{2}(1𝜾^1𝜾^2)\frac{1}{2}(1𝜾^1𝜾^3)\mathrm{}\frac{1}{2}(1𝜾^1𝜾^N)$$
This commutes with everything in $`𝒢_3^N`$ and satisfies $`𝜾^p𝜾^q𝑪=𝑪`$ for $`1p,qN`$, so that multiplication by it homomorphically maps $`𝒢_3^N`$ onto an ideal $`𝒢_3^N/𝑪`$ wherein all the unit pseudo-scalars have been identified,<sup>2</sup><sup>2</sup>2The notation $`𝒢_3^N/𝑪`$ is justified by the fact that the two-sided principle ideal $`𝒢_3^N(𝑪)`$ generated by $`𝑪`$ is isomorphic to the quotient algebra $`𝒢_3^N/\mathrm{ker}(𝑪)`$, where $`\mathrm{ker}(𝑪)\{𝒈𝒢_3^N𝒈𝑪=0\}`$. and which therefore has the correct dimension over $``$. As a subalgebra, this ideal is in fact isomorphic to the algebra of $`2^N\times 2^N`$ complex matrices, and hence capable of describing all the states and transformations of (ensembles of) $`N`$ qubit systems. In the following, we shall generally omit $`𝑪`$ from our expressions altogether, and use a single unit imaginary $`𝜾`$ as in conventional quantum mechanics.
On the “even” subalgebra $`(𝒢_3^+)^N`$, multiplication by the correlator turns out to be an algebra automorphism; this algebra can thus be written as
(3.20)
$$(𝒢_3^+)^N(𝒢_3^+)^N/𝑪(𝒢_3^N/𝑪)^+𝒮𝒰(2)^N,$$
where the “$`+`$” refers throughout to the subalgebra generated by expressions which are invariant under inversion in the origin, and $`𝒮𝒰(2)\text{}^N`$ to the algebra generated over $``$ by the Kronecker products of matrices in the group $`\mathrm{𝖲𝖴}(2)`$. This subalgebra has real dimension $`2^{2N}`$, but is mapped onto a left-ideal of dimension $`2^{N+1}`$ by right-multiplication with another idempotent which is given by the tensor product of those considered earlier, namely
(3.21)
$$𝑬_+𝑬_+^1𝑬_+^2\mathrm{}𝑬_+^N,$$
where $`𝑬_\pm ^q(1\pm 𝝈_𝗓^q)/2`$ for $`q=1,\mathrm{},N`$. Henceforth, the term “even subalgebra” will refer to $`(𝒢_3^+)\text{}^N`$ (suitably correlated) unless otherwise stated.
In terms of the usual matrix representation, right-multiplication of an element of the even subalgebra $`\underset{¯}{𝚿}`$ by $`\underset{¯}{𝑬}_+`$ likewise sets all but the first column to zero, so that $`𝚿𝑬_+`$ transforms like a “spinor” in $`^N`$ under left-multiplication by single particle rotations $`𝑹^q(𝒢_3^+)\text{}^N`$. Unlike the single particle case, however, this one column does not uniquely determine an element of the even subalgebra $`(𝒢_3^+)\text{}^N/𝑪`$. What has been proposed instead \[DLG93\] is to use the fact that $`𝑬_+`$ “absorbs” $`𝝈_𝗓`$’s to distribute copies of the latter across the correlator, converting it to what will here be called the *directional* correlator $`𝑫`$, i.e.
(3.22)
$$𝚿𝑪𝑬_+=𝚿𝑪\left((𝝈_𝗓^{\mathrm{\hspace{0.17em}1}})^{N1}𝝈_𝗓^{\mathrm{\hspace{0.17em}2}}\mathrm{}𝝈_𝗓^N\right)𝑬_+=𝚿𝑫𝑬_+,$$
where
(3.23)
$$𝑫\frac{1}{2}(1𝜾^1𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}𝜾^2𝝈_𝗓^{\mathrm{\hspace{0.17em}2}})\frac{1}{2}(1𝜾^1𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}𝜾^3𝝈_𝗓^{\mathrm{\hspace{0.17em}3}})\mathrm{}\frac{1}{2}(1𝜾^1𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}𝜾^N𝝈_𝗓^N).$$
It can be shown that right-multiplication by $`𝑫`$, unlike $`𝑪`$, reduces the dimensionality to $`2^{N+1}`$, thereby permitting the objects in this *reduced* even subalgebra $`(𝒢_3^+)\text{}^N/𝑫`$ to be regarded as spinors, analogous to $`𝒢_3^+`$ for a single qubit. In the corresponding left-ideal, $`𝑲𝜾^1𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}𝑫\mathrm{}𝜾^N𝝈_𝗓^N𝑫`$ serves as the unit imaginary, since $`𝑲^2=𝑫`$, but is required to always operate from the *right*. Henceforth, unless otherwise mentioned, we will regard spinors $`(𝒢_3^+)\text{}^N/𝑫=((𝒢_3^+)\text{}^N/𝑪)/(𝑪𝑫)`$ as a left-ideal in the $`𝑪`$-correlated even subalgebra, drop both $`𝑪`$ and the superscripts on the $`𝜾`$’s as above, and use $`𝑫`$ as a short-hand for $`𝑪𝑫=𝑫𝑪`$.
In the case of two qubits, for example, the identifications are induced by $`𝑫`$ are
$`|\mathrm{\hspace{0.17em}00}`$ $`1`$ $`\stackrel{𝑫}{}`$ $`𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}2}}`$ $`𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}`$ $`\stackrel{𝑫}{}`$ $`𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}2}}`$
(3.24) $`|\mathrm{\hspace{0.17em}01}`$ $`𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}`$ $`\stackrel{𝑫}{}`$ $`𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}`$ $`𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}`$ $`\stackrel{𝑫}{}`$ $`𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}`$
$`|\mathrm{\hspace{0.17em}10}`$ $`𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}`$ $`\stackrel{𝑫}{}`$ $`𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}2}}`$ $`𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}`$ $`\stackrel{𝑫}{}`$ $`𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}2}}`$
$`|\mathrm{\hspace{0.17em}11}`$ $`𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}`$ $`\stackrel{𝑫}{}`$ $`𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}`$ $`𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}`$ $`\stackrel{𝑫}{}`$ $`𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}`$
(where the two columns differ by operation with $`𝑲`$). From this it may be seen that any “spinor” in $`(𝒢_3^+)^2/𝑫`$ can be written as
(3.25)
$$\begin{array}{cc}\hfill 𝚿=& (\text{}(\alpha _0+\beta _0𝑲)𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}(\alpha _1+\beta _1𝑲)\hfill \\ & \text{}𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}(\alpha _2+\beta _2𝑲)+𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}(\alpha _3+\beta _3𝑲))𝑫\hfill \end{array}$$
(cf. \[SLD99\]). Alternatively, again using Eq. (3), a unit norm spinor may be factorized into a product of entities in the *correlated and reduced* even subalgebra, namely $`𝚿=𝑹^1𝑺^2𝑻𝑷𝑫𝑪`$, where
(3.26)
$$\begin{array}{cc}\hfill 𝑹^1& e^{𝜾\varphi 𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}/2}e^{𝜾\theta 𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}/2},\hfill \\ \hfill 𝑺^2& e^{𝜾\phi 𝝈_𝗓^{\mathrm{\hspace{0.17em}2}}/2}e^{𝜾\vartheta 𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}/2},\hfill \end{array}\begin{array}{cc}\hfill 𝑻& \mathrm{cos}(\varsigma /2)\mathrm{sin}(\varsigma /2)𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}𝑲,\hfill \\ \hfill 𝑷& e^{\tau 𝑲/2}.\hfill \end{array}$$
Thus when $`\varsigma =\pi `$, the factor $`𝑻`$ becomes
(3.27)
$$𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}𝑲=(𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}})(𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}})𝑲=e^{(\pi /2)𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}}e^{(\pi /2)𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}}𝑲,$$
so that the arguments of the exponentials involving $`𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}`$ and $`𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}`$ in the first two factors are shifted by $`\pi /2`$ while the total phase is shifted by $`\tau =\pi `$. It follows that $`𝑻`$ rotates the first two factors in the planes defined by their conjugate spinors $`[r_2^{};r_1^{}],[s_2^{};s_1^{}]`$. Thus on right-multiplying by $`𝑬_+`$ and expanding in the usual basis, we obtain (up to an overall phase)
(3.28)
$$\begin{array}{cc}\hfill \underset{¯}{𝚿}=& \mathrm{cos}(\frac{\varsigma }{2})e^{ı\tau /2}\left[\begin{array}{c}\mathrm{cos}(\frac{\theta }{2})e^{ı\varphi /2}\\ \mathrm{sin}(\frac{\theta }{2})e^{ı\varphi /2}\end{array}\right]\left[\begin{array}{c}\mathrm{cos}(\frac{\vartheta }{2})e^{ı\phi /2}\\ \mathrm{sin}(\frac{\vartheta }{2})e^{ı\phi /2}\end{array}\right]+\hfill \\ & \mathrm{sin}(\frac{\varsigma }{2})e^{ı\tau /2}\left[\begin{array}{c}\mathrm{sin}(\frac{\theta }{2})e^{ı\varphi /2}\\ \mathrm{cos}(\frac{\theta }{2})e^{ı\varphi /2}\end{array}\right]\left[\begin{array}{c}\mathrm{sin}(\frac{\vartheta }{2})e^{ı\phi /2}\\ \mathrm{cos}(\frac{\vartheta }{2})e^{ı\phi /2}\end{array}\right]\hfill \end{array}$$
This is known as the *Schmidt decomposition* \[EK95\]. It is useful in studying the *entanglement* of bipartite quantum systems, which (in conventional terms) means that $`|\psi ^2`$ cannot be written as a product $`|\psi ^1|\psi ^2|\psi ^1|\psi ^2|\psi ^1\psi ^2`$ for any one-particle spinors $`|\psi ^1,|\psi ^2`$. In fact it is just the singular value decomposition in disguise, since (for example) on arranging the entries of a two-qubit spinor $`|\psi =[\psi _1;\mathrm{};\psi _4]`$ in a $`2\times 2`$ matrix, we can write
(3.29)
$$\underset{¯}{𝚿}\left[\begin{array}{cc}\psi _1& \psi _3\\ \psi _2& \psi _4\end{array}\right]=\underset{¯}{𝑼}\underset{¯}{𝑽}\underset{¯}{𝑾}^{}=\underset{¯}{𝒖}^1v^{11}(\underset{¯}{𝒘}^1)^{}+\underset{¯}{𝒖}^2v^{22}(\underset{¯}{𝒘}^2)^{},$$
where $`\underset{¯}{𝑽}`$ is a $`2\times 2`$ diagonal matrix containing the singular values $`v^{11}v^{22}0`$ and $`\underset{¯}{𝑼}`$, $`\underset{¯}{𝑾}`$ are unitary matrices with columns $`\underset{¯}{𝒖}^k`$, $`\underset{¯}{𝒘}^k`$, respectively. Since the entries of the dyadic products $`\underset{¯}{𝒖}^1(\underset{¯}{𝒘}^1)^{}`$, $`\underset{¯}{𝒖}^2(\underset{¯}{𝒘}^2)^{}`$ are exactly the same as the Kronecker matrix products $`\underset{¯}{𝒖}^1\underset{¯}{𝒘}^1`$, $`\underset{¯}{𝒖}^2\underset{¯}{𝒘}^2`$, the equivalence with Eq. (3.28) follows with $`v^{11}\mathrm{cos}(\varsigma /2)`$, $`v^{22}\mathrm{sin}(\varsigma /2)`$, and the Kronecker products of the columns of $`\underset{¯}{𝑼}`$ and $`\underset{¯}{𝑾}`$ identified with conjugate pairs of single qubit spinors whose relative phases are given by $`\mathrm{exp}(\pm ı\tau /2)`$.
Clearly a two-qubit spinor is unentangled if and only if $`v^{11}=1`$, which is equivalent to $`\varsigma =0`$ or $`𝑻=1`$. Thus $`𝑻`$ describes the entanglement of the qubits, and is accordingly called the *tangler*. The geometric algebra approach clearly provides deeper insight into the structure of entanglement than does one based on mechanical matrix algebra. In particular, the fact that $`\stackrel{~}{𝚿}𝚿`$ is even and reversion-symmetric in the Dirac as well as the Pauli algebra implies that it is the sum of a scalar and a four-vector in the two-particle Dirac algebra $`𝒢_{2,6}`$. Since Lorentz transformations of the spinors cancel, this entity is in fact a Lorentz invariant, and dividing out the total phase $`𝑷`$ as $`𝑷(\stackrel{~}{𝚿}𝚿)\stackrel{~}{𝑷}`$ yields the *square* of the tangler directly. The availability of such powerful methods of manipulating entities in the multiparticle Dirac algebra promises to be useful in finding analogs of the Schmidt decomposition for three or more qubits.
## 4. Quantum Operations on Density Operators
Quantum computers operate on information stored in the states of quantum systems. The systems are usually assumed to be arrays of distinguishable qubits (two-state subsystems), whose basis states $`|\mathrm{\hspace{0.17em}0}`$ and $`|\mathrm{\hspace{0.17em}1}`$ correspond to the binary digits $`0`$ and $`1`$, respectively, while the operations are usually taken to be unitary. General unitary transformations of the qubits are built up from simpler ones that affect only a few qubits at a time, which are called *quantum logic gates*. The representation of these gates in suitable products of Clifford algebras has been described in Refs. \[SCH98, Vla01\]. The goal here will be to show how gates act upon spinors in the even subalgebra, and how they can be extended to a wider class of nonunitary quantum operations on density operators.
Given the isomorphism between the algebra of $`2^N\times 2^N`$ matrices over $``$ and $`𝒢_3^N/𝑪`$ relative to a choice of basis in each particle space, it is straightforward to interpret matrices in the former as geometric entities in the latter. A matrix $`\underset{¯}{𝑼}𝖴(2^N)`$, however, does not generally correspond to an entity $`𝑼`$ in the *even* subalgebra $`(𝒢_3^+)\text{}^N/𝑪`$, so that $`𝑼𝚿(𝒢_3^+)\text{}^N/𝑫`$ for a general spinor $`𝚿(𝒢_3^+)\text{}^N/𝑫`$. Nevertheless, letting $`𝑬_{}_q𝑬_{}^q`$ be the idempotent “opposite” to $`𝑬_+`$, and noting that this satisfies $`𝑬_+𝑬_{}=0`$, the product of $`𝑼𝚿`$ with $`𝑬_+`$ may be written as
(4.1)
$$𝑼𝚿𝑬_+=(𝑼𝚿𝑬_++\widehat{𝑼}𝚿𝑬_{})𝑬_+=2𝑼𝚿𝑬_+_+𝑬_+,$$
where the “hat” on $`\widehat{𝑼}`$ denotes its image under inversion in the origin (so that $`\widehat{𝑬}_+=𝑬_{}`$), and hence $`\underset{¯}{}_+`$ is a projection onto the even subalgebra. Because $`(𝒢_3^+)\text{}^N/𝑪`$ and $`𝖴(2^N)`$ are both $`(2^{2N})`$-dimensional, nothing is lost in this projection! Thus we can drop the right-factor of $`𝑬_+`$ as usual, and define the action of $`𝑼`$ on $`𝚿(𝒢_3^+)\text{}^N/𝑫`$ as
(4.2)
$$𝑼𝚿2𝑼𝚿𝑬_+_+.$$
More generally, the usual action of the Pauli matrices on spinors corresponds to the following action of the basis vectors on the reduced even subalgebra \[DLG93\]:
(4.3)
$$𝝈_\mu 𝚿𝝈_\mu 𝚿𝝈_𝗓,𝜾𝚿𝜾𝚿𝝈_𝗓$$
The simplest logic gate is the NOT of a single qubit, which operates on the computational basis as follows:
(4.4)
$$\underset{¯}{𝑵}|\mathrm{\hspace{0.17em}0}=|\mathrm{\hspace{0.17em}1}𝜾𝝈_𝗒,\underset{¯}{𝑵}|\mathrm{\hspace{0.17em}1}=|\mathrm{\hspace{0.17em}0}1$$
Thus it might appear reasonable to represent the NOT by $`𝑵𝜾𝝈_𝗒\mathrm{𝖲𝖴}(2)`$, but when $`𝜾𝝈_𝗒`$ is applied a superposition $`(1𝜾𝝈_𝗒)/\sqrt{2}(|\mathrm{\hspace{0.17em}0}+|\mathrm{\hspace{0.17em}1})/\sqrt{2}`$, we get $`(1+𝜾𝝈_𝗒)/\sqrt{2}(|\mathrm{\hspace{0.17em}0}|\mathrm{\hspace{0.17em}1})/\sqrt{2}`$ instead of $`(|\mathrm{\hspace{0.17em}0}+|\mathrm{\hspace{0.17em}1})/\sqrt{2}`$ again. For a single qubit this difference is just an overall rotation by $`\pi `$ about $`𝝈_𝗓`$, but a second qubit can be affected by this phase difference between the first qubit’s states. Therefore the correct representation of the NOT gate in $`\mathrm{𝖲𝖴}(2)`$ is actually $`𝑵\pm 𝜾𝝈_𝗑`$, which preserves this superposition up to an irrelevant overall phase shift: $`ı\underset{¯}{𝝈}_𝗑(|\mathrm{\hspace{0.17em}0}+|\mathrm{\hspace{0.17em}1})/\sqrt{2}`$
(4.5)
$$\begin{array}{cc}& (𝜾𝝈_𝗑)(1𝜾𝝈_𝗒)/\sqrt{2}=𝜾𝝈_𝗑(1𝜾𝝈_𝗒)/\sqrt{2}\hfill \\ \hfill =& 𝜾(𝝈_𝗑(1𝜾𝝈_𝗒)𝝈_𝗓/\sqrt{2})=𝜾(𝜾𝝈_𝗒+1)/\sqrt{2}\hfill \\ \hfill =& (𝜾𝝈_𝗒+1)𝜾𝝈_𝗓/\sqrt{2}ı(|\mathrm{\hspace{0.17em}0}+|\mathrm{\hspace{0.17em}1})/\sqrt{2}\hfill \end{array}$$
More interesting logical operations on the qubits must be able to transform the state of one *conditional* on that of another. The usual way in which this is done is via the c-NOT or *controlled-NOT* gate. As a matrix in $`\mathrm{𝖲𝖴}(4)`$, this is represented in the computational basis by
(4.6)
$$\underset{¯}{𝑵}^{2|1}\sqrt{ı}\left[\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right],$$
which makes it clear that this operation NOT’s the second qubit whenever the first is $`1`$. The corresponding operator in geometric algebra is
(4.7)
$$𝑵^{2|1}(1+𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}1}})/\sqrt{2}\left(𝑬_+^{\mathrm{\hspace{0.17em}1}}+𝑬_{}^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}\right).$$
This may also be written in exponential form as
(4.8)
$$e^{𝜾\pi 𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}/4}e^{𝜾\pi 𝑬_{}^{\mathrm{\hspace{0.17em}1}}𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}/2}=e^{𝜾\pi (𝑬_{}^{\mathrm{\hspace{0.17em}1}}(1𝝈_𝗑^{\mathrm{\hspace{0.17em}2}})/2\mathrm{\hspace{0.17em}1}/4)}.$$
Physical implementations of this operation by e.g. NMR typically expand the exponential into a product of relatively simple commuting factors which can be performed sequentially \[HSTC00\].
Note that since $`(1𝝈_𝗑^{\mathrm{\hspace{0.17em}2}})/2`$ is also an idempotent, $`𝑵^{2|1}`$ differs from $`𝑵^{1|2}`$ by a swap of the $`𝗑`$ and $`𝗓`$ axes for both qubits. This self-inverse operation, called the *Hadamard transform* $`𝑯`$, is simply a rotation by $`\pi `$ about the $`(𝝈_𝗑+𝝈_𝗓)/\sqrt{2}`$ axis. Sandwiching $`𝑵^{2|1}`$ by Hadamards $`𝑯^{\mathrm{\hspace{0.17em}2}}=𝜾(𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}+𝝈_𝗓^{\mathrm{\hspace{0.17em}2}})/\sqrt{2}`$ to just the second qubit gives
(4.9)
$$𝑯^{\mathrm{\hspace{0.17em}2}}𝑵^{2|1}𝑯^{\mathrm{\hspace{0.17em}2}}=e^{𝜾\pi (𝑬_{}^{\mathrm{\hspace{0.17em}1}}𝑯^{\mathrm{\hspace{0.17em}2}}(1𝝈_𝗑^{\mathrm{\hspace{0.17em}2}})𝑯^{\mathrm{\hspace{0.17em}2}}/21/4)}=e^{𝜾\pi (𝑬_{}^{\mathrm{\hspace{0.17em}1}}𝑬_{}^{\mathrm{\hspace{0.17em}2}}\mathrm{\hspace{0.17em}1}/4)},$$
so the c-NOT can also be viewed as a rotated phase shift of the state $`|\mathrm{\hspace{0.17em}11}`$ by $`\pi `$. The Hadamard gate has the important feature of transforming basis states into superpositions thereof; indeed, as an element of the even subalgebra, it actually represents the spinor of a uniform superposition directly:
(4.10)
$$\begin{array}{cc}\hfill \underset{¯}{𝑯}|\mathrm{\hspace{0.17em}0}=& ı(|\mathrm{\hspace{0.17em}0}+|\mathrm{\hspace{0.17em}1})/\sqrt{2}(𝜾𝝈_𝗓+𝜾𝝈_𝗑)/\sqrt{2}𝑫\hfill \\ \hfill \underset{¯}{𝑯}|\mathrm{\hspace{0.17em}1}=& ı(|\mathrm{\hspace{0.17em}0}|\mathrm{\hspace{0.17em}1})/\sqrt{2}(𝜾𝝈_𝗓𝜾𝝈_𝗑)/\sqrt{2}𝑫\hfill \end{array}$$
Thus, by using the relations (4.3), we can show that applying a Hadamard to one of two qubits in the state $`|\mathrm{\hspace{0.17em}11}`$ followed by a c-NOT gate to the other yields the entangled singlet state: $`\underset{¯}{𝑵}^{2|1}\underset{¯}{𝑯}^1|\mathrm{\hspace{0.17em}11}`$
(4.11)
$$\begin{array}{cc}& \frac{1}{2}\left(\text{}(1+𝜾)𝑬_+^1+(1𝜾)𝑬_{}^1𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}\right)\left((𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}1}})(𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}})\right)\hfill \\ \hfill =& \frac{1}{2}\left(\text{}\left((1+𝜾)(𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}1}})\right)(𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}})+\left((1𝜾)(𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}1}})\right)\left((𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}2}})(𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}})\right)\right)\hfill \\ \hfill =& \frac{1}{2}(1𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}1}})𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}\frac{1}{2}(𝜾𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}})𝜾𝝈_𝗓^{\mathrm{\hspace{0.17em}2}}\stackrel{𝑫}{}\frac{𝜾}{2}(𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}+𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝝈_𝗑^{\mathrm{\hspace{0.17em}1}})\hfill \end{array}$$
$`\sqrt{ı}(|\mathrm{\hspace{0.17em}10}|\mathrm{\hspace{0.17em}01})/\sqrt{2}\sqrt{ı}|\psi _{}`$. “Quantum” gates like $`𝑯`$ are not, of course, found in conventional boolean logic, and are an essential component of all quantum algorithms that are more efficient than their classical counterparts \[EJ98, CEMM98\]. Indeed, the c-NOT gate together with general single qubit rotations are known to generate $`\mathrm{𝖲𝖴}(2^N)`$, and hence are *universal* for quantum logic \[BBC<sup>+</sup>95\].
It turns out that unitary transformations are not the most general sort of operation that can be applied to a quantum system. Most such *quantum operations*, however, produce a statistical outcome, and the ensemble of possible outcomes must be described by a density operator. The previous definition (Eq. (2.20)) of the density operator of an ensemble of identical and noninteracting qubits may be extended to an ensemble of multi-qubit systems as follows:
(4.12)
$$\underset{¯}{𝝆}\overline{|\psi \psi |}𝝆\overline{(𝚿𝑫)𝑬_+(𝚿𝑫)^{}}=\overline{𝚿𝑬_+\stackrel{~}{𝚿}}𝑪$$
Suppressing the correlator $`𝑪`$ as usual, $`𝝆`$ may also be expressed in diagonal form as
(4.13)
$$𝝆=𝑹\left(_{k=0}^{2^N}\rho _k𝑬_{(k)}\right)𝑹^{}=_{k=0}^{2^N}\rho _k𝒓_k𝒓_k^{}.$$
where $`𝑹𝒢_3^N/𝑪`$ corresponds to a unitary matrix $`\underset{¯}{𝑹}𝖴(2^N)`$ (in the usual $`𝝈_𝗓`$ coordinate system), and $`0\rho _k1`$ are the eigenvalues of $`𝝆`$. The idempotents $`𝑬_{(k)}`$ are given by $`_q𝑬\text{}_{ϵ_k^q}^q|\chi _k^1\mathrm{}\chi _k^N\chi _k^1\mathrm{}\chi _k^N|`$, where $`ϵ_k^q12\chi _k^q`$ with $`\chi _k^q`$ equal to the $`q`$-th bit in the binary expansion of $`k\{0,\mathrm{},2^N1\}`$. It follows that $`|\rho _k𝒓_k𝑹𝑬_{(k)}`$ for $`\rho _k>0`$ are the spinors of the (unique, if $`\rho _k\rho _{\mathrm{}}kl`$) minimal ensemble that realizes $`𝝆`$, which therefore describes a pure state if and only if it has rank $`1`$ as an operator.
Note that by Eq. (4.2), the density operator transforms under unitary operations as
(4.14)
$$𝝆\overline{(𝑼𝚿)𝑬_+(𝑼𝚿)^{}}=𝑼\overline{𝚿𝑬_+\stackrel{~}{𝚿}}𝑼^{}=𝑼𝝆𝑼^{}.$$
Similarly, the ensemble-average expectation value of any observable $`𝑶=𝑶^{}𝒢_3^N/𝑪`$ is
(4.15)
$$\overline{\psi |\underset{¯}{𝑶}|\psi }2^N\overline{𝑬_+\stackrel{~}{𝚿}𝑶𝚿𝑬_+}_0=2^N𝑶\overline{𝚿𝑬_+\stackrel{~}{𝚿}}_02^N𝑶𝝆_0,$$
just as shown in Eq. (2.17) for single qubit ensembles. In contrast to the case of a single qubit, however, the geometric interpretation of these observables is not straightforward. While one can certainly express $`𝝆`$ as a finite ensemble average $`_kp_k𝚿_k𝑬_+\stackrel{~}{𝚿}_k`$ (where the $`p_k>0`$ are probabilities with $`_kp_k=1`$), this decomposition is highly nonunique. The minimal ensemble obtained by diagonalization, on the other hand, will generally include entangled spinors $`𝒓_k`$, for which the expectation value $`𝑶𝒓_k^{}𝒓_k^{}_0`$ cannot be expressed as a product of inner products of the factors of $`𝑶=𝑶^1\mathrm{}𝑶^N`$ with the polarization vectors of the individual qubits (indeed, $`𝑶`$ itself need not be factorizable!).
The best one can do is to expand $`𝝆`$ in the *product operator* basis consisting of all $`2^{2N}`$ products of the basis vectors $`𝝈_\mu ^q`$, i.e.
(4.16)
$$𝝆=\underset{\mu ^1,\mathrm{},\mu ^N\{0,𝗑,𝗒,𝗓\}}{}\rho _{\mu ^1\mathrm{}\mu ^N}𝝈_{\mu ^1}^{\mathrm{\hspace{0.17em}1}}\mathrm{}𝝈_{\mu ^N}^N,$$
where $`\rho _{\mu ^1\mathrm{}\mu ^N}`$ and $`𝝈_0^q1`$ for notational convenience. The utility of this basis is most simply demonstrated via a concrete example, namely NMR spectroscopy. Here one is given a liquid sample consisting of identical molecules whose nuclear spins are chemically distinguishable, and hence constitutes an ensemble of multi-qubit systems (see \[CPH98, HCST00, HSTC00\] and references therein). The energy of interaction between the spins and an external magnetic field along $`𝗓`$ is given by an observable called the Zeeman Hamiltonian, $`𝒁(\omega ^1𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}+\mathrm{}+\omega ^N𝝈_𝗓^N)/2`$, where $`\omega ^q`$ is the energy difference between the $`|\mathrm{\hspace{0.17em}0}`$ and $`|\mathrm{\hspace{0.17em}1}`$ states of the $`q`$-th spin in the field. In thermal equilibrium at room temperatures, the polarization of the spins relative to the strongest available fields is typically $`\alpha 10^6`$, and the density operator of the ensemble is essentially $`𝝆_{\mathrm{𝖾𝗊}}=2^N(1+\alpha (𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}+\mathrm{}+𝝈_𝗓^N))`$. Via a suitable pulse of radio-frequency radiation, this may be rotated to $`𝝆_{\mathrm{𝖾𝗊}}^{}2^N(1+\alpha (𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}+\mathrm{}+𝝈_𝗑^N))`$, which evolves under the interaction with the field as
(4.17)
$$e^{𝜾𝒁t}𝝆_{\mathrm{𝖾𝗊}}^{}e^{𝜾𝒁t}=\begin{array}{cc}& 2^N(1+\alpha (\mathrm{cos}(\omega ^1t)𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}\mathrm{sin}(\omega ^1t)𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}+\mathrm{}\hfill \\ & \mathrm{}+\mathrm{cos}(\omega ^Nt)𝝈_𝗑^N\mathrm{sin}(\omega ^Nt)𝝈_𝗒^N)).\hfill \end{array}$$
Thus on measuring the total magnetization $`M_𝗑`$ along the $`𝗑`$ axis, $`𝑶\gamma (𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}+\mathrm{}+𝝈_𝗑^N)`$ (where $`\gamma `$ is the nuclear gyromagnetic ratio), one obtains the sum of the projections of the rotating magnetization vectors of the spins along the $`𝗑`$-axis, i.e.
(4.18)
$$M_𝗑(t)=\alpha \gamma (\mathrm{cos}(\omega ^1t)+\mathrm{}+\mathrm{cos}(\omega ^Nt)),$$
whose Fourier transform reveals the contribution from each spin. The way in which the factors of product operators transform like vectors under rotations accounts in large part for the computational utility of the product operator basis. Of course, unless it is a natural part of the problem at hand (as in NMR), one is better off not chosing a basis at all!
A *normal quantum operation* is a linear transformation of the density operator that may be written in operator sum form as \[Kra83\]
(4.19)
$$𝝆\mathrm{\Omega }(𝝆)_k𝑸_k𝝆𝑸_k^{},$$
where the *Kraus operators* $`𝑸_k𝒢_3^N/𝑪`$ satisfy $`_k𝑸_k^{}𝑸_k=1`$. The term “normal” here<sup>3</sup><sup>3</sup>3 We prefer to avoid the more common but clumsy and matrix-bound term “trace-preserving”. refers to the fact that such an operation preserves the scalar part of $`𝝆`$, since
(4.20)
$$\mathrm{\Omega }(𝝆)_0=_k𝑸_k𝝆𝑸_k^{}_0=𝝆_k𝑸_k^{}𝑸_k_0=2^N.$$
It is also easily seen that such quantum operations are *positive*, in that they preserve the positive-definiteness of $`𝝆`$; in fact, these operations have a yet stronger property known as *complete positivity*, meaning that if the qubits to which $`\mathrm{\Omega }`$ applies are embedded in a larger system, then applying $`\mathrm{\Omega }`$ to just those qubits preserves the positive-definiteness of the larger system’s density operator. That this is a nontrivial extension of positivity is shown by the two-qubit “partial transpose” operator $`\mathrm{\Omega }_𝖳^1`$, which carries $`𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}`$ but leaves all the other operator factors unchanged; this is clearly positive on density operators not involving the second qubit, but acts on the density operator of the singlet state (Eq. (LABEL:eq:singlet)) as
(4.21)
$$\begin{array}{cc}& 𝝍_{}\begin{array}{cc}\hfill & \frac{1}{4}(𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}+𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}𝝈_𝗒^{\mathrm{\hspace{0.17em}2}})𝑬_+^1𝑬_+^2(𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}+𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}𝝈_𝗒^{\mathrm{\hspace{0.17em}2}})\hfill \\ \hfill =& \frac{1}{4}(1𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}𝝈_𝗓^{\mathrm{\hspace{0.17em}2}})\hfill \end{array}\hfill \\ \hfill & \mathrm{\Omega }_𝖳^1(𝝍_{})=\frac{1}{4}(1𝝈_𝗑^{\mathrm{\hspace{0.17em}1}}𝝈_𝗑^{\mathrm{\hspace{0.17em}2}}+𝝈_𝗒^{\mathrm{\hspace{0.17em}1}}𝝈_𝗒^{\mathrm{\hspace{0.17em}2}}𝝈_𝗓^{\mathrm{\hspace{0.17em}1}}𝝈_𝗓^{\mathrm{\hspace{0.17em}2}}),\hfill \end{array}$$
which has eigenvalues $`[1/2,\mathrm{\hspace{0.17em}1}/2,\mathrm{\hspace{0.17em}1}/2,1/2]`$.
A quantum operation $`\mathrm{\Omega }`$ is called *unital* if it preserves the identity itself, i.e. $`\mathrm{\Omega }(1)=1`$, or equivalently, $`_k𝑸_k𝑸_k^{}=1`$. Perhaps the most important example of a normal unital operation is found in the *contraction*<sup>4</sup><sup>4</sup>4Otherwise known as the “partial trace”. by a single qubit $`q\{1,\mathrm{},N\}`$, which may be written in operator sum form as \[SCH98\]:
(4.22)
$$2𝝆^q𝑬_+^q𝝆𝑬_+^q+𝑬_{}^q𝝆𝑬_{}^q+𝝈_𝗑^q(𝑬_+^q𝝆𝑬_+^q+𝑬_{}^q𝝆𝑬_{}^q)𝝈_𝗑^q$$
This may also be expressed by dropping all terms in the product operator expansion of $`𝝆`$ depending on $`q`$, and multiplying the remaining terms by a factor of $`2`$. Note that, while $`\underset{¯}{}^q`$ is normal and unital, this factor means that the contraction itself is neither. The factor is nevertheless required if the result is to be interpreted as a density operator for the remaining qubits, since the contraction by the second qubit of the above singlet state is $`𝝍_{}^2=1/4`$ (not $`1/2`$).
This example also illustrates an important way in which general quantum operations are realized in practice, despite the fact that the universe as a whole evolves unitarily. As shown previously, the superposition state with spinor $`𝚿^1=(1𝜾𝝈_𝗒^{\mathrm{\hspace{0.17em}1}})/\sqrt{2}`$ is converted into the singlet state with density operator $`𝝍_{}`$ by letting it interact with a second qubit so as to effect the c-NOT operation $`𝑵^{2|1}`$. The contraction then corresponds to “discarding” the second qubit (i.e. ensuring that it does not further interact with the first and hence can be ignored), which yields the density operator $`1/2`$ of the totally mixed state for the first qubit. Since the basis states $`𝑬_\pm ^{\mathrm{\hspace{0.17em}1}}`$ are unaffected by $`𝑵^{2|1}`$, the net quantum operation on the first qubit corresponds to what is known in quantum communications theory as the *phase damping channel*
(4.23)
$$𝝆(1p)𝝆+p𝑬_+𝝆𝑬_++p𝑬_{}𝝆𝑬_{}$$
with damping parameter $`p=1`$. Phase damping is also known as $`T_2`$ relaxation in NMR, and as *decoherence* in quantum information processing; it is widely believed to be the dominant mechanism by which classical statistical mechanics arises from the underlying unitary dynamics \[GJK<sup>+</sup>96\].
To illustrate the utility of geometric algebra in the study of general quantum operations, an eigenvalue characterization of normal, unital, one-bit quantum operations $`\mathrm{\Omega }`$ will now be derived. This characterization was originally given by Fujiwara & Algoet \[FA99\], although the derivation here parallels that more recently obtained using matrix methods King & Ruskai \[KR00\]. This derivation will regard the Kraus operators $`𝑸_k𝒢_3`$ as “complex quaternions” $`𝑨_k+𝜾𝑩_k`$ with $`𝑨_k,𝑩_k𝒢_3^+`$, and consider the action of an arbitrary operation $`\mathrm{\Omega }`$ on the scalar and vector parts of $`𝝆(1+𝒓)/2`$ separately.
First, the action on $`1`$ is
(4.24)
$$\begin{array}{cc}\hfill \mathrm{\Omega }(1)=& _k(𝑨_k+𝜾𝑩_k)(𝑨_k+𝜾𝑩_k)^{}\hfill \\ \hfill =& _k(𝑨_k\stackrel{~}{𝑨}_k+𝑩_k\stackrel{~}{𝑩}_k)+𝜾_k(𝑩_k\stackrel{~}{𝑨}_k𝑨_k\stackrel{~}{𝑩}_k)).\hfill \end{array}$$
The first summation is symmetric with respect to spatial reversion and inversion, i.e. scalar, while the second (excluding the $`𝜾`$) is reversion antisymmetric but inversion symmetric, i.e. a bivector. Writing $`𝑨\alpha +𝜾𝒂`$ and $`𝑩\beta +𝜾𝒃`$, so that $`𝑨_k\stackrel{~}{𝑨}_k+𝑩_k\stackrel{~}{𝑩}_k=\alpha _k^2+𝒂_k^2+\beta _k^2+𝒃_k^2`$, we may further expand the second summation as follows:
(4.25)
$$\begin{array}{cc}& _k\left[(\beta _k+𝜾𝒃_k)(\alpha _k𝜾𝒂_k)(\alpha _k+𝜾𝒂_k)(\beta _k𝜾𝒃_k)\right]\hfill \\ \hfill =& _k\begin{array}{cc}& [(\alpha _k\beta _k+𝜾(\alpha _k𝒃_k\beta _k𝒂_k)+𝒃_k𝒂_k)\text{}\hfill \\ & \text{}(\alpha _k\beta _k+𝜾(\beta _k𝒂_k\alpha _k𝒃_k)+𝒂_k𝒃_k)]\hfill \end{array}\hfill \\ \hfill =& _k\left[\mathrm{\hspace{0.17em}2}𝜾(\alpha _k𝒃_k\beta _k𝒂_k)2𝒂_k𝒃_k\right]\hfill \end{array}$$
Thus $`\mathrm{\Omega }`$ is unital if and only if
(4.26)
$$_k\left(\alpha _k𝒃_k\beta _k𝒂_k\right)=_k𝒂_k\times 𝒃_k,$$
(where “$`\times `$” denotes the cross product) and
(4.27)
$$_k\left(\alpha _k^2+𝒂_k^2+\beta _k^2+𝒃_k^2\right)=1.$$
Similarly, the action on $`𝒓_3`$ is
(4.28)
$$\begin{array}{cc}\hfill \mathrm{\Omega }(𝒓)=& _k(𝑨_k+𝜾𝑩_k)𝒓(\stackrel{~}{𝑨}_k𝜾\stackrel{~}{𝑩}_k)\hfill \\ \hfill =& _k\left(𝑨_k𝒓\stackrel{~}{𝑨}_k+𝑩_k𝒓\stackrel{~}{𝑩}_k\right)+𝜾_k\left(𝑩_k𝒓\stackrel{~}{𝑨}_k𝑨_k𝒓\stackrel{~}{𝑩}_k\right).\hfill \end{array}$$
The first summation is over different dilation/rotations of $`𝒓`$; the second summation (excluding the $`𝜾`$) is reversion and inversion antisymmetric, i.e. a trivector, and may be further expanded as above:
(4.29)
$$\begin{array}{cc}& _k\left[(\beta _k+𝜾𝒃_k)𝒓(\alpha _k𝜾𝒂_k)(\alpha _k+𝜾𝒂_k)𝒓(\beta _k𝜾𝒃_k)\right]\hfill \\ \hfill =& _k\begin{array}{cc}& [(\alpha _k\beta _k𝒓+𝜾(\alpha _k𝒃_k𝒓\beta _k𝒓𝒂_k)+𝒃_k𝒓𝒂_k)\text{}\hfill \\ & \text{}(\alpha _k\beta _k𝒓+𝜾(\beta _k𝒂_k𝒓\alpha _k𝒓𝒃_k)+𝒂_k𝒓𝒃_k)]\hfill \end{array}\hfill \\ \hfill =& _k\left[\mathrm{\hspace{0.17em}2}𝜾\alpha _k𝒃_k𝒓2𝜾\beta _k𝒂_k𝒓+2𝒃_k𝒓𝒂_k\right]\hfill \end{array}$$
Multiplying through by $`𝜾/2`$ converts this to
(4.30)
$$_k\left[\alpha _k𝒃_k𝒓\beta _k𝒂_k𝒓(𝜾(𝒂_k𝒃_k))𝒓\right],$$
which vanishes if and only if
(4.31)
$$_k\left(\alpha _k𝒃_k\beta _k𝒂_k\right)𝒓=_k(𝒃_k\times 𝒂_k)𝒓.$$
If $`\mathrm{\Omega }`$ is normal, this must be true for all $`𝒓`$, which is equivalent to
(4.32)
$$_k\left(\alpha _k𝒃_k\beta _k𝒂_k\right)=_k𝒃_k\times 𝒂_k.$$
A comparison with Eq. (4.26) shows further that $`\mathrm{\Omega }`$ is both unital and normal if and only if $`_k\alpha _k𝒃_k=_k\beta _k𝒂_k`$ and $`_k𝒂_k\times 𝒃_k=0`$.
If we regard a normal $`\mathrm{\Omega }`$ as an affine transformation, i.e.
(4.33)
$$\mathrm{\Omega }(\frac{1}{2}(1+𝒓))=\frac{1}{2}(1+𝒕+_\mu (𝝈_\mu 𝒓)𝒔_\mu )$$
($`𝒔_\mu ,𝒕_3`$), we see from the derivation leading up to (4.26) that
(4.34)
$$\begin{array}{cc}\hfill 𝒕=& 2_k\left(\alpha _k𝒃_k\beta _k𝒂_k𝒂_k\times 𝒃_k\right)\hfill \\ \hfill =& 4_k(\alpha _k𝒃_k\beta _k𝒂_k)=4_k𝒃_k\times 𝒂_k(\text{by (}\text{4.32}\text{))}.\hfill \end{array}$$
Similarly, the above vectors $`𝒔_\mu =\mathrm{\Omega }(𝝈_\mu )_1`$ are
(4.35)
$$𝒔_\mu =_k\left(𝑨_k𝝈_\mu \stackrel{~}{𝑨}_k+𝑩_k𝝈_\mu \stackrel{~}{𝑩}_k\right)(\mu \{𝗑,𝗒,𝗓\}),$$
i.e. a sum of independent dilation/rotations of each basis vector.
It follows that normal and unital quantum operations $`\mathrm{\Omega }`$ may be characterized by finding conditions for the linear map $`𝒓\mathrm{\Omega }(𝒓)_1=\mathrm{\Omega }(𝒓)`$ to be written as a sum of dilation/rotations of $`𝒓`$. To this end, we expand $`𝑨_k𝒓\stackrel{~}{𝑨}_k`$ as
(4.36)
$$\begin{array}{cc}\hfill (\alpha _k+𝜾𝒂_k)𝒓(\alpha _k𝜾𝒂_k)=& \alpha _k^2𝒓+\alpha _k𝜾(𝒂_k𝒓𝒓𝒂_k)+𝒂_k𝒓𝒂_k\hfill \\ \hfill =& \alpha _k^2𝒓+2\alpha _k𝒓\times 𝒂_k+2(𝒓𝒂_k)𝒂_k𝒂_k^2𝒓,\hfill \end{array}$$
with a similar expansion for $`𝑩_k𝒓\stackrel{~}{𝑩}_k`$. Thus on assuming that $`\mathrm{\Omega }`$ is diagonal, i.e. $`𝒔_\mu \lambda _\mu 𝝈_\mu =\mathrm{\Omega }(𝝈_\mu )`$, we get
(4.37)
$$\lambda _\mu 𝝈_\mu =\begin{array}{cc}& 𝝈_\mu _k(\alpha _k^2+\beta _k^2𝒂_k^2𝒃_k^2)+2𝝈_\mu \times _k(\alpha _k𝒂_k+\beta _k𝒃_k)\hfill \\ & +\mathrm{\hspace{0.17em}2}_k\left((𝝈_\mu 𝒂_k)𝒂_k+(𝝈_\mu 𝒃_k)𝒃_k\right).\hfill \end{array}$$
Dotting both sides by $`𝝈_\mu `$ now yields
(4.38)
$$\begin{array}{cc}\hfill \lambda _\mu =& _k\left(\alpha _k^2+\beta _k^2𝒂_k^2𝒃_k^2+2(𝝈_\mu 𝒂_k)^2+2(𝝈_\mu 𝒃_k)^2\right)\hfill \\ \hfill =& 12_k\left(𝒂_k^2+𝒃_k^2(𝝈_\mu 𝒂_k)^2(𝝈_\mu 𝒃_k)^2\right)\hfill \\ \hfill =& 1+2_k\left((𝝈_\mu 𝒂_k)^2+(𝝈_\mu 𝒃_k)^2\right),\hfill \end{array}$$
so we have simple expressions for the eigenvalues. Now consider the vector obtained from the first line of this equation, i.e. $`_{\mu \{𝗑,𝗒,𝗓\}}\lambda _\mu 𝝈_\mu `$
(4.39)
$$\begin{array}{cc}\hfill =& _\mu 𝝈_\mu _k\left(\alpha _k^2+\beta _k^2𝒂_k^2𝒃_k^2+2(𝝈_\mu 𝒂_k)^2+2(𝝈_\mu 𝒃_k)^2\right)\hfill \\ \hfill =& 𝒑_0_k\left(\alpha _k^2+\beta _k^2\right)+_\mu 𝒑_\mu _k\left((𝝈_\mu 𝒂_k)^2+(𝝈_\mu 𝒃_k)^2\right),\hfill \end{array}$$
where
(4.40) $`𝒑_0`$ $`𝝈_𝗑+𝝈_𝗒+𝝈_𝗓,`$ $`𝒑_𝗑`$ $`𝝈_𝗑𝝈_𝗒𝝈_𝗓,`$
$`𝒑_𝗒`$ $`𝝈_𝗑+𝝈_𝗒𝝈_𝗓,`$ $`𝒑_𝗓`$ $`𝝈_𝗑𝝈_𝗒+𝝈_𝗓.`$
Since the coefficients of the $`𝒑`$’s are nonnegative and sum to $`1`$ by (4.27), this shows that the vector $`_\mu \lambda _\mu 𝝈_\mu `$ lies within the tetrahedron $`𝒑_0,𝒑_𝗑,𝒑_𝗒,𝒑_𝗓`$, which is the condition on the eigenvalues found by Fujiwara and Algoet \[FA99\] as well as by King and Ruskai \[KR00\].
It is also known that an arbitrary linear map $`\mathrm{\Omega }`$ has an operator sum representation if and only if it is completely positive \[Sch96\], so the above can also be viewed as a characterization of complete positivity for normal and unital maps of a single qubit’s density operator. Finally, it is worth stressing once again that, because of the isomorphisms which exist between the Pauli algebra and the even subalgebra of the Dirac algebra, every step of the above derivation carries with it a natural interpretation in space-time, and is in fact even easier to carry out when the full power of the Dirac algebra is used.
In conclusion, it is hoped that the forgoing has given the reader a taste of the new insights which geometric algebra can provide into quantum information processing — and an appetite for more!
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# Fluctuation-dissipation relationship in chaotic dynamics
## I Introduction
Although deterministic in principle classically chaotic motion is stochastic in nature. Ever since the early numerical study of Chirikov mapping revealed that the motion of a phase space variable can be characterized by a simple random walk diffusion equation, attempts have been made to describe the chaotic motion in terms of Langevin or Fokker-Planck equations . It is therefore easy to comprehend a close connection between classical chaos and statistical mechanics. Two distinct situations arise in this context. The first one concerns whether classical chaos may serve as a basis for classical statistical mechanics since the ultimate justification of the postulates of statistical mechanics like Boltzmann hypothesis of molecular chaos, ergodicity or the postulate of equal a priori probability rests on the dynamics of each particle . The second one concerns the following : Given that the classical chaotic motion is stochastic how and to what extent one can realize the formulation of statistical mechanics for useful description of classical chaos keeping in mind that one essentially deals here with a few-degree-of-freedom system. The present paper addresses the second issue.
The emergence of stochastic behaviour of the classically chaotic system is due to the loss of correlation of initially nearby trajectories. This is reflected in the nature of the largest Lyapunov exponent whose calculation rests on the linear equation of motion for the separation of these trajectories. When chaos has fully set in, the time dependence of the linear stability matrix or Jacobian of the system in the equation of motion in the tangent space can be described as a stochastic process since the phase space variables behave as stochastic variables. In a number of recent studies we have shown that this fluctuation of the Jacobian is amenable to a theoretical description in terms of the theory of multiplicative noise. This allows us to realize a number of important results of nonequilibrium statistical mechanics, like Kubo relation , fluctuation-decoherence relation , exponential divergence of quantum fluctuations , thermodynamically inspired quantities, e. g. , entropy production in chaotic dynamics. Based on a Fokker-Planck description in the tangent space where the drift and the diffusion coefficients explicitly depend on the phase space variables or dynamical properties of the system, we show that a connection between the two moments in terms of the stochastic parameters which characterize the long time limit of the dynamical system can be established in the spirit of fluctuation-dissipation relation. We verify the theoretical proposition by numerical experiments on a simple dissipative system.
The rest of the paper is organized as follows : In Sec. II we introduce a Fokker-Planck description of the dynamical system in the tangent space and identify the drift and diffusion coefficients which are the functions of fluctuations of the phase space variables. This is followed by solving the Fokker-Planck equation for the steady state distribution required for the calculation of long time averages in Sec III. In Sec. IV the dynamical stochastic parameters which characterize the long time behaviour of the system are introduced. The first one of them is a well-known stochastic parameter closely related to Kolmogorov entropy. With the help of these stochastic parameters we establish a connection between the drift and diffusion coefficients of the Fokker-Planck equation in the spirit of fluctuation-dissipation relation in nonequilibrium statistical mechanics. In Sec. V we illustrate the general method by an explicit numerical example to verify the theoretical proposition. The paper is concluded in Sec. VI.
## II A Fokker-Planck equation for dissipative chaotic dynamics
We are concerned here with a general N-degree-of-freedom system whose Hamiltonian is given by
$$H=\underset{i=1}{\overset{N}{}}\frac{p_i^2}{2m_i}+V(\{q_i\},t),i=1\mathrm{}N$$
(1)
where $`\{q_i,p_i\}`$ are the co-ordinate and momentum of the i-th degree-of-freedom, respectively, which satisfy the generic form of equations
$$\dot{q}_i=\frac{H}{p_i}\mathrm{and}\dot{p}_i=\frac{H}{q_i}.$$
(2)
We now make the Hamiltonian system dissipative by introducing $`\gamma p_i`$ on the right hand side of the second of Eqs.(2). For simplicity we assume $`\gamma `$ to be the same for all the N degrees of freedom. By invoking the symplectic structure of the Hamiltonian dynamics as
$`z_i=\{\begin{array}{cc}q_i\hfill & \mathrm{for}i=1\mathrm{}N,\hfill \\ p_{iN}\hfill & \mathrm{for}i=N+1,\mathrm{}2N.\hfill \end{array}`$
and defining I as
$`I=\left[\begin{array}{cc}0& E\\ E& \gamma E\end{array}\right]`$
where E is an $`N\times N`$ unit matrix, and $`0`$ is an $`N\times N`$ null matrix, the equation of motion for the dissipative system can be written as
$$\dot{z}_i=\underset{j=1}{\overset{2N}{}}I_{ij}\frac{H}{z_j}.$$
(5)
We now consider two nearby trajectories, $`z_i,\dot{z}_i`$ and $`z_i+X_i`$, $`\dot{z}_i+\dot{X}_i`$ at the same time $`t`$ in $`2N`$ dimensional phase space. The time evolution of separation of these trajectories is then determined by
$$\dot{X}_i=\underset{j=1}{\overset{2N}{}}J_{ij}(t)X_j$$
(6)
in the tangent space $`\{X_i\}`$, where
$`J_{ij}={\displaystyle \underset{k}{}}I_{ik}{\displaystyle \frac{^2H}{z_kz_j}}.`$ (7)
Therefore the $`2N\times 2N`$ linear stability matrix $`\underset{¯}{J}`$ assumes the following form
$$\underset{¯}{J}=\left[\begin{array}{cc}0& E\\ \underset{¯}{M(t)}& \gamma E\end{array}\right]$$
(8)
where M is an $`N\times N`$ matrix. Note that the time dependence of stability matrix $`\underset{¯}{J}(t)`$ is due to the second derivative $`\frac{^2H}{z_kz_j}`$ which is determined by the equation of motion (3). The procedure for calculation of $`X_i`$ and the related quantities is to solve the trajectory equation (3) simultaneously with Eq.(4). Thus when the dissipative system described by (3) is chaotic, J(t) becomes (deterministically) stochastic due to the fact that $`z_i`$-s behave as stochastic variables and the equation of motion (6) in the tangent space can be interpreted as a stochastic equation .
In the next step we shall be concerned with a stochastic description of $`\underset{¯}{J}(t)`$ or $`\underset{¯}{M(t)}`$. For convenience we split up M into two parts as
$$\underset{¯}{M}=\underset{¯}{M_0}+\underset{¯}{M_1(t)}$$
(9)
where $`\underset{¯}{M_0}`$ is independent of variables $`\{z_i\}`$ and therefore behaves a sure or constant part and $`\underset{¯}{M_1}`$ is determined by the variables $`\{z_i\}`$ for $`i=1\mathrm{}2N`$. $`\underset{¯}{M_1}`$ refers to the fluctuating part. We now rewrite the equation of motion (6) in tangent space as
$`\dot{X}`$ $`=`$ $`\underset{¯}{J}X`$ (10)
$`=`$ $`L(\left\{X_i\right\},\left\{z_i\right\})`$ (11)
where $`X`$ and $`L`$ are the vectors with $`2N`$ components. Corresponding to (9) $`L`$ in (11) can be split up again to yield
$`\dot{X}=L^0(X)+L^1(X,\{z_i(t)\}),i=1\mathrm{}2N.`$ (12)
Eq.(6) indicates that Eq.(11) is linear in $`\{X_i\}`$. Eqs.(6), (7) and (8) express the fact the first $`N`$ components of $`L^1`$ are zero and the last $`N`$ components of $`L^1`$ are the functions of $`\{X_i\}`$ for $`i=1\mathrm{}N`$. The fluctuation in $`L_i^1`$ is caused by the chaotic variables $`\left\{z_i\right\}`$-s. This allows us to write the following relation (which will be used later on),
$$_XL^1\varphi (\{X_i\})=L^1_X\varphi (\{X_i\})$$
(13)
where $`\varphi (\{X_i\}`$) is any function of $`\{X_i\}`$. $`_X`$ refers to differentiation with respect to components $`\{X_i\}`$ (explicitly $`X_i=\mathrm{\Delta }q_i`$ for $`i=1\mathrm{}N`$ and $`X_i=\mathrm{\Delta }p_i`$ for $`i=N+1\mathrm{}2N`$).
Note that Eq.(12) by virtue of (11) is a linear stochastic differential equation with multiplicative noise where the noise is due to $`\{z_i\}`$ determined by equation of motion (5). This is the starting point of our further analysis.
Eq.(12) determines a stochastic process with some given initial conditions $`\{X_i(0)\}`$. We now consider the motion of a representative point $`X`$ in $`2N`$ dimensional tangent space ($`X_1\mathrm{}X_{2N}`$) as governed by Eq.(12). The equation of continuity which expresses the conservation of points determines the variation of density function $`\varphi (X,t)`$ in time as given by
$$\frac{\varphi (X,t)}{t}=_𝐗L(t)\varphi (X,t).$$
(14)
Expressing $`A_0`$ and $`A_1`$ as
$$A_0=_XL^0\mathrm{and}A_1=_XL^1$$
(15)
we may rewrite the equation of continuity as
$$\frac{\varphi (X,t)}{t}=[A_0+\alpha A_1(t)]\varphi (X,t).$$
(16)
It is easy to recognize that while $`A_0`$ denotes the sure part $`A_1`$ contains the multiplicative fluctuations through $`\{z_i(t)\}`$. $`\alpha `$ is a parameter introduced from outside to keep track of the order of fluctuations in the calculations. At the end we put $`\alpha =1`$.
One of the main results for the linear equations of the form with multiplicative noise may now be in order . The average equation of $`\varphi `$ obeys \[ $`P(x,t)\varphi `$\],
$$\dot{P}=\left\{A_0+\alpha A_1+\alpha ^2_0^{\mathrm{}}𝑑\tau A_1(t)\mathrm{exp}(\tau A_0)A_1(t\tau )\mathrm{exp}(\tau A_0)\right\}P(x,t).$$
(17)
The above result is based on second order cumulant expansion and is valid when fluctuations are small but rapid and the correlation time $`\tau _c`$ is short but finite or more precisely
$$A_1(t)A_1(t^{})=0\mathrm{for}\left|tt^{}\right|>\tau _c$$
(18)
We have, in general, $`A_10`$. Here $`\mathrm{}`$ implies $`\zeta _i\zeta _j=\zeta _i\zeta _j\zeta _i\zeta _j`$.
The Eq.(17) is exact in the limit $`\tau _c0`$. Making use of relation (15) in (17) we obtain
$`{\displaystyle \frac{P}{t}}`$ $`=`$ $`\{L^0\alpha L^1+\alpha ^2{\displaystyle _0^{\mathrm{}}}d\tau L^1(t)\mathrm{exp}(\tau L^0)`$ (20)
$`L^1(t\tau )\mathrm{exp}(\tau L^0)\}P.`$
The above equation can be transformed into the following Fokker-Planck equation ($`\alpha =1`$) for probability density function $`P(X,t)`$, (the details are given in the Appendix A);
$$\frac{P(X,t)}{t}=.FP+\underset{i,j}{}𝒟_{ij}\frac{^2P}{X_iX_j}$$
(21)
where,
$$F=L^0+L^1+Q$$
(22)
and $`Q`$ is a $`2N`$-dimensional vector whose components are defined by
$`Q_j={\displaystyle _0^{\mathrm{}}}R_j^{}𝑑\tau d_1(\tau )d_2(\tau )`$ (23)
Here the determinants $`Det_1(\tau )`$, $`Det_2(\tau )`$ and $`R_j^{}`$ are given by
$`Det_1(\tau )`$ $`=`$ $`\left|{\displaystyle \frac{dX^\tau }{dX}}\right|\mathrm{and}Det_2(\tau )=\left|{\displaystyle \frac{dX}{dX^\tau }}\right|`$ (24)
$`\mathrm{and}R_j^{}`$ $`=`$ $`{\displaystyle \underset{i}{}}L_i^1(X,t){\displaystyle \frac{}{X_i}}{\displaystyle \underset{k}{}}L_k^1(X^\tau ,t\tau ){\displaystyle \frac{X_j}{X_k^\tau }}.`$ (25)
It is easy to recognize $`F`$ as an evolution operator. Because of the dissipative perturbation we note that div $`F<0`$.
The diffusion coefficient $`𝒟_{ij}`$ in Eq.(21) is defined as
$$𝒟_{ij}=_0^{\mathrm{}}\underset{k}{}L_i^1(X,t)L_k^1(X^\tau ,t\tau )\frac{dX_j}{dX_k^\tau }Det_1(\tau )Det_2(\tau )d\tau $$
(26)
We have followed closely van Kampen’s approach to generalized Fokker-Planck equation (21). Before concluding this section several critical remarks regarding this derivation need attention:
First, the stochastic process $`\underset{¯}{M_1(t)}`$ determined by $`\{z_i\}`$ is obtained exactly by solving equations of motion (5) for the chaotic motion of the system. It is therefore necessary to emphasize that we have not assumed any special property of noise, such as, $`\underset{¯}{M_1(t)}`$ is Gaussian or $`\delta `$-correlated. We reiterate Van Kampen’s emphasis in this approach.
Second, the only assumption made about the noise is that its correlation time $`\tau _c`$ is short but finite compared to the coarse-grained timescale over which the average quantities evolve.
Third, we take care of fluctuations upto second order which implies that the deterministic noise is not too strong.
Eq.(21) is the required Fokker-Planck equation in the tangent space $`\{X_i\}`$. Note that the drift and diffusion coefficients are determined by the phase space $`\{z_i\}`$ properties of the chaotic system and directly depend on the correlation functions of the fluctuations of the second derivatives of the Hamiltonian (7).
## III The Steady state distribution and the calculation of averages
In what follows we shall be concerned with the long time limit of the dynamical system. Thus the steady state distribution of the tangent space co-ordinates $`X_i(i=1\mathrm{}2N)`$ are specially relevant for the present purpose. To make all these co-ordinates dimensionless we use the following transformations in Eq.(17)
$`\tau ^{}`$ $`=`$ $`\omega ^{}t,`$
$`y_i`$ $`=`$ $`{\displaystyle \frac{X_i}{d_0}}\mathrm{for}i=1\mathrm{}N,`$
$$y_i=\frac{X_i}{\omega ^{}d_0}\mathrm{for}i=N+1\mathrm{}2N,$$
(27)
where $`\omega ^{}`$ is a scaling constant having dimension of reciprocal of time (a possible choice is the linearized frequency of the dynamical system) and $`\tau ^{}`$ becomes a dimensionless variable. $`d_0`$ is a constant (to be specified later) having the dimension of length. The resulting Fokker-Planck Eq.(17) reduces to
$$\frac{P(y,\tau ^{})}{\tau ^{}}=.F^{}(y)P+\underset{i,j}{}𝒟_{}^{}{}_{ij}{}^{}(y)\frac{^2P}{y_iy_j}.$$
(28)
Note that Eq.(23) is independent of $`d_0`$ since $`F(X)`$ is linear in $`\{X_i\}`$ and $`𝒟`$($`X`$) is quadratic in $`\{X_i\}`$. Next we consider the stationary state of the system ($`\frac{P}{\tau ^{}}=0`$) and make use of the following linear transformation ( with $`\alpha _{2N}=1`$ )
$$U=\underset{i=1}{\overset{2N}{}}\alpha _iy_i$$
(29)
in Eq.(23) to obtain the equation for steady state probability distribution $`P_s(U)`$ :
$$\frac{}{U}\lambda UP_s(U)+𝒟_s\frac{^2P_s}{U^2}=0.$$
(30)
$`\alpha _i`$-s ($`i=1\mathrm{}2N1`$) are the constants to be determined. Here
$$\lambda U=\underset{i}{}\alpha _iF_i^{}(y)$$
(31)
and
$$𝒟_s=\underset{i,j}{}𝒟_{}^{}{}_{ij}{}^{}\alpha _i\alpha _j,$$
(32)
and we disregard the time dependence of $`𝒟^{}`$ under weak noise approximation, to treat $`𝒟^{}`$ as a constant in the usual way.
Putting (24) in (26) and comparing the coefficients of $`y_i`$ on both sides we obtain $`2N`$ algebraic equation (for $`\alpha _i\mathrm{}\mathrm{}\alpha _{2N1}`$ and $`\lambda `$). The set $`\{\alpha _i\}`$ and $`\lambda `$ are therefore known.
The exact steady state solution, $`P_s`$ has the well known Gaussian form which is given by
$$P_s(\{y_i\})=N\mathrm{exp}\left(\frac{\lambda }{2𝒟_s}\underset{i,j}{}\alpha _i\alpha _jy_iy_j\right),$$
(33)
where $`N`$ is the normalization constant. Eq.(28) expresses the probability distribution of tangent space co-ordinates of the dynamical system in the long time limit. The important relevant quantity which measures the separation of initially nearby trajectories when the system has attained the stationary state can be computed by calculating the average of $`_{i=1}^{2N}y_i^2`$. Making use of the distribution (28) we obtain
$$\underset{i=1}{\overset{2N}{}}y_i^2=\frac{𝒟_s}{\lambda }\underset{i=1}{\overset{2N}{}}\frac{1}{\alpha _i^2}$$
(34)
Note that the average as calculated above is a function of $`D_s`$, $`\lambda `$ and $`\alpha _i`$-s which are dependent on the phase space properties of the dynamical system.
## IV Stochastic parameters, connection between $`𝒟_s`$ and $`\lambda `$ ; Fluctuation-dissipation relation
Eq.(25) is a steady state Fokker-Planck equation in tangent space with linear drift and constant diffusion coefficients where the co-ordinates have been expressed as dimensionless variables $`\{y_i\}`$. $`\lambda `$ and $`𝒟_s`$ are the first and second moments, respectively, of the underlying stochastic process. Our objective here is seek for a connection between the two moments. In standard nonequilibrium statistical mechanics this connection is expressed by the fluctuation-dissipation relation through temperature, an equilibrium parameter characterizing the equilibrium state. Our approach here is to follow a somewhat similar procedure. This implies that we search for the stochastic parameters which characterize the long time limit of the nonlinear dynamical system. We show that an appropriate relation between $`𝒟_s`$ and $`\lambda `$ can be established through these parameters.
An important parameter proposed many years ago by Casartelli et. al. (a precursor for the largest Lyapunov exponent used as a measure of regularity or chaoticity of a nonlinear dynamical system) is the long time average of $`\mathrm{ln}\frac{d(t)}{d_0}`$ where $`d_0`$ is the separation of the two initially nearby trajectories and $`d(t)`$ is the corresponding separation at some time $`t`$. To express $`d(t)`$ (having dimension of length) we write $`d(t)=[_{i=1}^N(X_i)^2+_{i=N+1}^{2N}(\frac{X_i}{\omega ^{}})^2]^{\frac{1}{2}}`$. $`d(t)`$ is determined by solving numerically Eqs. (3) and (4) simultaneously or their appropriately transformed version for the initial condition $`z_0`$ corresponding to Eq.(3). In going from j-th to j+1-th step of iteration in course of time evolution any of the components of $`X`$ say $`X_i`$ has to be initialized as $`X_i^{j0}=\frac{X_i^j}{d_j}d_0`$. This initialization implies that at each step iteration starts with same magnitude of $`d_0`$ but the direction of $`d_0`$ for step j+1 is that of $`d(t)`$ for j-th step (considered in terms of the ratio $`\frac{X_i^j}{d_j}`$). For a pictorial illustration we refer to Fig.1 of Ref. . j-th time of iteration implies $`t=jT`$ ($`j=1,2\mathrm{}\mathrm{})`$ and $`T`$ is the characteristic time which corresponds to the shortest ensemble averaged period of nonlinear dynamical system. Thus following Casartelli et. al. a stochastic parameter can be defined by the following time average of $`\mathrm{ln}\frac{d_j}{d_0}`$ as
$$\sigma _n(t,z_0,d_0)=\frac{1}{n}\underset{j}{\overset{n}{}}\mathrm{ln}\frac{d_j}{d_0}$$
(35)
It has been shown that as $`n\mathrm{}`$ $`\sigma _n`$ has a definite value. For the disordered system it is positive and for the regular system it is zero. The difference of $`\sigma _n`$ from the largest Lyapunov exponent is also noteworthy. Our object here is to generalize (30) by defining the other higher order moments (higher than the first $`\sigma _n\mathrm{}`$). To express these quantities we define first
$$\sigma ^{}=\mathrm{ln}\frac{d(t)}{d_0}$$
(36)
We now make use of the transformation (22) to express $`d(t)`$ as a dimensionless quantity in terms of $`\sigma ^{}`$ as follows:
$$\mathrm{ln}\underset{i=1}{\overset{2N}{}}y^2=2\sigma ^{}.$$
(37)
The method of cumulant expansion on the other hand tells us that the average of the sum of $`y_i^2`$ can be written as
$$\underset{i=1}{\overset{2N}{}}y_i^2=\mathrm{exp}\left(\underset{m}{}A_m\right)m=1,2,3,\mathrm{}$$
(38)
where $`A_m`$-s result from cumulants of the stochastic quantity $`2\sigma ^{}`$. $`A_m`$-s are calculated dynamically from the following relations
$`A_1`$ $`=`$ $`m_1,A_2={\displaystyle \frac{1}{2!}}[m_2m_1^2],A_3={\displaystyle \frac{1}{3!}}[m_33m_1m_2+2m_1^3]`$ (39)
$`A_4`$ $`=`$ $`{\displaystyle \frac{1}{4!}}[m_43m_2^24m_1m_3+12m_1^2m_26m_1^4]\mathrm{etc}.`$ (40)
where $`m_k=\frac{2^k}{n}_{j=1}^n\left(\mathrm{ln}\frac{d_j}{d_0}\right)^k`$ \[$`k=1,2,3,4,\mathrm{}`$\]. In the spirit of Ref. we enquire, whether these moments/cumulants reach their steady state values in the long time limit. We have numerically examined the dependence of $`m_k`$-s on various parameters. The parameters are $`n`$, the time, $`d_0`$, the measure of initial separation, the characteristic time $`T`$ (j-th time of iteration implies $`t=jT,j=1,2,\mathrm{}\mathrm{}`$). Our observation is that the limit $`m_k`$ or limit $`A_m`$ as $`n\mathrm{}`$ seems to exist in all cases. We have examined these limits for a number of test cases, e. g. , for Lorentz system, Henon-Heiles system and others. In Fig.(1) we exhibit a typical representative long time behaviour of the cumulants $`A_m(m=1\mathrm{to}\mathrm{\hspace{0.33em}4})`$ for a driven double well potential system discussed in the next section. It is apparent that they attain their long time limits as $`n\mathrm{}`$. Secondly, the first two cumulants are much higher compared to others The first moment is the stochastic parameter defined by Casartelli et. al. as a quantity closely related to Kolmogorov entropy. We are therefore led to believe that the quantities $`A_m`$-s characterize the long time limit or the steady state of a dynamical system.
The relations (33) and (29) can now be combined to give
$$𝒟_s=\frac{\lambda }{_{i=1}^{2n}}\frac{1}{\alpha _i^2}\mathrm{exp}\left(\underset{m}{}A_m\right).$$
(41)
The above relation is the central result of this paper. This establishes a connection between the drift and the diffusion coefficients of the Fokker-Planck equation (25) through the stochastic parameters characterizing long time behaviour of the nonlinear dynamical system. It must be emphasized that both the drift $`\lambda `$ and the diffusion $`𝒟_s`$ coefficients arise from the deterministic stochasticity implied in the dynamical equation of motion (3). The relation (35) is therefore reminiscent of the familiar fluctuation-dissipation relation.
A few points regarding the relation (35) are in order. It is important to note that the fluctuation-dissipation relation in conventional nonequilibrium statistical mechanics is valid for a stochastic system for which the noise is internal. The spiritual root of this relation lies at the dynamic balance between the input of energy into the system from the fluctuations of the surrounding and the output of energy from the system due to its dissipation into the surrounding. The system-reservoir model developed over the last few decades suggests that the coupling between the system and the reservoir is responsible for a common origin of drift and diffusion. In the present theory this common mechanism is the fluctuations of the phase space variables (or second derivative of the Hamiltonian) inherent in both the drift $`\lambda `$ and the diffusion $`D_s`$ coefficients of the Fokker-Planck equation. We point out that the relation is still valid for the pure Hamiltonian system ($`\gamma =0`$). For this reason the relation (35) is somewhat formal in contrast to the standard fluctuation-dissipation relation.
## V An example and numerical verification
To illustrate the theory developed above, we now choose a driven double-well oscillator system with Hamiltonian
$$H=\frac{p_1^2}{2}+aq_1^4bq_1^2+ϵq_1\mathrm{cos}\mathrm{\Omega }t$$
(42)
where $`p_1`$ and $`q_1`$ are the momentum and position variables of the system. $`a`$ and $`b`$ are the constants characterizing the potential. $`ϵ`$ includes the effect of coupling constant and the driving strength of the external field with frequency $`\mathrm{\Omega }`$. This model has been extensively used in recent years for the study of chaotic dynamics .
The dissipative equations of motion for the tangent space variables $`X_1`$ and $`X_2`$ corresponding to $`q_1`$ and $`p_1`$ (Eq.8) read as follows:
$$\frac{d}{dt}\left[\begin{array}{c}X_1\\ X_2\end{array}\right]=\underset{¯}{J}\left[\begin{array}{c}X_1\\ X_2\end{array}\right],\left\{\begin{array}{c}\mathrm{\Delta }q_1=X_1\\ \mathrm{\Delta }p_1=X_2\end{array}\right\}.$$
(43)
where $`\underset{¯}{J}`$ as expressed in our earlier notation $`z_1=q_1`$ and $`z_2=p_1`$ is given by
$`\left(\begin{array}{cc}0& 1\\ \zeta (t)+2b& \gamma \end{array}\right),`$
where $`\zeta (t)=12az_1^2`$. Eq.(43) is thus rewritten as
$$\frac{d}{dt}\left(\begin{array}{c}X_1\\ X_2\end{array}\right)=L^0+L^1$$
(45)
with
$`L^0=\left(\begin{array}{c}X_2\\ 2bX_1\gamma X_2\end{array}\right)\mathrm{and}L^1=\left(\begin{array}{c}0\\ \zeta (t)X_1\end{array}\right),`$
where $`L^0`$ and $`L^1`$ are the constant and the fluctuating parts(vectors), respectively. The fluctuation in $`L^1`$, i.e., in $`\zeta (t)`$ is due to stochasticity of the following chaotic dissipative dynamical equations of motion;
$$\dot{z}_1=z_2\mathrm{and}\dot{z}_2=az_1^3+2bz_1ϵ\mathrm{cos}\mathrm{\Omega }t\gamma z_2.$$
(47)
The result of Eq.(A6) can then be applied and after some algebra the Fokker-Planck equation (17) for the dissipative driven double-well oscillator assumes the following form:
$`{\displaystyle \frac{P}{t}}`$ $`=`$ $`X_2{\displaystyle \frac{P}{X_1}}\omega ^2X_1{\displaystyle \frac{P}{X_2}}+\gamma {\displaystyle \frac{}{X_2}}(X_2P)+𝒟_{21}{\displaystyle \frac{^2P}{X_2X_1}}+𝒟_{22}{\displaystyle \frac{^2P}{X_2^2}}`$ (48)
where
$`𝒟_{21}=X_1^2{\displaystyle _0^{\mathrm{}}}\zeta (t)\zeta (t\tau )\tau e^{\gamma \tau }𝑑\tau `$
and
$$𝒟_{22}=X_1^2_0^{\mathrm{}}\zeta (t)\zeta (t\tau )e^{\gamma \tau }𝑑\tau X_1X_2_0^{\mathrm{}}\zeta (t)\zeta (t\tau )\tau e^{\gamma \tau }𝑑\tau $$
(49)
with
$$\omega ^2=2b+c+c_2,c_2=_0^{\mathrm{}}\zeta (t)\zeta (t\tau )\tau e^{\gamma \tau }𝑑\tau \mathrm{and}c=\zeta .$$
(50)
The similarity of the equation (40) to generalized Kramers’ equation can not be overlooked. This suggests a clear interplay of chaotic diffusive motion and dissipation in the dynamics.
Using the transformation (22) Eq.(40) can be written as
$`{\displaystyle \frac{P}{\tau ^{}}}`$ $`=`$ $`y_2{\displaystyle \frac{P}{y_1}}\overline{\omega }^2y_1{\displaystyle \frac{P}{y_2}}+\overline{\gamma }{\displaystyle \frac{}{y_2}}(y_2P)+𝒟_{}^{}{}_{21}{}^{}{\displaystyle \frac{^2P}{y_2y_1}}+𝒟_{}^{}{}_{22}{}^{}{\displaystyle \frac{^2P}{y_2^2}}`$ (51)
where
$`\overline{\omega }^2={\displaystyle \frac{\omega ^2}{\omega ^2}},\overline{\gamma }={\displaystyle \frac{\gamma }{\omega ^{}}},𝒟_{}^{}{}_{21}{}^{}={\displaystyle \frac{y_1^2(0)}{\omega _{}^{}{}_{}{}^{2}}}{\displaystyle _0^{\mathrm{}}}\zeta (\tau ^{})\zeta (\tau ^{}\tau )\tau e^{\gamma \tau }𝑑\tau \mathrm{and}`$
$$𝒟_{}^{}{}_{22}{}^{}=\frac{y_1^2(0)}{\omega _{}^{}{}_{}{}^{2}}_0^{\mathrm{}}\zeta (\tau ^{})\zeta (\tau ^{}\tau )e^{\gamma \tau }𝑑\tau \frac{y_1(0)y_2(0)}{\omega ^{}}_0^{\mathrm{}}\zeta (\tau ^{})\zeta (\tau ^{}\tau )\tau e^{\gamma \tau }𝑑\tau $$
(52)
and the time dependence of $`y_1`$ and $`y_2`$ in the diffusion coefficients have been frozen under weak noise approximation.
Now using the linear transformation (24) in Eq.(43) we obtain in the stationary state
$$\frac{}{U}\lambda UP_s+𝒟_s\frac{^2P_s}{U^2}=0$$
(53)
where
$$U=\alpha _1y_1+y_2\mathrm{and}\lambda U=\alpha _1y_2\overline{\omega }^2y_1+\overline{\gamma }y_2$$
(54)
and
$$𝒟_s=𝒟_{}^{}{}_{22}{}^{}$$
(55)
where for simplicity it has been assumed that $`𝒟_{}^{}{}_{21}{}^{}`$ is much small compared to the Markovian contribution $`𝒟_{}^{}{}_{22}{}^{}`$.
Comparing the coefficients of $`y_1`$ and $`y_2`$ on both sides of Eq.(46) we obtain
$`\lambda \alpha _1=\overline{\omega }^2\mathrm{and}\lambda =\alpha _1+\overline{\gamma }`$
Therefore we have
$$\alpha _1=\frac{\overline{\gamma }\sqrt{\overline{\gamma }^2+4\overline{\omega }^2}}{2}\mathrm{and}\lambda =\frac{\overline{\gamma }+\sqrt{\overline{\gamma }^2+4\overline{\omega }^2}}{2}.$$
(56)
Here the negative value of $`\alpha _1`$ is taken to make $`\lambda `$ positive for a physically allowed solution of the steady state distribution (49). The solution of Eq.(45) is given by
$$P_s=N\mathrm{exp}\left(\frac{\lambda }{2𝒟_s}(\alpha _{1}^{}{}_{}{}^{2}y_1^2+2\alpha _1y_1y_2+y_2^2)\right).$$
(57)
With the help of above distribution the average quantities in tangent space can be calculated. Thus we have
$$y_1^2+y_2^2=\frac{𝒟_s}{\lambda }\left(\frac{1}{\alpha _1^2}+1\right).$$
(58)
The fluctuation dissipation relation (35) can then be obtained by combining (50) with (33) as follows ;
$$𝒟_s=\frac{\lambda }{(\frac{1}{\alpha _1^2}+1)}\mathrm{exp}\left(\underset{m}{}A_m\right).$$
(59)
$`\lambda `$ and $`\alpha _1`$ are to be calculated using (48). For these we require explicit numerical evaluation of $`\overline{\omega }^2`$ as defined in (43) and (44). The dissipative chaotic motion is governed by equations (37) and (39). We choose the following values of the parameters $`a=0.5`$, $`b=10`$, $`ϵ=10`$, $`\mathrm{\Omega }=6.07`$ and $`\gamma =0.001`$. The coupling-cum-field strength $`ϵ`$ has been varied from set to set. We choose the initial conditions $`z_1(0)=3.5`$ and $`z_2(0)=0`$ which ensures strong global chaos. Note that $`c_2`$ as expressed in (42) and in the diffusion coefficients are the integrals over the correlations of $`\zeta (t)`$ ($`\zeta (t)`$ is the fluctuating part of the second derivative of the potential $`V(z)`$ and is given by $`\zeta (t)=12az_1^2`$). To calculate the correlation function $`\zeta (t)\zeta (t\tau )`$ and the average $`\zeta (t)`$ it is necessary to determine long time series in $`\zeta (t)`$ by numerically solving the classical equation of motion (39). The next step is to carry out the averaging over the time series. For further details of the numerical procedure we refer to the earlier work . On the other hand the cumulants $`A_m(m=1,2,3,4)`$ (as defined in (34) and (35)) are calculated from Eqs.(37) and (39) directly. The method has already been outlined in Sec.(IV) and in Ref. . We then plot the theoretically calculated values of $`𝒟_s`$ from the evaluation of $`\lambda `$, $`\alpha _1`$ and the cumulants for several values of the coupling constant $`ϵ`$ (Eq.36) and compare them with the diffusion coefficients obtained from the direct numerical integration of Eqs.(39) and (37) with the appropriate transformation (22) for the corresponding values of $`ϵ`$. The result is shown in Fig. 2. It may be noted that the theoretical and numerical results are in good agreement. The validity of the fluctuation-dissipation relation as proposed in Eq.(35) is therefore reasonably satisfactory.
## VI Conclusions
The crucial question of instability of classical motion essentially rests on the linear stability matrix or Jacobian matrix associated with the equations of motion. While the linear stability analysis around the fixed points is based on the assumption of constancy of this matrix we take full account of the time dependence of the quantity in the chaotic regime by considering it to be a stochastic process, since the phase variables behave stochastically. Based on a Fokker-Planck description in the tangent space we trace the origin of chaotic diffusion and drift in the correlation of fluctuations of the linear stability matrix.
The main conclusions of this study are the following :
(i) We show that a class of dynamical stochastic parameters which attain their steady state values in the long time limit of the dynamical system may be used to characterize the dynamical steady state of the system. The first one of them which was proposed by Casartelli et. al. several years ago as a measure of the chaoticity of the system is closely related to Kolmogorov entropy.
(ii) We establish a connection between the drift and the diffusion coefficients of the Fokker-Planck equation and the dynamical stochastic parameters in the spirit of fluctuation-dissipation relation. The realization of this relation in chaotic dynamics therefore carries the message that although comprising a few degrees of freedom a chaotic system may behave as a statistical mechanical system (although in a somewhat different sense).
The theoretical relations proposed here are generic for N-degree-of-freedom chaotic Hamiltonian system with or without dissipation and have been verified by numerical analysis of a driven nonlinear dissipative system. We hope that the present approach will find useful application in searching for the related thermodynamically inspired quantities in few-degree-of-freedom systems.
###### Acknowledgements.
B. C. Bag is indebted to the Council of Scientific and Industrial Research (C.S.I.R.), Government of India, for a fellowship.
## A The derivation of the Fokker-Planck equation
We first note that the operator $`e(\tau L^0)`$ provides the solution of the equation \[Eq.(16), $`\alpha =0`$\]
$$\frac{f(X,t)}{t}=_XL^0f(X,t)$$
(A1)
$`f`$ signifies the “unperturbed” part of $`P`$ which can be found explicitly in terms of characteristic curves. The equation
$$\dot{X}=L^0(X)$$
(A2)
determines for a fixed $`t`$ a mapping from $`X(\tau =0)`$ to $`X(\tau )`$, i. e., $`XX^\tau `$ with inverse $`(X^\tau )^\tau =X`$ . The solution of (A1) is
$$f(X,t)=f(X^t,0)\left|\frac{dX^t}{dX}\right|=e\left[tF_0\right]f(X,0).$$
(A3)
$`\left|\frac{d(X^t)}{d(X)}\right|`$ being a Jacobian determinant. The effect of $`e(tL^0)`$ on $`f(X)`$ is as
$$e(tL^0)f(X,0)=f(X^t,0)\left|\frac{dX^t}{dX}\right|.$$
(A4)
This simplification in Eq.(20) yields
$`{\displaystyle \frac{P}{t}}`$ $`=`$ $`\{L^0\alpha L^1+\alpha ^2{\displaystyle _0^{\mathrm{}}}d\tau |{\displaystyle \frac{dX^\tau }{dX}}|`$ (A6)
$`L^1(X,t)_\tau L^1(𝐱^\tau ,t\tau )\left|{\displaystyle \frac{dX}{dX^\tau }}\right|\}P.`$
Now to express the Jacobian, $`X^\tau `$ and $`_\tau `$ in terms of $``$ and $`X`$ we solve Eq.(A2) for short time (this is consistent with the assumption that the fluctuations are rapid ).
We now write the solution of Eq.(A2) \[using Eqs.(4-6)\] as follows ;
$$\left(\begin{array}{c}X_1^\tau \\ \mathrm{}\\ X_N^\tau \end{array}\right)=\tau \left(\begin{array}{c}X_{N+1}\\ \mathrm{}\\ X_{2N}\end{array}\right)+\left(\begin{array}{c}X_1\\ \mathrm{}\\ X_N\end{array}\right)=\left(\begin{array}{c}\overline{G}_1(X)\\ \mathrm{}\\ \overline{G}_N(X)\end{array}\right)$$
(A7)
and
$$\left(\begin{array}{c}X_{N+1}^\tau \\ \mathrm{}\\ X_{2N}^\tau \end{array}\right)=e^{\gamma \tau }\left(\begin{array}{c}X_{N+1}\\ \mathrm{}\\ X_{2N}\end{array}\right)\tau \left(\begin{array}{c}G_{N+1}(X)\\ \mathrm{}\\ G_{2N}(X)\end{array}\right)=\left(\begin{array}{c}\overline{G}_{N+1}(X)\\ \mathrm{}\\ \overline{G}_{2N}(X)\end{array}\right)$$
(A8)
Here the terms of $`O(\tau ^2)`$ are neglected. Since the vector $`X^\tau `$ is expressible as a function of $`X`$ we write
$$X^\tau =\overline{G}(X),$$
(A9)
and the following simplification holds good;
$`L^1(X^\tau ,t\tau )._\tau `$ $`=`$ $`L^1(\overline{G}(X),t\tau )._\tau `$ (A10)
$`=`$ $`{\displaystyle \underset{k}{}}L_k^1(\overline{G}(X),t\tau ){\displaystyle \frac{}{X_k^\tau }}`$ (A11)
$`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \underset{k}{}}L_k^1(\overline{G}(X),t\tau )g_{jk}{\displaystyle \frac{}{X_j}};j,k=1\mathrm{}2N`$ (A12)
where
$$g_{jk}=\frac{X_j}{X_k^\tau }$$
(A13)
In view of Eqs.(A7) and (A8) we note:
$`\mathrm{if}j=k\mathrm{then}g_{jk}`$ $`=`$ $`1,k=1\mathrm{}N`$
$`=`$ $`e^{\gamma \tau },k=N+1\mathrm{}2N`$
$`\mathrm{if}jk\mathrm{then}g_{jk}`$ $``$ $`\tau e^{\gamma \tau }\mathrm{or}\mathrm{\hspace{0.33em}\hspace{0.33em}0}`$
Thus $`g_{jk}`$ is a function of $`\tau `$ only. Let
$$R_j=\underset{k}{}L_k^1(\overline{G}(X),t\tau )g_{jk}$$
(A14)
From Eqs.(11), (12) and (A9) we write
$$L_i^1(X^\tau ,t\tau )=L_i^1(\overline{G}(X),t\tau )=0\mathrm{for}i=1\mathrm{}N$$
(A15)
So the conditions (A14), (A15) and (A7) imply that
$`R_j(X,t\tau )`$ $`=`$ $`R_j(X_1\mathrm{}X_N,t\tau )\mathrm{for}j=1\mathrm{}N`$ (A16)
$`R_j(X,t\tau )`$ $`=`$ $`R_j(X_1\mathrm{}X_{2N},t\tau )\mathrm{for}j=N+1\mathrm{}2N`$ (A17)
We next carry out the following simplifications of $`\alpha ^2`$-term in Eq.(A6). We make use of the relation (13) to obtain
$`L^1(X,t).{\displaystyle \underset{j}{}}R_j{\displaystyle \frac{}{X_j}}P(X,t)`$ $`=`$ $`{\displaystyle \underset{i}{}}L_i^1(X,t){\displaystyle \frac{}{X_i}}{\displaystyle \underset{j}{}}R_j{\displaystyle \frac{}{X_j}}P(X,t)`$ (A18)
$`=`$ $`{\displaystyle \underset{i,j}{}}L_i^1(X,t)R_j{\displaystyle \frac{^2}{X_iX_j}}P(X,t)`$ (A20)
$`+{\displaystyle \underset{j}{}}R_j^{}{\displaystyle \frac{}{X_j}}P(X,t)`$
where
$$R_j^{}=\underset{i}{}L_i^1(X,t)\frac{}{X_i}R_j$$
(A21)
The conditions (A15) and (A17) imply that
$`R_j^{}`$ $`=`$ $`0\mathrm{for}j=1\mathrm{}N`$ (A22)
$`R_j^{}`$ $`=`$ $`R_j^{}(X_1\mathrm{}X_N,t\tau )0;\mathrm{for}j=N+1\mathrm{}2N`$ (A23)
By (A23) one has
$$R^{}.P(X,t)=.R^{}P(X,t)$$
(A24)
Making use of Eqs.(13), (A12), (A20) and (A24) in Eq.(A6) we obtain the Fokker-Planck equation (21).
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